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# Hydrostatic theory of superfluid 3He-B Associated internet page: http://boojum.hut.fi/research/theory/btex.html
## I Introduction
The superfluid phases of liquid <sup>3</sup>He show complex behavior, which still can be understood theoretically. Many phenomena have been studied in a pure form in <sup>3</sup>He, and the knowledge can then be applied to other physical systems. For example, several structures of quantized vorticity have been seen in both A and B phases of <sup>3</sup>He . The effect of impurities has recently been studied in many laboratories by aerogel immersed in liquid <sup>3</sup>He . Recent experiments on the Josephson effect show unexpected behavior . Theoretical understanding of all these phenomena requires good quantitative understanding of the basic properties of superfluid <sup>3</sup>He. This provides the motivation for the present paper.
The purpose of hydrodynamics is to determine the behavior of a fluid on length and time scales that are long compared to some microscopic lengths and times . Hydrostatics is a subfield of hydrodynamics. It is restricted to study the equilibrium properties of the fluid. In simple fluids hydrostatics reduces to statements about the pressure variation in the fluid, which either rotates uniformly or is exposed to some external field. The problem becomes more difficult, if the fluid has some broken symmetry. Particular examples of these are liquid crystals and superfluids <sup>4</sup>He and <sup>3</sup>He. In both superfluids, the equilibrium mass current belongs to the scope of the hydrostatic theory. The order parameter of <sup>3</sup>He has also other degrees of freedom. The structure of those, which is often called texture, also has to be incorporated. In <sup>3</sup>He the hydrostatic theory is limited to length scales that are large in comparison to the superfluid coherence length $`\xi _010`$nm. The theory is valid at all temperatures. The hydrostatic theory can still be applied when the motion of the quasiparticles at low temperatures becomes ballistic rather than diffusive. The hydrostatic theory can be generalized to hydrodynamic theory by using conservation laws and adding the transport coefficients.
Our purpose is to make a systematic presentation of the hydrostatic theory for the B phase of superfluid <sup>3</sup>He. Part of the reason is that presently the various results are scattered over a large literature. Previous reviews treat hydrostatics only as a side topic and cover only a small part of the subject . The existing papers are often unclear whether they treat general temperatures or are restricted to the neighborhood of the superfluid transition temperature $`T_\mathrm{c}`$. We note that the hydrostatics of the A phase is better presented in the existing literature than the B phase considered here. In addition to reviewing, we present several results that have not been published before.
We will begin with a general formulation of the hydrostatics of <sup>3</sup>He-B. Our approach is general enough to allow an external magnetic field and uniform rotation, both in the leading order. We write down an energy functional that consists of bulk terms and boundary conditions. All structures on the length scale of $`\xi _0`$, such as surface layer or vortex lines, have to be treated as boundary conditions. The theory is found to split into two pieces: one for the superfluid velocity and the other for the texture. The former is identical to the hydrostatics of superfluid <sup>4</sup>He, whereas the latter can be solved only after the superfluid velocity is determined.
The coefficients of the hydrostatic energy can either be obtained experimentally or be calculated by some more microscopic theory. The calculation is discussed in sections III and IV. The former considers the quasiclassical theory of <sup>3</sup>He . The coefficients can be calculated using the weak-coupling quasiclassical theory at arbitrary temperature. We also discuss a “trivial strong-coupling” (TSC) model, where the weak-coupling coefficients are improved by scaling the energy gap. A different approach is studied in section IV, where the hydrostatic coefficients are related to the parameters of the phenomenological Ginzburg-Landau (GL) theory at $`TT_\mathrm{c}`$.
Before any quantitative tests of the theory, we still have to determine the parameters that the quasiclassical theory needs as input. This is discussed in Section V, where we analyze several experiments. We find that certain quantities are reasonably well fitted using TSC model, but errors of 50% may occur for other quantities.
## II Hydrostatic free energy
The superfluidity in a Fermi system arises from formation of Cooper pairs. A macroscopic number of pairs occupies the same pair state in the superfluid state . The relative orbital wave function of a pair has p-wave symmetry in <sup>3</sup>He and the spin state is a triplet . The state of the pairs is thus described by an order parameter, which is a complex $`3\times 3`$ matrix $`A_{\alpha i}`$. It gives the projections of the Cooper-pair wave function on the three p-wave orbitals ($`p_x`$, $`p_y`$, and $`p_z`$, index $`i`$) and on the three spin triplet states ($`|+`$, $`\mathrm{i}|+`$, and $`|+`$, index $`\alpha `$). In unperturbed B phase the order parameter has the form
$$A_{\alpha j}=\mathrm{\Delta }e^{\mathrm{i}\varphi }R_{\alpha j}$$
(1)
with real $`\mathrm{\Delta }`$, $`\varphi `$ and $`R_{\alpha j}`$. Here the amplitude $`\mathrm{\Delta }`$ has a fixed (temperature and pressure dependent) value and $`R_{\alpha i}`$ is constrained to be a rotation matrix, i.e. $`R_{\alpha i}R_{\alpha j}=\delta _{ij}`$. (Summation over repeated indices is assumed.) The phase $`\varphi `$ and the more detailed form of the spin-orbit rotation matrix $`R_{\alpha i}`$ are not fixed on the scale of the superfluid condensation energy. These soft variables allow a dissipationless flow of both mass and spin. The order parameter can be interpreted as the wave-function of the center of mass of a pair. Using standard quantum mechanics, we can then define a mass-flow velocity
$$𝐯_\mathrm{s}=\frac{\mathrm{}}{2m}\mathbf{}\varphi ,$$
(2)
where $`m=5.009710^{27}`$ kg is the mass of a <sup>3</sup>He atom. In a similar way one can also define a spin-flow velocity
$$𝐯_{\mathrm{s},\alpha }^{\mathrm{spin}}=\frac{\mathrm{}}{4m}ϵ_{\alpha \beta \gamma }R_{\beta i}\mathbf{}R_{\gamma i}\frac{\mathrm{}}{4m}R_{\alpha i}ϵ_{ijk}R_{\beta j}\mathbf{}R_{\beta k},$$
(3)
where $`ϵ_{ijk}`$ is the maximally antisymmetric tensor. For example, if the axis of the spin-orbit rotation is constant and parallel to $`z`$, the pairs with spin states $`|`$, $`|`$, and $`|+`$ flow with velocities $`𝐯_\mathrm{s}+𝐯_{\mathrm{s},z}^{\mathrm{spin}}`$, $`𝐯_\mathrm{s}𝐯_{\mathrm{s},z}^{\mathrm{spin}}`$, and $`𝐯_\mathrm{s}`$, respectively. The three-dimensional rotation matrices are conveniently parametrized by an angle $`\theta `$ and an axis $`\widehat{𝐧}`$ of rotation as
$$R_{ij}(\widehat{𝐧},\theta )=\mathrm{cos}\theta \delta _{ij}+(1\mathrm{cos}\theta )\widehat{n}_i\widehat{n}_j\mathrm{sin}\theta ϵ_{ijk}\widehat{n}_k,$$
(4)
where $`\widehat{𝐧}\widehat{𝐧}=1`$. We note also the trivial identity $`\widehat{𝐚}\stackrel{}{R}\widehat{𝐚}=\mathrm{cos}\theta +(1\mathrm{cos}\theta )(\widehat{𝐚}\widehat{𝐧})^2`$, where $`\widehat{𝐚}`$ is an arbitrary unit vector.
The soft variables $`\varphi `$ and $`R_{\alpha i}`$ are determined by the interaction of the order parameter with various perturbations. The perturbations can be divided into external fields and boundary conditions. Experimentally, the most common field is the magnetic one $`𝐇`$. It would also be straightforward to include the electric field, but we will neglect it here because its coupling is very weak . The motion of the <sup>3</sup>He container can also be treated as an external field. In equilibrium the normal fluid component (velocity $`𝐯_\mathrm{n}`$) will follow the motion of the container, and the only allowed motions are uniform translation and rotation. The former is automatically taken into account because of Galilean invariance of the theory. The rotation will appear as a field $`\times 𝐯_\mathrm{n}`$, which equals twice the angular velocity.
It is possible to construct a hydrostatic theory for any magnitude of the external fields. Here we assume the fields are small enough that the order parameter is not distorted strongly from the bulk form (1). Strong perturbations, such as a surface or a vortex core, are treated as boundary conditions. These cause the order parameter to deviate substantially from the bulk form within a length scale of the coherence length $`\xi _010`$ nm, but on longer length scales also they act as weak perturbations.
The degrees of freedom $`\varphi `$ and $`R_{\alpha i}`$ differ crucially in the following respect. The mass flow can be written as $`𝐣_\mathrm{s}=\stackrel{}{\rho }_\mathrm{s}𝐯_\mathrm{s}`$, where $`\stackrel{}{\rho }_\mathrm{s}`$ is a phenomenological tensor. In the absence of external fields, $`\stackrel{}{\rho }_\mathrm{s}`$ must be a scalar because of the isotropy of the unperturbed $`A_{\alpha i}`$ (1). The mass current has also to be conserved: $`𝐣_\mathrm{s}=0`$. Thus we arrive at the Laplace equation $`^2\varphi =0`$. Adding the boundary conditions, this completely determines $`\varphi `$. The external fields cause only a small correction to this. Throughout the rest of this paper we neglect the small correction and start out from the assumption that the Laplace equation for $`\varphi `$ is already solved, and thus $`𝐯_\mathrm{s}`$ is known.
The problem that remains is to determine the rotation matrix $`R_{\alpha i}`$. What makes this different from $`\varphi `$ is that there exists interaction between the nuclear dipole moments of the <sup>3</sup>He atoms. It is of the form
$$F_\mathrm{D}=\lambda _\mathrm{D}d^3r(R_{ii}R_{jj}+R_{ij}R_{ji})=4\lambda _\mathrm{D}d^3r\mathrm{cos}\theta (1+2\mathrm{cos}\theta ).$$
(5)
Although this interaction is weak, it partly removes the degeneracy with respect to the rotation matrix. This means that the spin current is not conserved, but decays on a scale $`\xi _\mathrm{D}10\mu `$m. On the same scale, the rotation angle $`\theta `$ becomes fixed to $`\mathrm{arccos}(1/4)104^{}`$, which corresponds to the minimum of $`F_\mathrm{D}`$ (5), but the degeneracy with respect to the rotation axis $`\widehat{𝐧}`$ remains. Because no conservation laws exist, the rotation axis $`\widehat{𝐧}`$ is more susceptible to all kinds of perturbations than $`\varphi `$. The subject of the rest of this paper is to study the texture, i.e. $`\widehat{𝐧}(𝐫)`$ on a length scale $`\xi _\mathrm{D}`$.
We write down the free energy functional that governs the texture. The form of the energy terms is based on symmetry properties alone. The functional is valid in the limit of low fields and velocities, small gradients of the order parameter and weak coupling between the spin and orbital parts of the order parameter. The last condition is practically always satisfied because the coupling is due to the dipole-dipole interaction (5), which is small compared to the superfluid condensation energy by factor $`10^6(1T/T_c)^1`$. We neglect all constant terms, i.e. terms that do not depend on $`\widehat{𝐧}`$. The leading terms in the expansion can be written as follows
$$F_{\mathrm{DH}}=ad^3r(\widehat{𝐧}𝐇)^2$$
(6)
$$F_{\mathrm{DV}}=\lambda _{\mathrm{DV}}d^3r[\widehat{𝐧}(𝐯_\mathrm{s}𝐯_\mathrm{n})]^2$$
(7)
$$F_{\mathrm{HV}}=\lambda _{\mathrm{HV}}d^3r[𝐇\stackrel{}{R}(𝐯_\mathrm{s}𝐯_\mathrm{n})]^2$$
(8)
$$F_{\mathrm{HV1}}=\lambda _{\mathrm{HV1}}d^3r𝐇\stackrel{}{R}\mathbf{}\times 𝐯_\mathrm{n}$$
(9)
$$F_\mathrm{G}=d^3r\left[\lambda _{\mathrm{G1}}\frac{R_{\alpha i}}{r_i}\frac{R_{\alpha j}}{r_j}+\lambda _{\mathrm{G2}}\frac{R_{\alpha j}}{r_i}\frac{R_{\alpha j}}{r_i}\right].$$
(10)
The dipole-field term $`F_{\mathrm{DH}}`$ is discussed in Refs. , the gradient term $`F_\mathrm{G}`$ in Refs. , the dipole-velocity $`F_{\mathrm{DV}}`$ and field-velocity $`F_{\mathrm{HV}}`$ in Ref. , and the first-order field-velocity term $`F_{\mathrm{HV1}}`$ in Ref. . Because of Galilean invariance only the combination $`𝐯_\mathrm{s}𝐯_\mathrm{n}`$ appears in Equations (7) and (8). The superfluid velocity does not appear in the gyromagnetic term (9) because $`𝐯_\mathrm{s}`$ is curl free. Terms that are linear in $`\widehat{𝐧}𝐇`$, $`\widehat{𝐧}(𝐯_\mathrm{s}𝐯_\mathrm{n})`$ or $`𝐇\stackrel{}{R}(𝐯_\mathrm{s}𝐯_\mathrm{n})`$ are prohibited by parity and time-reversal symmetry. Equations (6)-(10) serve as definitions of the parameters $`a`$, $`\lambda _{\mathrm{DV}}`$, $`\lambda _{\mathrm{HV}}`$, $`\lambda _{\mathrm{HV1}}`$, $`\lambda _{\mathrm{G1}}`$ and $`\lambda _{\mathrm{G2}}`$, which depend on temperature and pressure. The calculation of these parameters are discussed in Sections III and IV. For some of the parameters we use names given in Ref. instead of the more systematic names introduced here; for example $`\lambda _{\mathrm{DH}}a`$ and $`\lambda _{\mathrm{SH}}d`$.
The derivation of the gradient energy (10) deserves special consideration. Originally, one starts from a general expression that is quadratic in the spin velocity (3). Making use of the properties of the rotation matrices, it is possible to simplify the energy to a form that is bilinear in the rotation matrices $`R_{\alpha i}`$. In addition to the two terms in (10), this form contains a third term of the form $`_iR_{\alpha j}_jR_{\alpha i}`$. The form (10) can then be obtained by partial integration which converts the third term into the form $`_iR_{\alpha i}_jR_{\alpha j}`$.
It should be noted that the partial integration of the gradient energy (10) produces a surface term that is similar to $`F_{\mathrm{SG}}`$ below. Thus the value of the surface coefficient $`\lambda _{\mathrm{SG}}`$ is unique only if the form of the bulk gradient energy is properly defined. Here the uniqueness of $`\lambda _{\mathrm{SG}}`$ is guaranteed by restricting the bulk gradient energy to the form (10).
The gradient term can also be expressed explicitly as a function of $`\widehat{𝐧}`$ using the representation (4) . The needed identities are given in the Appendix. Our preference is to keep the shorter form (10) because a numerical algorithm can directly be based on it.
The dipole length $`\xi _\mathrm{D}`$ is defined by $`\xi _\mathrm{D}=\sqrt{\lambda _{\mathrm{G2}}/\lambda _\mathrm{D}}`$. It is conventional to define dipole velocity $`v_\mathrm{D}`$ and dipole field $`H_\mathrm{D}`$ by writing $`\lambda _{\mathrm{HV}}=2a/(5v_\mathrm{D}^2)`$ and $`\lambda _{\mathrm{DV}}=aH_\mathrm{D}^2v_\mathrm{D}^2`$. We can also define a magnetic coherence length $`\xi _\mathrm{H}=\sqrt{65\lambda _{\mathrm{G2}}/(8aH^2)}`$, which is inversely proportional to the field. The parameters defined here are temperature dependent. Near $`T_\mathrm{c}`$ they reduce to constants that are commonly used. For example, $`\xi _\mathrm{H}R_\mathrm{c}H_\mathrm{B}/H`$ defined in Ref. .
In addition to the bulk terms (6)-(10), there are boundary terms. These energy terms originate from regions where the order parameter is strongly distorted from the form (1). We are here interested in two cases: surfaces and vortex cores. The boundary terms below are valid in the limit that the length scale of the distorted region ($`\xi _0`$) is small compared to the dipole length $`\xi _\mathrm{D}`$. In reality this is well the case. It guarantees that the rotation angle $`\theta `$ is not affected by the boundary. The form of the allowed boundary terms depends on the symmetry of the order parameter at the boundary.
We assume that the surface structure has the maximal symmetry, i.e. time-inversion symmetry, rotation symmetry around the surface normal and reflection symmetry in planes perpendicular to the surface. (We note that also less symmetric states are possible .) We also assume that the curvature of the surface is small. Such a surface gives rise to the energy terms
$$F_{\mathrm{SH}}=d_Sd^2r(𝐇\stackrel{}{R}\widehat{𝐬})^2$$
(11)
$$F_{\mathrm{SHV1}}=\lambda _{\mathrm{SHV1}}_Sd^2r𝐇\stackrel{}{R}\widehat{𝐬}\times (𝐯_\mathrm{s}𝐯_\mathrm{n})$$
(12)
$$F_{\mathrm{SG}}=\lambda _{\mathrm{SG}}_Sd^2r\widehat{s}_jR_{\alpha j}\frac{R_{\alpha i}}{r_i}$$
(13)
$$F_{\mathrm{SD}}=_Sd^2r[b_4(\widehat{𝐬}\widehat{𝐧})^4b_2(\widehat{𝐬}\widehat{𝐧})^2].$$
(14)
Here $`\widehat{𝐬}`$ is a unit vector that is perpendicular to the surface and points towards the superfluid. The surface-field term $`F_{\mathrm{SH}}`$ is discussed in Ref. , the surface-dipole term $`F_{\mathrm{SD}}`$ in Ref. , and the first-order surface-field-velocity term $`F_{\mathrm{SHV1}}`$ in Ref. . There are two contributions to the surface-gradient coefficient, $`\lambda _{\mathrm{SG}}=\lambda _{\mathrm{SG}}^\mathrm{a}+\lambda _{\mathrm{SG}}^\mathrm{b}`$. The former comes from the equilibrium spin current that flows spontaneously along any surface. In fact, the surface spin current $`J_{\alpha i}^{\mathrm{ss}}=\lambda _{\mathrm{SG}}^\mathrm{a}R_{\alpha j}ϵ_{ijk}\widehat{s}_k`$. (Note that $`J^{\mathrm{ss}}`$ contains the factor $`\mathrm{}/2`$ for each fermion and thus has the unit J/m.) The other contribution $`\lambda _{\mathrm{SG}}^\mathrm{b}`$ comes from the partial integration that depends on the chosen form of $`F_\mathrm{G}`$ . Note that there exists only one surface-gradient term (13) because $`R_{\alpha i}R_{\alpha j}`$ is antisymmetric in $`i`$ and $`j`$. For the same reason the term (13) does not depend on the normal derivative. We have constructed the definitions (6)-(14) so that all surface ($`d`$, $`\lambda _{\mathrm{SHV1}}`$, $`\lambda _{\mathrm{SG}}`$, $`b_2`$, $`b_4`$) and bulk coefficients are non-negative, at least in the Ginzburg-Landau region.
The order parameter is strongly distorted from the bulk form (1) in the cores of quantized vortex lines. Therefore the cores must be treated as boundary regions. We describe a vortex line by unit vector $`\widehat{𝐥}`$ that is parallel to the line and points in the direction of the circulation $`\times 𝐯_\mathrm{s}`$. The maximal point-symmetry operations of a vortex are a) reflection in plane perpendicular to $`\widehat{𝐥}`$, b) rotation around $`\widehat{𝐥}`$ (combined with a phase shift) and c) reflection in plane containing $`\widehat{𝐥}`$. The last one has to be combined with time inversion because otherwise the circulation would change direction. Assuming that the order parameter in the core has all these symmetries, we get the phenomenological terms
$$F_{\mathrm{LH}}=\lambda _{\mathrm{LH}}_Ld^3r(𝐇\stackrel{}{R}\widehat{𝐥})^2$$
(15)
$$F_{\mathrm{LH1}}=\lambda _{\mathrm{LH1}}_Ld^3r𝐇\stackrel{}{R}\widehat{𝐥}$$
(16)
$$F_{\mathrm{LD}}=\lambda _{\mathrm{LD}}_Ld^3r[(\widehat{𝐥}\widehat{𝐧})^2+\text{corrections}].$$
(17)
The line-field term $`F_{\mathrm{LH}}`$ is discussed in Ref. , the first-order line-field term $`F_{\mathrm{LH1}}`$ in Ref. , and line-dipole term $`F_{\mathrm{LD}}`$ in Ref. . Here $`L`$ denotes the region where vortices are present. In $`F_{\mathrm{LD}}`$ only the dominant term is written explicitly.
It is well known that the vortex cores do not have the maximal symmetry . In the A-phase-core vortex the symmetry (a) is broken. Because this can take place in two different ways, we have to assign to each vortex line a new variable $`q`$ that equals either +1 (left-handed vortex) or -1 (right handed vortex). This allows the line-gradient term
$$F_{\mathrm{LG}}=\lambda _{\mathrm{LG}}_Ld^3rq\widehat{l}_jR_{\alpha j}\frac{R_{\alpha i}}{r_i},$$
(18)
where $`\mathrm{}`$ denotes the average because $`q`$ may change from one vortex to another. Similar to the surface term $`F_{\mathrm{SG}}`$, $`F_{\mathrm{LG}}`$ arises from spontaneous spin currents. For an isolated vortex these currents form closed loops in the plane perpendicular to $`\widehat{𝐥}`$. All vortices in <sup>3</sup>He-B also have axial spin currents but they do not couple to external spin velocity in the lowest order because the net current vanishes.
The double-core vortex also allows the term (18). Additionally, the circular symmetry (b) is broken leaving only discrete symmetry in rotations by $`\pi `$. Thus an additional unit vector $`\widehat{𝐛}`$ perpendicular to $`\widehat{𝐥}`$ is needed to describe the vortex. This gives rise to line-anisotropy terms
$$F_{\mathrm{LAH}}=\lambda _{\mathrm{LAH}}_Ld^3r(𝐇\stackrel{}{R}\widehat{𝐛})^2$$
(19)
$$F_{\mathrm{LAD}}=\lambda _{\mathrm{LAD}}_Ld^3r(\widehat{𝐛}\widehat{𝐧})^2+\text{corrections}.$$
(20)
We note that there is flexibility in the definitions of the different terms. For example, the superflow around a vortex has to be counted into term $`F_{\mathrm{LH}}`$ (15) in the region where the order parameter is strongly distorted but at larger distances it also can be included as the bulk term $`F_{\mathrm{HV}}`$ (8).
## III Connection to the quasiclassical theory
The energy terms (5)-(20) contain a number of phenomenological coefficients. They should either be determined experimentally or calculated from a more microscopic theory than the hydrodynamic one. Pursuing the latter, there exists the quasiclassical theory . This theory bypasses the difficult many-body problem of strongly interacting <sup>3</sup>He atoms by concentrating in the low-energy range. It uses an expansion, where the relevant expansion parameter for the superfluid phases is the transition temperature divided by the Fermi temperature, $`T_\mathrm{c}/T_\mathrm{F}0.001`$. The lowest nontrivial order in this expansion is known as the weak-coupling theory. It effectively contains the Bardeen-Cooper-Schrieffer theory as a special case, but it also reduces to the Landau Fermi-liquid theory in the normal state. This theory is adequate for some properties of superfluid <sup>3</sup>He, especially at low pressures, but it fails, for example, to stabilize the A phase. For many purposes it is important to continue the expansion to the next order in $`T_\mathrm{c}/T_\mathrm{F}`$. We call this the strong-coupling theory. (Serene and Rainer use the name “weak-coupling plus”, but we think this is too modest since there seems to be very little hope to calculate further orders in the expansion.) We will not go into the details of the quasiclassical theory, which is extensively discussed by Serene and Rainer .
It is important to realize that the quasiclassical theory is not microscopic in the sense that it would depend only on fundamental constants. Instead, it needs several parameters as input. This is especially a problem in the strong coupling case, which needs as input the scattering amplitude of quasiparticles (in the normal state) that is not accurately known. Additionally, the needed calculations are rather complicated at general temperature. There are two practical ways to proceed. The first is to restrict to the temperature region close to $`T_\mathrm{c}`$ and use the Ginzburg-Landau theory. This approach will be described in Section IV. The second way is to work at arbitrary temperature but to use the weak-coupling approximation in the quasiclassical theory. The latter approach is discussed in this section. At the end of this section we discuss how to improve the weak-coupling results by including a trivial strong-coupling correction.
In the weak-coupling theory, the properties of the normal state are included via spin symmetric and antisymmetric Fermi-liquid parameters, $`F_l^\mathrm{s}`$ and $`F_l^\mathrm{a}`$ ($`l=0`$, 1, 2, $`\mathrm{}\mathrm{}`$). We assume that the pairing interaction is effective in the p-wave channel only. The symmetry between particle and hole types of quasiparticles is consistently assumed in the quasiclassical theory. It turns out that all results presented below depend only on five Fermi-liquid parameters: $`F_1^\mathrm{s}`$ and $`F_l^\mathrm{a}`$ with $`l=0`$, 1, 2, and 3. (Infinite number of coefficients is needed in the hydrostatics of the A phase .) In addition, the results depend on the mass density $`\rho `$ of <sup>3</sup>He liquid, on the superfluid transition temperature $`T_\mathrm{c}`$ and on the magnetic dipole-dipole interaction parameter $`g_\mathrm{D}`$.
In the Bardeen-Cooper-Schrieffer model $`g_\mathrm{D}`$ has the expression (in SI units)
$$g_\mathrm{D}=\frac{\mu _0}{40}\overline{R^2}\left(\mathrm{}\gamma N(0)\pi k_\mathrm{B}T\underset{ϵ_n=ϵ_\mathrm{c}}{\overset{ϵ_\mathrm{c}}{}}\frac{1}{\sqrt{ϵ_n^2+\mathrm{\Delta }^2}}\right)^2,$$
(21)
where $`\overline{R^2}`$ is a renormalization constant and $`ϵ_\mathrm{c}`$ a high energy cut-off. (Note that our definition of $`g_\mathrm{D}`$ is different from that in Ref. .) The Matsubara energies $`ϵ_n`$ and the weak-coupling energy gap $`\mathrm{\Delta }(T)`$ are defined in the Appendix, the gyromagnetic ratio of the <sup>3</sup>He nucleus $`\gamma =2.0410^8(\mathrm{T}\mathrm{s})^1`$, and the density of states at the Fermi energy $`2N(0)=(1+\frac{1}{3}F_1^\mathrm{s})(3m^2\rho /\pi ^4\mathrm{}^6)^{1/3}`$. It is very convenient that the dependence of $`g_\mathrm{D}`$ on temperature is so weak that we can safely ignore it. In the weak-coupling approximation the constancy of $`g_\mathrm{D}(T)`$ would be exact if the cut-off energy $`ϵ_\mathrm{c}`$ in (21) were the same in the gap equation (54). (For the standard choice $`ϵ_\mathrm{c}\mathrm{}`$ in the gap equation and $`\overline{R^2}=1`$, the relative variation of $`g_\mathrm{D}(T)`$ is less than $`10^5`$.) In the trivial strong-coupling model (see below) Eq. (21) gives the maximum variation at the melting pressure, where $`g_\mathrm{D}`$ decreases monotonically by 1.3% when $`T`$ decreases from $`T_\mathrm{c}`$ to zero (assuming $`\overline{R^2}=1`$). Because of uncertainties associated with $`ϵ_\mathrm{c}`$ and $`\overline{R^2}`$ in Eq. (21), we prefer to extract $`g_\mathrm{D}`$ from experiments, as will be discussed in section V.
For completeness, we give the results for nuclear magnetic susceptibility $`\chi `$ , superfluid density $`\rho _\mathrm{s}`$, and $`\lambda _\mathrm{D}`$ (5)
$$\chi =2\mu _0N(0)\left(\frac{\mathrm{}\gamma }{2}\right)^2\frac{\frac{2}{3}+(\frac{1}{3}+\frac{1}{5}F_2^\mathrm{a})Y}{1+F_0^\mathrm{a}(\frac{2}{3}+\frac{1}{3}Y)+\frac{1}{5}F_2^\mathrm{a}(\frac{1}{3}+(\frac{2}{3}+F_0^\mathrm{a})Y)}$$
(22)
$$\rho _\mathrm{s}=\rho \frac{1Y}{1+\frac{1}{3}F_1^\mathrm{s}Y}$$
(23)
$$\lambda _\mathrm{D}=g_\mathrm{D}\mathrm{\Delta }^2.$$
(24)
All the following coefficients can be understood as corrections to these. Here $`Y(T)=1Z_3(T)`$ is the Yoshida function, and the functions $`Z_j(T)`$ are defined in the Appendix.
The basic principle for calculating the hydrostatic parameters is explained in Section VI of Ref. . For the coefficient of the dipole-field energy $`F_{\mathrm{DH}}`$ the main part of the work, the calculation of the gap distortion, is explained in detail in Ref. . The result is
$$a=\frac{5g_\mathrm{D}}{2}\left[\frac{\frac{1}{2}\mathrm{}\gamma \mu _0(1+\frac{1}{5}F_2^\mathrm{a})}{1+F_0^\mathrm{a}(\frac{2}{3}+\frac{1}{3}Y)+\frac{1}{5}F_2^\mathrm{a}(\frac{1}{3}+(\frac{2}{3}+F_0^\mathrm{a})Y)}\right]^2\left[5\frac{3Z_5}{Z_3}\frac{3F_2^\mathrm{a}Z_3}{5(1+\frac{1}{5}F_2^\mathrm{a})}\right].$$
(25)
The coefficient of the dipole-velocity energy (7) can be calculated in a similar way and we obtain
$$\lambda _{\mathrm{DV}}=5g_\mathrm{D}\left(\frac{m^{}v_\mathrm{F}}{1+\frac{1}{3}F_1^\mathrm{s}Y}\right)^2\left(1\frac{3Z_5}{2Z_3}\right).$$
(26)
Here $`m^{}`$ is the effective mass given by $`m^{}/m=1+F_1^\mathrm{s}/3`$. The Fermi velocity $`v_\mathrm{F}`$ is related to basic parameters by $`v_\mathrm{F}=\mathrm{}(3\pi ^2\rho /m)^{1/3}/m^{}`$. As far as we know, the expressions (25) and (26) have not been published before. Rather tedious calculation is needed for the coefficient in the field-velocity energy (8). This is done in Ref. , and we quote the result
$`\lambda _{\mathrm{HV}}`$ $`=`$ $`{\displaystyle \frac{\rho }{\mathrm{\Delta }^2}}{\displaystyle \frac{m^{}/m}{(1+\frac{1}{3}F_1^\mathrm{s}Y)^2}}\left[{\displaystyle \frac{\frac{1}{2}\mathrm{}\gamma \mu _0(1+\frac{1}{5}F_2^\mathrm{a})}{1+F_0^\mathrm{a}(\frac{2}{3}+\frac{1}{3}Y)+\frac{1}{5}F_2^\mathrm{a}(\frac{1}{3}+(\frac{2}{3}+F_0^\mathrm{a})Y)}}\right]^2`$ (28)
$`\times \left[Z_3{\displaystyle \frac{9}{10}}Z_5+{\displaystyle \frac{9}{10}}{\displaystyle \frac{Z_5^2}{Z_3}}{\displaystyle \frac{3}{2}}Z_7+{\displaystyle \frac{3F_2^\mathrm{a}Z_3}{50(1+\frac{1}{5}F_2^\mathrm{a})}}(3Z_52Z_3)\right].`$
The gyromagnetic coefficient $`\lambda _{\mathrm{HV1}}=0`$ because of particle-hole symmetry . The gradient energy coefficients are calculated in Ref. , and can also be found in Appendix F of Ref. . They are
$$\lambda _{\mathrm{G2}}=\frac{\mathrm{}^2\rho }{40mm^{}}\frac{(1+\frac{1}{3}F_1^\mathrm{a})(1+\frac{1}{7}F_3^\mathrm{a})(1Y)}{1+\frac{1}{3}F_1^\mathrm{a}(\frac{2}{5}+\frac{3}{5}Y)+\frac{1}{7}F_3^\mathrm{a}(\frac{3}{5}+(\frac{2}{5}+\frac{1}{3}F_1^\mathrm{a})Y)}$$
(29)
$$\frac{\lambda _{\mathrm{G1}}}{\lambda _{\mathrm{G2}}}=2+\frac{(\frac{1}{3}F_1^\mathrm{a}\frac{1}{7}F_3^\mathrm{a})(1Y)}{(1+\frac{1}{7}F_3^\mathrm{a})(1+\frac{1}{3}F_1^\mathrm{a}Y)}$$
(30)
The structure of the order parameter near surfaces has been studied for a long time (for example in Refs. and ), but the surface terms have been evaluated only quite recently. The gyromagnetic surface term (12) vanishes identically because of particle-hole symmetry. As discussed above, the surface-gradient (13) term has two contributions: $`\lambda _{\mathrm{SG}}=\lambda _{\mathrm{SG}}^\mathrm{a}+\lambda _{\mathrm{SG}}^\mathrm{b}`$. For the part that arises from the partial integration in the derivation of (10) we find $`\lambda _{\mathrm{SG}}^\mathrm{b}=2\lambda _{\mathrm{G2}}`$. The other part $`\lambda _{\mathrm{SG}}^\mathrm{a}`$ coming from spontaneous spin currents has recently been calculated in Ref. . The same reference evaluates also the surface dipole coefficients in (14). The field term (11) has not been calculated. Until this is done, we can use an extrapolation of the Ginzburg-Landau result $`d=\frac{\mu _0}{2}(\chi _\mathrm{n}\chi )\xi _{\mathrm{GL}}d_0`$, where $`d_0=d/g_\mathrm{H}\mathrm{\Delta }^2\xi _{\mathrm{GL}}`$ is plotted by solid lines in Fig. 1. The Ginzburg-Landau coherence length $`\xi _{\mathrm{GL}}`$ can be extrapolated to general temperature by $`\xi _{\mathrm{GL}}(T)=\mathrm{}v_\mathrm{F}/\sqrt{10}\mathrm{\Delta }(T)`$. (Note that no strong-coupling correction to the weak-coupling $`\mathrm{\Delta }`$ is allowed in this equation.) $`\chi _\mathrm{n}`$ is the susceptibility in the normal state \[given by (22) with $`Y=1`$\].
For accurate calculation of the vortex terms, one needs a calculation of the order parameter in the vortex core. This has been done at general temperature by Fogelström and Kurkijärvi , but they do not give explicit values of $`\lambda _{\mathrm{LH}}`$. However, at not too high rotation velocities, the most of the contribution to $`\lambda _{\mathrm{LH}}`$ comes from outside of the vortex core, and therefore a reasonable estimate at all temperatures is
$$\lambda _{\mathrm{LH}}\frac{1}{2}\lambda _{\mathrm{HV}}|𝐯_\mathrm{s}𝐯_\mathrm{n}|^2_L\frac{\mathrm{}}{2m}\mathrm{\Omega }\lambda _{\mathrm{HV}}(\mathrm{ln}\frac{R}{r_\mathrm{L}}\frac{3}{4}),$$
(31)
where $`\mathrm{\Omega }`$ is the angular velocity of rotation, $`R=\sqrt{\mathrm{}/2m\mathrm{\Omega }}`$ the unit cell radius of a vortex, and $`r_\mathrm{L}`$ the radius of the vortex core. Because $`\lambda _{\mathrm{LH}}`$ is rather insensitive to $`r_\mathrm{L}`$, we may use $`r_\mathrm{L}\xi _{\mathrm{GL}}`$ , and use the same extrapolation of $`\xi _{\mathrm{GL}}`$ as for the surface term.
The weak coupling approximation used above is not expected to be accurate at high pressures, where strong coupling corrections are largest. As mentioned above, accurate strong-coupling calculations are very cumbersome, and introduce the scattering amplitude, which is poorly known. However, there is a simple procedure that is expected to take into account a major part of the strong coupling effects. This “trivial strong coupling correction” is to scale the energy gap $`\mathrm{\Delta }`$ by a temperature and pressure dependent factor. This factor is tabulated by Serene and Rainer as a function of the temperature and the relative jump $`\mathrm{\Delta }C/C_\mathrm{n}`$ of the specific heat at $`T_\mathrm{c}`$. The latter can be related to pressure according to the measurements by Greywall . Such scaling of $`\mathrm{\Delta }`$ affects all hydrostatic coefficients (21)-(31) directly and/or via modification of the functions $`Z_j`$ and $`Y`$.
## IV Connection to the Ginzburg-Landau theory
The hydrostatic theory was based on expansion of the free energy in small gradients and external fields. The Ginzburg-Landau (GL) theory is based on additional expansion in the amplitude of the order parameter . The expansion can be limited to a small number of terms near the superfluid transition temperature $`T_\mathrm{c}`$, where the order parameter is small. In most superconductors and in <sup>3</sup>He, the GL theory gives reliable results in the neighborhood of $`T_\mathrm{c}`$ because the fluctuation range, where it becomes invalid, consists of a negligible temperature range just at $`T_\mathrm{c}`$.
The order parameter in <sup>3</sup>He is $`3\times 3`$ matrix $`A_{\alpha j}`$. The GL theory consists of writing down the terms in the free energy that are allowed by known symmetries. The superfluid condensation energy must be invariant in separate rotations in the spin and orbital spaces. This allows the leading terms
$`F_{\mathrm{cond}}`$ $`=`$ $`{\displaystyle }d^3r[\alpha A_{\mu i}^{}A_{\mu i}+\beta _1A_{\mu i}^{}A_{\mu i}^{}A_{\nu j}A_{\nu j}+\beta _2A_{\mu i}^{}A_{\mu i}A_{\nu j}^{}A_{\nu j}`$ (33)
$`+\beta _3A_{\mu i}^{}A_{\nu i}^{}A_{\nu j}A_{\mu j}+\beta _4A_{\mu i}^{}A_{\nu i}A_{\nu j}^{}A_{\mu j}+\beta _5A_{\mu i}^{}A_{\nu i}A_{\nu j}A_{\mu j}^{}].`$
The zero of the coefficient $`\alpha =\alpha ^{}(1T/T_\mathrm{c})`$ defines the transition temperature $`T_\mathrm{c}`$. Other coefficients $`\beta _i`$ ($`i=1\mathrm{}5`$) as well as $`\alpha ^{}`$ can be taken as constants in the expansion of $`F_{\mathrm{cond}}`$ to order $`(1T/T_\mathrm{c})^2`$. In the presence of nonuniform order parameter one needs the gradient energy
$`F_\mathrm{G}`$ $`=`$ $`K{\displaystyle d^3r\left[(\gamma 2\eta )(_iA_{\mu i})^{}_jA_{\mu j}+(_iA_{\mu j})^{}_iA_{\mu j}+(2\eta 1)(_iA_{\mu j})^{}_jA_{\mu i}\right]}`$ (34)
$`=`$ $`K{\displaystyle d^3r\left[(\gamma 1)(_iA_{\mu i})^{}_jA_{\mu j}+(_iA_{\mu j})^{}_iA_{\mu j}\mathrm{i}(2\eta 1)ϵ_{kij}(\mathbf{}\times 𝐯_\mathrm{n})_kA_{\mu i}^{}A_{\mu j}\right]}`$ (35)
with the Galilean-invariant derivative $`\mathbf{}=\mathbf{}+2\mathrm{i}m𝐯_\mathrm{n}/\mathrm{}`$. In addition there are the energy caused by the magnetic field $`𝐇`$
$$F_\mathrm{H}=d^3r\left(\mathrm{i}g_{\mathrm{H1}}ϵ_{\kappa \mu \nu }H_\kappa A_{\mu i}^{}A_{\nu i}+g_\mathrm{H}H_\mu A_{\mu i}^{}A_{\nu i}H_\nu +g_\mathrm{H}^{}H^2A_{\mu i}^{}A_{\mu i}\right),$$
(36)
and the energy of the magnetic dipole-dipole interaction
$$F_\mathrm{D}=g_\mathrm{D}d^3r(A_{ii}^{}A_{jj}+A_{ij}^{}A_{ji}\frac{2}{3}A_{\mu i}^{}A_{\mu i}).$$
(37)
We neglect all terms in the free energy that are independent of $`A_{\alpha j}`$. The gradient energy (35) is parametrized using two dimensionless parameters $`\gamma `$ and $`\eta `$, which are related to parameters introduced by Serene and Rainer as $`\gamma =K_\mathrm{L}/K_\mathrm{T}`$ and $`\eta =K_\mathrm{C}/K_\mathrm{T}`$. The two different forms (35) are equivalent, as can be verified by partial integration. In contrast to the hydrostatic case (10), the surface term in the partial integration vanishes here because of the boundary condition $`\widehat{s}_iA_{\mu i}=0`$.
The parameters of the GL theory have been calculated using the weak-coupling quasiclassical theory, and the results are well known (see Refs. , for example). There are two alternatives to incorporate the strong-coupling effects. One is to determine the coefficients purely experimentally. The $`\beta _i`$’s, or at least some combinations of them, have been determined using experiments in the superfluid phases . The other alternative is to consider the GL theory as a limiting case of the strong-coupling quasiclassical theory near $`T_\mathrm{c}`$ . Here the problem of the poorly-known scattering amplitude is encountered again, but fortunately there exists model calculations for the most important coefficients. We give here a short summary of the results.
There is a small correction to $`\alpha ^{}`$ arising from finite lifetime of quasiparticles . There are several suggestions for the $`\beta _i`$’s that are based on different theoretical assumptions about the scattering amplitude and measurements in the normal state of <sup>3</sup>He . Although the strong-coupling corrections generally are small, they can be quite substantial in some combinations of $`\beta _i`$’s. For example, $`\beta _{345}\beta _3+\beta _4+\beta _5`$ may differ 50% from its weak-coupling value. The corrections to $`K`$, $`\gamma `$, $`\eta `$, $`g_\mathrm{H}`$, and $`g_\mathrm{H}^{}`$ are calculated by Serene and Rainer in Ref. . They find that $`\eta `$ is unchanged from its weak-coupling value 1, but $`\gamma `$ increases from its weak coupling value 3 to $`3.1`$ at the melting pressure. $`g_\mathrm{H}^{}`$ is found to vanish even after strong-coupling corrections, and therefore it is dropped in the following. The parameter $`g_{\mathrm{H1}}`$ vanishes in the quasiclassical theory because of particle-hole symmetry, but this term is still kept because it is important in several situations. Its value is best extracted from measurements of the splitting of the A transition in magnetic field . We have assumed that the nontrivial corrections to the dipole energy (37) are small, and therefore use the same coefficient $`g_\mathrm{D}`$ as already discussed in Sect. III.
The calculations in the GL theory are considerably simpler than in the general quasiclassical theory. Essentially all the hydrostatic parameters appearing in equations (6)-(20) have been calculated. We list below the bulk hydrostatic coefficients as functions of GL parameters.
$$a=\frac{5g_\mathrm{D}g_\mathrm{H}}{4\beta _{345}}$$
(38)
$$\lambda _{\mathrm{DV}}=\frac{5m^2g_\mathrm{D}(\gamma 1)K}{\mathrm{}^2\beta _{345}}$$
(39)
$$\lambda _{\mathrm{HV}}=\frac{2m^2g_\mathrm{H}(\gamma 1)K}{\mathrm{}^2\beta _{345}}$$
(40)
$$\lambda _{\mathrm{HV1}}=\frac{mg_{\mathrm{H1}}(\gamma 4\eta +1)K}{\mathrm{}(\beta _43\beta _1\beta _{35})}$$
(41)
$$\lambda _{\mathrm{G2}}=K\mathrm{\Delta }^2=\frac{\alpha K}{2(3\beta _{12}+\beta _{345})}$$
(42)
$$\frac{\lambda _{\mathrm{G1}}}{\lambda _{\mathrm{G2}}}=\gamma 1.$$
(43)
The first equality in (42) and (43) can be obtained trivially by substituting the B-phase order parameter (1) into the gradient energy (35). The amplitude $`\mathrm{\Delta }`$ of the order parameter is obtained by substitution into $`F_{\mathrm{cond}}`$ (33) and minimization with respect to $`\mathrm{\Delta }`$. Equations (38)-(40) can be obtained by solving the GL equations in simple cases of axially distorted B phase. For example, the coefficient (38) can be obtained by first calculating the anisotropy of the gap due to a magnetic field and then evaluating the dipole energy for this gap. The gyromagnetic coefficient $`\lambda _{\mathrm{HV1}}`$ (41) has been calculated by Mineev . Because of deviation of $`\gamma `$ from 3, it is considerably larger than anticipated in Refs. .
Accurate determination of the surface terms requires a self-consistent solution of the order parameter near a wall. In the absence of fields, the order parameter $`\stackrel{~}{A}_{\alpha i}`$, which is normalized to unit matrix in the bulk, has real components $`\stackrel{~}{A}_{xx}(x)`$ and $`\stackrel{~}{A}_{yy}(x)=\stackrel{~}{A}_{zz}(x)`$ near a surface located in the $`yz`$ plane. The surface coefficients are then obtained by integration:
$$d=g_\mathrm{H}\mathrm{\Delta }^2\xi _{\mathrm{GL}}_0^{\mathrm{}}\frac{dx}{\xi _{\mathrm{GL}}}(\stackrel{~}{A}_{yy}^2\stackrel{~}{A}_{xx}^2)$$
(44)
$$\lambda _{\mathrm{SG}}=K\mathrm{\Delta }^2_0^{\mathrm{}}𝑑x2(\gamma 1)\stackrel{~}{A}_{yy}\frac{d\stackrel{~}{A}_{xx}}{dx}$$
(45)
$$b_2=g_\mathrm{D}\mathrm{\Delta }^2\xi _{\mathrm{GL}}_0^{\mathrm{}}\frac{dx}{\xi _{\mathrm{GL}}}\frac{5}{4}(\stackrel{~}{A}_{xx}^26\stackrel{~}{A}_{xx}\stackrel{~}{A}_{yy}+5\stackrel{~}{A}_{yy}^2)$$
(46)
$$b_4=g_\mathrm{D}\mathrm{\Delta }^2\xi _{\mathrm{GL}}_0^{\mathrm{}}\frac{dx}{\xi _{\mathrm{GL}}}\frac{25}{8}(\stackrel{~}{A}_{yy}\stackrel{~}{A}_{xx})^2,$$
(47)
where $`\xi _{\mathrm{GL}}=\sqrt{K/\alpha }`$. The surface term $`F_{\mathrm{SHV1}}`$ can be found by calculating the order parameter in the presence of phase gradient: $`A_{\alpha j}(x,y)=\mathrm{\Delta }R_{\alpha i}\mathrm{exp}(\mathrm{i}ky)\stackrel{~}{A}_{ij}(x)`$. The coefficient is then given by
$$\lambda _{\mathrm{SHV1}}=\frac{2m}{\mathrm{}}g_{\mathrm{H1}}\mathrm{\Delta }^2\xi _{\mathrm{GL}}^2_0^{\mathrm{}}\frac{dx}{\xi _{\mathrm{GL}}^2}\underset{k0}{lim}\frac{1}{\mathrm{i}k}[\stackrel{~}{A}_{xj}\stackrel{~}{A}_{yj}^{}\stackrel{~}{A}_{yj}\stackrel{~}{A}_{xj}^{}].$$
(48)
The surface coefficients $`d`$, $`b_2`$, $`b_4`$ , and $`\lambda _{\mathrm{SHV1}}`$ have been estimated before using simple models for the order parameter. We calculate them here by solving the order parameter numerically using the Sauls-Serene values for the coefficients $`\beta _i`$ . For pressures below 1.2 MPa we smoothly interpolate the parameters to the weak-coupling values at zero pressure. We also assume the weak-coupling value $`\gamma =3`$. The calculations are done using boundary conditions appropriate for both specular and diffuse scattering of quasiparticles . (In the latter all the components of the order parameter have to vanish at the surface.) The integrals in Equations (44)-(48) (excluding the prefactors $`g_\mathrm{H}\mathrm{\Delta }^2\xi _{\mathrm{GL}}`$ etc.) are dimensionless, and they are plotted in Figures 1 and 2. Because $`b_2>2b_4>0`$, $`F_{\mathrm{SD}}`$ (14) is minimized by $`\widehat{𝐧}\widehat{𝐬}`$ .
All the vortex terms except $`\lambda _{\mathrm{LG}}`$ (18) have been calculated in Ref. .
All the results of this and the previous section are, of course, identical in the limit where both theories are valid: weak coupling near $`T_\mathrm{c}`$.
## V Determination of parameters
In his section we analyze a few experiments in order to deduce the values of parameters $`F_0^\mathrm{a}`$, $`F_2^\mathrm{a}`$, and $`g_\mathrm{D}`$. We apply trivial strong-coupling corrections to the gap $`\mathrm{\Delta }`$, as explained at the end of Section III. For the molar volume $`v=mN_\mathrm{A}/\rho `$ as a function of pressure we assume the fit in Ref. .
The parameter $`F_0^\mathrm{a}`$ can be obtained from measurements in the normal state: the specific heat and the nuclear susceptibility. For the specific heat we use the measurements by Greywall . The susceptibility has been measured by Ramm et al and by Hensley et al with essentially identical results. Unfortunately, it has been measured only below 2.9 MPa, and depending on the extrapolation alone, the relative error in $`1+F_0^\mathrm{a}`$ may be as large as 10% at the melting pressure. Examples are the simple fit $`F_0^\mathrm{a}=0.909+0.0055v\mathrm{cm}^3`$ and the nonmonotonic $`F_0^\mathrm{a}(p)`$ fit in Ref. .
The reduced nuclear susceptibility in the superfluid state, $`\chi (T)/\chi _\mathrm{n}`$ (22), has been measured by Corruccini and Osheroff , by Ahonen et al , and by Scholz et al . (There has been a discrepancy between the susceptibility measured by NMR and by a SQUID magnetometer , but that is probably caused by difficulties in calibration .) $`\chi (T)/\chi _\mathrm{n}`$ depends only on two parameters, $`F_0^\mathrm{a}`$ and $`F_2^\mathrm{a}`$. (We treat the trivial strong-coupling corrections as fixed, and use only the low-field limit of the Scholz data.) If we assume $`F_0^\mathrm{a}`$ given by Ramm and Hensley et al, only $`F_2^\mathrm{a}`$ remains to be fitted. We find that a nonzero pressure-independent value of $`F_2^\mathrm{a}`$ does not improve the fit essentially compared to $`F_2^\mathrm{a}0`$. Since we believe that $`F_2^\mathrm{a}`$ cannot have strong pressure dependence, the simple choice $`F_2^\mathrm{a}0`$ seems most attractive to us. Scholz finds $`F_2^\mathrm{a}1`$ with a weak-coupling fit , but this tendency is largely removed by the inclusion of trivial strong-coupling corrections. Note that the susceptibility data could be equally consistent with $`F_2^\mathrm{a}0.7`$, say, but that would imply a systematic reduction of $`F_0^\mathrm{a}`$ by -0.025 from the results by Ramm and Hensley et al. Therefore we take $`F_2^\mathrm{a}0`$ in the following. For $`F_0^\mathrm{a}`$ we use the simple fit given above because it is in better agreement with the reduced susceptibility at the melting pressure than the fit by Halperin and Varoquaux. Note that at zero pressure our choice is not far off from the relation between $`F_0^\mathrm{a}`$ and $`F_2^\mathrm{a}`$ based on the ”catastrophic relaxation” by Bunkov et al . There exists also other attempts to get $`F_2^\mathrm{a}`$ .
The dipole constant $`g_\mathrm{D}`$ has to be extracted from experiments because its value cannot be calculated accurately in the quasiclassical theory (Sect. III). The most straightforward way to get $`g_\mathrm{D}`$ is to measure the B phase longitudinal NMR frequency $`\mathrm{\Omega }_{}`$. A direct measurement of $`\mathrm{\Omega }_{}`$ has been made by Bloyet et al and Candela et al . These experiments were done at low temperatures in the collisionless regime. According to the collisionless theory in a small magnetic field
$$\mathrm{\Omega }_{}^2=\frac{45\mathrm{\Delta }^2g_\mathrm{D}}{\mathrm{}^2N(0)}\left(\frac{1}{\lambda }+\frac{2}{3}F_0^\mathrm{a}+\frac{1}{15}F_2^\mathrm{a}\right)$$
(49)
where the function $`\lambda (T)`$ is defined in the Appendix. An alternative is to extract $`\mathrm{\Omega }_{}`$ from transverse NMR frequency in surface-oriented texture. These measurements have been done by Osheroff et al. , by Ahonen, Krusius, and Paalanen (results tabulated in Ref. ), by Spencer, Alexander and Ihas , and by Hakonen et al. . Because the external magnetic field in these experiments reduces the frequency difference between normal and superfluid precession, these experiments have to be analyzed using hydrodynamic theory. There $`\mathrm{\Omega }_{}`$ is related to $`g_\mathrm{D}`$ via
$$\mathrm{\Omega }_{}^2=\frac{15\mu _0\gamma ^2\mathrm{\Delta }^2g_\mathrm{D}}{\chi }.$$
(50)
Both relations (49) and (50) imply that the temperature dependence of $`\mathrm{\Omega }_{}`$ is fully determined by the energy gap $`\mathrm{\Delta }`$ and the $`\lambda `$ function (56) or the susceptibility $`\chi `$ (22).
A third way to get $`g_\mathrm{D}`$ is so-called g shift of the transverse NMR frequency $`\omega `$ from the Larmor frequency $`\omega _0`$. According to Ref.
$$\frac{\omega \omega _0}{\omega _0}=\frac{4\mu _0a}{5\chi }.$$
(51)
The g shift was measured by Osheroff at the melting pressure and by Kycia et al below 2.17 MPa. \[We note that these measurements are done in such a high field that the expressions given in this paper are no more reliable near $`T_\mathrm{c}`$. However, Kycia et al measure the g-shift as a function of magnetization, and this plot is found field independent, both experimentally and theoretically . The only consequence is that the temperature $`T`$ in Fig. 3 is the temperature according to the weak-field susceptibility (22), which differs from the true temperature near $`T_\mathrm{c}`$.\]
We plot $`g_\mathrm{D}`$ obtained by all three methods in Fig. 3.
Let us first ignore the g-shift data (lines). It can be seen that the data for $`g_\mathrm{D}`$ at each pressure is almost independent of temperature, as required by theory. We note that in order to reach this constancy it really is necessary that the longitudinal and transverse data of $`\mathrm{\Omega }_{}`$ are analyzed with collisionless and hydrodynamic theories, respectively. Equally important is that we use trivial strong-coupling corrections. The value of $`g_\mathrm{D}`$ (but not the temperature dependence) also depends on $`F_1^\mathrm{s}`$ and $`T_\mathrm{c}`$, for which we use the measurements by Greywall .
The $`g_\mathrm{D}`$ data as a function of pressure is plotted in Fig. 4. It contains all the same data as Fig. 3 and some additional data.
It can be seen that the longitudinal and transverse measurements of $`\mathrm{\Omega }_{}`$ agree very well at intermediate pressures, but there is a difference at both high and low pressures. Theoretically the ratio of collisionless and hydrodynamic $`\mathrm{\Omega }_{}`$ (at the same pressure) depends only on the reduced temperatures $`T_1/T_\mathrm{c}`$ and $`T_2/T_\mathrm{c}`$ of the two measurements and on $`F_0^\mathrm{a}`$, $`F_2^\mathrm{a}`$, and $`\mathrm{\Delta }C/C_\mathrm{n}`$. The difference in the effective $`g_\mathrm{D}`$ obtained by the two methods seems at low pressures much larger than the expected uncertainties of the parameters, and remains unexplained. Both the positive slope of $`g_\mathrm{D}(T)`$ (Fig. 3) at high pressures and the difference between transverse and longitudinal data could be reduced by giving $`F_2^\mathrm{a}`$ a negative value ($`0.7`$ at high pressures), but only at the expense of impaired fit of the $`\chi (T)/\chi _\mathrm{n}`$.
We believe that the nontrivial strong coupling corrections are small in the expressions for $`\mathrm{\Omega }_{}`$ (49)-(50), which are based on expectation values in an unperturbed order parameter. \[There are nontrivial corrections to $`\chi `$ , but as long as (22) can reproduce (possibly with incorrect $`F_2^\mathrm{a}`$) the measured $`\chi `$, the value obtained for $`g_\mathrm{D}`$ is unaffected.\] The collisionless $`\mathrm{\Omega }_{}`$ data is fitted by solid line in Fig. 4. The values obtained for $`g_\mathrm{D}`$ depend on $`F_0^\mathrm{a}`$, $`T_\mathrm{c}`$, $`N(0)`$, $`F_2^\mathrm{a}`$, and $`\mathrm{\Delta }C/C_\mathrm{n}`$, and must be revised if more accurate values of these become available, for example, via improved measurement of the temperature .
We can also estimate $`g_\mathrm{D}`$ based on the simple model of Eq. (21). Assuming $`\overline{R^2}=1`$ and making the sum at $`T=T_\mathrm{c}`$ gives
$$g_\mathrm{D}=\frac{\mu _0}{40}\left(\mathrm{}\gamma N(0)\mathrm{ln}\frac{1.1339ϵ_c}{k_\mathrm{B}T_\mathrm{c}}\right)^2.$$
(52)
We take the cut-off energy $`ϵ_c`$ proportional to the Fermi temperature defined by $`T_\mathrm{F}=3\rho /4N(0)k_\mathrm{B}m`$. As shown by dotted line in Fig. 4, the resulting expression fits nicely the experimental data for the constant of proportionality $`ϵ_c/k_\mathrm{B}T_\mathrm{F}=0.45`$. The agreement may be accidental, however, because there is no fundamental justification for the approximations made.
Let us next study the $`g_\mathrm{D}`$ data based on the g-shift. It is also rather independent of the temperature (lines in Fig. 3). Fig. 4 shows that the $`g_\mathrm{D}`$ data (bars) at different pressures are well consistent with each other: the low pressure data extrapolates well to the melting pressure data by Osheroff. However, $`g_\mathrm{D}`$ deduced from the g shift differs essentially from the determinations based on $`\mathrm{\Omega }_{}`$, especially at high pressures. The reason for this is that the expression for $`a`$ (25) has substantial strong-coupling corrections that are not included in the scaling of the energy gap $`\mathrm{\Delta }`$. This can be seen by comparing the $`TT_\mathrm{c}`$ limit of trivial strong-coupling $`a`$ (25), denoted by $`a_{\mathrm{TSC}}`$, with the Ginzburg-Landau limit $`a_{\mathrm{GL}}`$ (38). We find
$$\frac{a_{\mathrm{GL}}}{a_{\mathrm{TSC}}}=\frac{a_{\mathrm{GL}}}{a_{\mathrm{WC}}}=\frac{g_\mathrm{H}}{g_\mathrm{H}^{\mathrm{WC}}}\frac{\beta _{345}^{\mathrm{WC}}}{\beta _{345}}.$$
(53)
where WC denotes weak-coupling. It is well known that $`\beta _{345}`$ differs substantially from its weak coupling value . This explains the difference in the apparent $`g_\mathrm{D}`$ deduced from g shift and $`\mathrm{\Omega }_{}`$. The new thing in the present analysis compared to Ref. is that the difference is not limited to the Ginzburg-Landau region, but because of the weak dependency of the apparent $`g_\mathrm{D}`$ on temperature (Fig. 3), it persists almost unchanged at all temperatures.
We conclude this section with a comparison of theoretical and experimental dipole velocity $`v_\mathrm{D}`$. Theoretically this quantity is related to coefficients $`a`$ (25) and $`\lambda _{\mathrm{HV}}`$ (28) by the relation $`v_\mathrm{D}^2=2a/(5\lambda _{\mathrm{HV}})`$. It has been measured by Nummila et al . Originally they compared their result to a theory that turned out to be in error, see discussion in Refs. . The comparison with the present theory is given in Fig. 5. We have used trivial strong-coupling theory, parameters as described above and $`g_\mathrm{D}`$ from solid line in Fig. 4.
Because both $`a`$ (38) and $`\lambda _{\mathrm{HV}}`$ (40) are proportional to $`\beta _{345}^1`$ in the Ginzburg-Landau region, the uncertainty discussed in connection with $`a`$ is expected to cancel out in $`v_\mathrm{D}`$. There also exists a direct measurement of $`\lambda _{\mathrm{HV}}`$ . It also shows deviation from the trivial strong-coupling model, but the differences are not of similar type as for $`a`$, and are presently not understood.
Above we have discussed all input parameters of the trivial-strong-coupling hydrostatic theory except $`F_1^\mathrm{a}`$ and $`F_3^\mathrm{a}`$. Out of the bulk terms only the gradient coefficients (29) and (30) depend on these. There is several independent evidence that $`F_1^\mathrm{a}1`$ at high pressures but $`F_3^\mathrm{a}`$ is not known.
## VI Conclusions
We have presented a summary of the hydrostatic theory in superfluid <sup>3</sup>He-B. Several new analytic and numerical results were included. Some experimental data was analyzed in order to extract the parameters of the theory. A particular goal was to understand how well the B phase is described by the trivial strong-coupling model. We found that some quantities ($`\chi `$, $`g_\mathrm{D}`$, $`v_\mathrm{D}`$) can successfully be calculated, but there are other quantities ($`a`$, $`\lambda _{\mathrm{HV}}`$) that may be wrong by 50% in this model. The parameter $`g_\mathrm{D}`$ has direct relevance also for the A phase, where it has been used in comparison between theory and experiment (see Ref. , for example).
The first application of the results calculated here has been the comparison of the ratio $`d/a`$ \[from equations (44) and (38)\] to experiments in Ref. . More recently, Kopu et al applied the hydrostatic theory to a rotating cylindrical container. They calculated the NMR line shape and studied the optimal conditions for observing single vortex lines. Another application is the Josephson $`\pi `$ state observed recently . This effect was found to depend crucially on the texture at the Josephson junction . The texture is also essential in several experiments of superfluid <sup>3</sup>He in aerogel. For example, the identification of the B phase was based on its texture-dependent NMR spectrum . For the present, the textural parameters in aerogel have been evaluated only in the Ginzburg-Landau region in the homogeneous scattering model . These developments demonstrate that there still are open problems in superfluid <sup>3</sup>He and in many cases a proper understanding of the hydrostatic theory is a prerequisite for solving them.
## Acknowledgments
I thank the Academy of Finland for financial support.
## Appendix
We give here some equations that complete the theory presented above. The weak-coupling energy gap $`\mathrm{\Delta }(T)`$ is determined by the equation
$$\mathrm{ln}\frac{T}{T_\mathrm{c}}+\pi T\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\left[\frac{1}{|ϵ_n|}\frac{1}{\sqrt{ϵ_n^2+\mathrm{\Delta }^2}}\right]=0,$$
(54)
where the Matsubara energies $`ϵ_n=\pi T(2n1)`$ with $`n=0,\pm 1,\mathrm{}\pm \mathrm{}`$. The $`Z_j(T)`$ functions are defined by
$$Z_j=\pi k_\mathrm{B}T\mathrm{\Delta }^{j1}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}(ϵ_n^2+\mathrm{\Delta }^2)^{j/2},$$
(55)
$`Y(T)=1Z_3(T)`$. The $`\lambda (T)`$ function , which equals to $`1f(T)`$ defined in Ref. , can be written as
$$\lambda =\pi k_\mathrm{B}T\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\frac{\mathrm{\Delta }}{\sqrt{ϵ_n^2+\mathrm{\Delta }^2}(\sqrt{ϵ_n^2+\mathrm{\Delta }^2}+\mathrm{\Delta })}.$$
(56)
The numerical calculation of the functions is discussed in Ref. .
The gradient energies (10) and (13) can be written in different forms using the identities
$`_iR_{\alpha j}_iR_{\alpha j}`$ $`=`$ $`4(1\mathrm{cos}\theta )(_i\widehat{n}_j)^2=4(1\mathrm{cos}\theta )\{(\times \widehat{𝐧})^2+(\widehat{𝐧})^2+[(\widehat{𝐧})\widehat{𝐧}\widehat{𝐧}(\widehat{𝐧})]\}`$ (57)
$`_iR_{\alpha i}_jR_{\alpha j}`$ $`=`$ $`(1\mathrm{cos}\theta )[2(\times \widehat{𝐧})^2+(1\mathrm{cos}\theta )(\widehat{𝐧})^2(1\mathrm{cos}\theta )(\widehat{𝐧}\times \widehat{𝐧})^22\mathrm{sin}\theta (\widehat{𝐧})(\widehat{𝐧}\times \widehat{𝐧})]`$ (58)
$`_iR_{\alpha j}_jR_{\alpha i}`$ $`=`$ $`(1\mathrm{cos}\theta )\{2(\times \widehat{𝐧})^2+(1\mathrm{cos}\theta )(\widehat{𝐧})^2(1\mathrm{cos}\theta )(\widehat{𝐧}\times \widehat{𝐧})^2`$ (60)
$`2\mathrm{sin}\theta (\widehat{𝐧})(\widehat{𝐧}\times \widehat{𝐧})+2[(\widehat{𝐧})\widehat{𝐧}\widehat{𝐧}(\widehat{𝐧})]\}`$
$`\widehat{s}_iR_{\alpha i}_jR_{\alpha j}`$ $`=`$ $`(1\mathrm{cos}\theta )\widehat{𝐬}[(\widehat{𝐧})\widehat{𝐧}\widehat{𝐧}(\widehat{𝐧})]\}.`$ (61)
Using these it can be seen that $`F_\mathrm{G}`$ (10) has a pure divergence term $`[(\widehat{𝐧})\widehat{𝐧}\widehat{𝐧}(\widehat{𝐧})]`$. The prefactor of this term is half of the value of that by Smith, Brinkman and Engelsberg . With present definitions the other half is transferred to the surface term (13).
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# A proper ballistic calculation of tunneling conductance for real junctions
\[
## Abstract
Employing an ab initio Screened Korringa-Kohn-Rostoker (SKKR) band structure method for a metal-vacuum-metal junction, we find that the tunnel conductance is different when it is calculated across the barrier and far from it. We attribute this difference to an artefact of the ballistic approach which overestimates the role of specular reflections, and its inability to pick up contributions from localized interface states. To reconcile the ballistic approach with experiment, we propose that the tunnel conductance should be calculated as if it is measured directly across the barrier. In this case the predicted tunneling magnetoresistance is larger.
. \]
Measurements of the resistance and magnetoresistance (MR) of magnetic planar tunnel junctions are usually made by passing a current with voltage probes far removed from the interfaces between the electrodes and insulating barrier. Therefore it makes sense to calculate these transport properties as the transmission probability from propagating eigenstates in one electrode to those of the other. To maintain a steady state current the electrodes are connected to reservoirs, and it is understood that propagating eigenstates are determined far from the electrode/barrier interfaces. While this procedure is reasonable, a consensus based on data and intuitive grounds has evolved that the current in tunnel junctions is controlled by the electronic structure at the interfaces, e.g., the local density of states (LDOS) . This conflicts with the above description which uses the electronic structure in the electrodes far from the interfaces, which is different from that at the interface. In this letter we point out that calculations made with the Landauer-Büttiker or other formalism for purely ballistic transport over the whole junction cannot be directly compared to data on real junctions for at least two reasons. The ballistic conductance overestimates the role of specular reflections from the barrier as seen by the electrodes, and, if they are present, overlooks contributions from states localized near the interfaces that are coupled to itinerant states in the electrodes by diffusive and relaxation processes in real junctions. The first causes the ballistic conductance to decrease much more rapidly than it would with diffusive electrodes; the second provides additional conduction channels. Under these conditions the electronic structure at the electrode/barrier interfaces controls the tunneling current. As we will show, to model the conductance of real planar tunnel junctions, i.e., with diffusive electrodes, it should be calculated as if it is measured directly across the barrier. In other words, to reconcile the ballistic approach with experiment one has to change the boundary conditions on the transport calculation.
We have calculated the conductance of transition-metal/vacuum tunnel junctions by using the Caroli formalism. While real junctions have insulating barriers, we have taken a vacuum as it is the simplest insulator for which we can do an ab-initio calculation. Normally, one sets the chemical potentials far from the barrier, i.e., deep in the leads. In such an approach one does not specify the chemical potentials at intermediate planes across which one chooses to evaluate the current, and one keeps the ballistic information about the propagators throughout the entire junction; not just between these planes. In our approach we set the chemical potentials at the planes we are looking at $`\mu _{\alpha (\beta )}\mu _{L(R)}`$, and so far as transport is concerned we are able to isolate the region between $`\alpha `$ and $`\beta `$ from the remainder of the junction; see Fig.1. Thus, as we see below, one is able to talk about the conductance of a part of a full junction.
In the linear response region, the tunneling conductance is
$$G=\frac{2\pi e^2}{\mathrm{}}\text{Tr}\left\{\rho ^\alpha (ϵ_F)[t^{}(ϵ_F)]^{\alpha \beta }\rho ^\beta (ϵ_F)t(ϵ_F)^{\beta \alpha }\right\},$$
(1)
where $`t^{\alpha \beta }`$ is the t-matrix between $`\alpha `$ and $`\beta `$; the trace Tr is over all (site and angular momentum, or energy level) indices; and the density of states (DOS) at the Fermi level $`\rho ^\alpha `$ and $`\rho ^\beta `$ used in this formalism are those at a surface created by cutting the junction so that the two parts thereby created are isolated from one another.
In order to obtain the t-matrix, we first determine the propagator $`G`$ from a full junction calculation and then backward derive the t-matrix from Dyson’s equation $`G=g+gVG=g+gtg`$, where $`g`$ is the propagator in the absence of the perturbation $`V`$. In our case, $`V`$ joins diffusive electrodes with a ballistically conducting system consisting of $`p`$ magnetic monolayers of a transition metal, the electrodes, on each side of a vacuum barrier of six monolayers as shown in Fig.1. We adopted an empty lattice structure for the vacuum layers which has the same properties as the metallic electrodes except that it contains no atoms. By definition $`g^{\alpha \beta }`$ is zero and, due to the fact that $`V`$ takes the form of a nearest neighbor interaction (in a principal layer description in which we use 2 atomic layers we have indeed nearest and next nearest neighbor interactions ), we find
$$G^{\alpha \beta }=g^{\alpha \alpha }t^{\alpha \beta }g^{\beta \beta }.$$
(2)
Upon inversion this yields the t-matrix, which differs for different $`\alpha `$ and $`\beta `$. By inserting this result in Eq. (1) to calculate the conductance we treat the transport between $`\alpha `$ and $`\beta `$ ballistically inasmuch as we explicitly use the propagator $`G^{\alpha \beta }`$, while the transport in the regions of the electrodes to the left of $`\alpha `$ and the right of $`\beta `$ are treated diffusively, because we do not keep track of the momentum there. Rather their effect on conduction is taken into account in a “mean field-like” manner by a self energy term in the propagator . In the limit of large $`p`$ the conductance converges to a system that is independent of $`p`$, i.e. the full ballistic junction.
The conductance for bcc(100) Fe/vacuum/Fe, and for fcc(100) Co/vacuum/Co tunnel junctions has been calculated from band structures obtained from the spin-polarized scalar-relativistic Screened Korringa-Kohn-Rostoker (SKKR) method, and the atomic sphere approximation (ASA) is used. Here we present the results for Fe; they are further coroborated by those on Co. The lattice parameter for Fe is 5.27 a.u. (atomic units). Two atomic layers are included in each screened principal layer, and the screening potential is set to $`2`$ Ry inside each atomic cell. The Gunnarsson-Lundqvist exchange-correlation potential is used, and energy integration is performed by means of Gaussian quadrature with 16 points on a semi-circle in the upper half complex energy plane. For self-consistent calculations of the bulk metal, the free metal surface and the metal-vacuum-metal interface potentials, 45 $`𝐤_{}`$ points are used in the irreducible wedge of two dimensional Brillouin Zone (2DBZ), which enables the Fermi level to be converged up to $`10^7`$ Ry. For more details on this method see Ref. . We used a small imaginary part of the energy of $`\epsilon =0.5mRyd`$ in the propagators in order to converge our results in a reasonable time.
In Fig.2 we show the ballistic conductance for Fe(p)/vac(6)/Fe(p) tunnel junctions calculated either across the barrier, i.e. $`p=0`$, or up to $`p=7`$ Fe layers from the barrier. The conductance in all channels, majority, minority and antiparallel, decrease as we are increasing the number of monolayers of the electrodes we include in the ballistic calculation; this is particularly pronounced in the minority channel. Most of the decay in the ballistic conductance is between $`p=0`$ and $`1`$, because this is the region where the electronic structure is changing most. It stands to reason that the usual approach, in which the entire junction is treated ballistically, produces a lower conductance than the one we would calculate just across the barrier. The full ballistic approach includes the additional specular reflections near the barrier that comes from the bending of the band bottoms which represents the charge that has leaked out of the metallic electrodes into the vacuum barrier. For example, we find 0.7 of a total of 3 minority electrons leak out of the surface layer of the Fe electrode; for the majority band only 0.1 of the 5 electrons leak out. In Fig.2 one sees that the effect of the leakage is far more pronounced in the minority band than the majority, and this explains why the conductance drops much more in the minority band than the majority when we take into account the reflections due to the bending of the band bottoms.
In real junctions transport is diffusive, particularly in the electrodes, so that the total resistance is the sum of resistances of each part; in this case one should limit a ballistic calculation to just the barrier. As the resistance of the barrier in tunnel junctions is about $`10^5`$ times greater than that of the electrodes, even for junctions with resistances as low as $`50\mathrm{\Omega }\mu m^2`$, it is reasonable to say that the barrier determines the resistance ($`p=0)`$ and the electrodes give negligible contributions, i.e., in real planar junctions the conductance does not vary much as one goes away from the barrier. The decay of the ballistic conductance for increasing $`p`$ is thus an artefact of having considered the transport in the electrodes ballistically.
We note that the tunneling magnetoresistance MR is larger when it is measured across the barrier. While we obtain MR ratios $`(G_PG_{AP})/G_P`$ of 82% and 38% for the Fe and Co junctions when we use the conductances calculated at $`p\mathrm{}`$, these ratios increase to 86% and 65% when measured directly across the barrier, i.e., at the interfaces $`p=0`$.
In addition, as we will show now for the case of the Fe/vacuum junction, diffuse scattering and relaxation processes may couple localized states at the interfaces to propagating states in the electrodes. As in Fe these localized states are dominant in the minority channel close to the Fermi energy, they give an additional rise to the MR. The complete eigenvalue spectrum of a semi-infinite solid contains the continuum of bulk states as well as additional states localized at the surface . In ferromagnetic metals such as Fe(100) the existence of localized states at the surface has been known for some time ; more recently surface states were also observed in Gd(0001) . If the energy of the surface states lie in the gaps between the bands of the bulk states they form true localized states; otherwise they are resonant states - admixtures of itinerant and localized states. Tunneling is certainly affected by these resonant states at the interfaces of the junction; see for example Fig.2b in Ref. where these resonances appear along the $`\mathrm{\Gamma }M`$ direction away from the zone center. The true localized states, on the contrary, are orthogonal to resonant and itinerant states; therefore ballistic transport will be unaffected by them. However, it has been shown that the conductance through the localized states at Fe/vacuum and Gd/vacuum interfaces can be substantial because these states have orbits that point out from the surface into the barrier. Therefore it is the diffusive nature of the transport and the ambient relaxation that allows the localized states to contribute to conduction.
At $`T=0K`$ electrons can scatter elastically from the itinerant states only to localized states at the Fermi energy and vice versa; this could come from impurities and roughness of the interface. The scattering produced by impurities, e.g., in Fe, was calculated by Mertig to be about $`1\mu \mathrm{\Omega }cm/atomic\%`$; therefore, for $`1atomic\%`$ impurities, we can roughly estimate the elastic scattering produced at the interfaces is equivalent to the scattering rate of the order of $`1/\tau _{imp}=10^{14}\mathrm{sec}^1`$. For the tunnel junctions studied to date with resistances in the range of $`10^310^8\mathrm{\Omega }\mu m^2`$, as well as for the Fe/vac/Fe junction we will discuss, the tunneling rate is in the range $`10^610^{10}\mathrm{sec}^1`$, so that the diffuse scattering in sufficient to have localized states contribute to participate in conduction. At finite temperatures relaxation processes, such electron-electron, electron-phonon, and electron-magnon interactions, can couple the localized to itinerant states; our rough estimates tell us that the relaxation is faster than the tunneling rate of $`10^6\mathrm{sec}^1`$when $`T>0.7K`$, while one needs $`T20K`$ for the localized states to participate in the conduction when the tunneling rate is $`10^{10}\mathrm{sec}^1`$. We conclude that while tunneling electrons create holes in the localized interface states, they recombine almost instantaneously due to scattering by phonons, magnons, other electrons, or interfacial disorder so that for the calculation of the current one can always assume a Fermi distribution even for the localized states. On the contrary for currents perpendicular to the plane of the layers (CPP) in metallic multilayered structures the rate at which electrons traverse a layer is determined by the Fermi velocity, which is the of the order of $`10^{16}\mathrm{sec}^1`$. As the relaxation mechanisms are much slower, it is reasonable to calculate the conductance in metallic multilayers by neglecting relaxation to localized states, even though they appear at interfaces in much the same way as surface states .
At the Fe(100)/vac interface surface states exist in the minority channel at the Fermi level for $`k_{}0`$; for $`k_{}0`$ localized states exist above $`ϵ_F`$ ; therefore those with $`k_{}0`$ contribute to tunneling if one applies a bias. In our ASA calculations there are only surface resonant states at $`ϵ_F`$, and we have found localized states in the minority channel about $`k_{}0`$ just below the Fermi level at $`ϵϵ_F0.05eV`$. We also calculated the surface density of states at $`0.05eV`$ below the Fermi level, as well as the conductance $`0.05eV`$ below the Fermi level both at the interface and in the bulk, i.e., $`p=0`$ and $`4`$ in terms of the in-plane momentum $`k_{}`$. The large DOS at the surface about $`k_{}=0`$ at $`0.05eV`$ below the Fermi level, which is absent at $`ϵ_F`$ and in the bulk $`p=4`$ indicates the presence of the localized surface state. On comparing the conductances in, only the conductance for the barrier $`p=0`$ and at $`ϵ_F0.05eV`$ has a strong contribution from the localized states about $`k_{}=0`$; all the other conductances have “holes” about $`k_{}=0`$. One notes that the average of the conductance for $`p=0`$ at $`ϵ_F0.05eV`$ is four times larger than at $`ϵ_F`$. The conclusion that can be drawn is that if localized states exist about $`k_{}0`$ for $`Eϵ_F`$ they would contribute to the conductance measured across the barrier, but do not contribute to the ballistic conductance away from the interfaces. In real junctions where localized states are mixed with resonant and itinerant at the surface they contribute to the conduction as measured across electrodes far from the barrier; therefore it is only the ballistic conductance calculated for the barrier itself, $`p=0`$, that captures the contribution from localized states if they exist. In general the contribution of localized states to tunneling will depend on their coupling to the states in the barrier, i.e. their chemical bonding.
In conclusion, when one compares the conductance of real magnetic tunnel junctions with diffusive electrodes to a calculation of conductance where it is assumed that transport is ballistic throughout, e.g., for large $`p`$, we find it overestimates the role of specular reflections as seen by the electrodes, and, if they are present, overlooks contributions from states localized near the interfaces that are coupled to itinerant states in the electrodes by diffusive and relaxation processes in real planar junctions. Both mechanisms present in all planar tunnel junctions conspire to maintain the conductance relatively constant as we go away from the barrier, i.e., the decrease in the ballistic conductance does not apply to realistic planar junctions. For these reasons the only conductance one obtains in a ballistic calculation that can be compared to real junctions is that across the barrier $`p=0`$. This should not suggest that ballistic calculations far away from the barrier are meaningless; one can certainly think of systems such as two iron whiskers separated by a vacuum or MgO barrier where conductance in the electrodes is purely ballistic at very low temperatures.
We have used vacuum whereas the barriers in the tunnel junctions studied to date have been insulators; while this changes the conductance one calculates, it does not alter the conclusion we arrive at, i.e., if one does a ballistic calculation it should be that of only the barrier. For finite bias one probes a larger region about the Fermi level so that localized surface states away from $`ϵ_F`$ contribute to conduction; their contribution to the conductance will be included if one does the calculation directly across the barrier rather than in the electrodes.
We would like to acknowledge and thank William Butler for sharing with us his unpublished results on the tunneling conductance of Fe/vac/Fe junctions, Matthias Bode for bringing to our attention his spin polarized tunneling results, and Phivos Mavropolous and Nickos Papanikolaou for helpful discussions. This work was supported by the Defense Advanced Research Projects Agency and Office of Naval Research (Grant No. N00014-96-1-1207 and Contract Nos. MDA972-96-C-0014, and MDA972-99-C-0009 ), the National Science Foundation (Grant No. INT-9602192), NATO (Grant No. CRG 960340), the Hungarian National Science Foundation (OTKA T030240), and the TMR Network ERBFMXCT-960089.
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# Abstract
## Abstract
We study the diffusion of Brownian particles on the surface of a sphere and compute the distribution of solid angles enclosed by the diffusing particles. This function describes the distribution of geometric phases in two state quantum systems (or polarised light) undergoing random evolution. Our results are also relevant to recent experiments which observe the Brownian motion of molecules on curved surfaces like micelles and biological membranes. Our theoretical analysis agrees well with the results of computer experiments.
Let a diffusing particle start from the north pole of a sphere at time $`\tau =0`$. We join the final position of the particle at time $`\beta `$ to its initial one (the north pole) by the shorter geodesic. This rule is well defined, unless the final position is exactly at the south pole, a zero probability event. The path followed by the diffusing particle (closed by the geodesic rule) encloses a solid angle $`\mathrm{\Omega }`$. The question we address is: At time $`\beta `$, what is the distribution $`P^\beta (\mathrm{\Omega })`$ of solid angles?
An experimental motivation for this question comes from recent time-resolved fluorescence studies on the Brownian motion of rod-like molecules on curved surfaces such as micelles and lipid vesicles. The experimentally measured fluorescence anisotropy is affected by the curvature of the surface due to geometric phase effects. The problem of diffusion on curved surfaces reappears in Nuclear Magnetic Resonance (NMR) and Electron Spin Resonance (ESR) studies on micelles and lipid vesicles. The curvature of the surface affects the spin relaxation times. In biological systems, the translational diffusion of solutes bound to various curved surfaces influences metabolic rates and transmission rates of chemical signals. Recently, fluorescence anisotropy decay of molecules on curved surfaces has been studied theoretically and using Monte Carlo simulations. This study shows that the geometric phase crucially affects the anisotropy decay if the molecules are tilted relative to the surface normal. We therefore need to theoretically evaluate the distribution of geometric phases or equivalently, the distribution of solid angles in order to understand diffusion processes on curved surfaces.
In this paper we present an answer to the question posed above and compare our theoretical results with Monte Carlo simulations. The theoretical analysis presented here is based on earlier work . In the distribution of solid angles for closed Brownian paths on the sphere is computed. However, closed Brownian paths are a set of measure zero among Brownian paths and a comparison of the results of with real or computer experiments is hampered by poor statistics. We adapt the methods of to allow for open Brownian paths and compute the distribution of solid angles for such paths, closing the path by the geodesic rule in order to make the solid angle a well-defined quantity. The qualitative idea behind and the present paper is to use a magnetic field as a “counter”, to measure the solid angle enclosed in a Brownian Motion.
We first illustrate our method by solving a similar problem on the plane. Let a diffusing particle start from the origin $`\stackrel{}{r}=0`$ of the plane at time $`\tau =0`$. Let us suppose that the particle arrives at $`\stackrel{}{r_f}`$ at time $`\beta `$. Join $`\stackrel{}{r_f}`$ to the origin by a straight line. The open Brownian path $`\stackrel{}{r}(\tau )`$ closed by a straight line encloses an area $`A`$. What is the probability distribution $`P(A)`$ of areas? By “area” we mean the algebraic area, including sign. Area enclosed to the left of the diffusing particle counts as positive and area to the right as negative. This problem has been posed and solved earlier. We illustrate our method by solving this problem before moving to the problem of main interest in this paper.
Let $`\{\stackrel{}{r}(\tau )=\{x(\tau ),y(\tau )\},0\tau \beta ,\stackrel{}{r}(0)=\stackrel{}{0}\}`$ be a realization of a Brownian path on the plane which starts at the origin. As is well known, Brownian paths are distributed according to the Wiener Measure : if $`[\stackrel{}{r}(\tau )]`$ is any functional on paths, the expectation value of $``$ is given by
$$<[\stackrel{}{r}(\tau )]>_𝒲:=\frac{d\stackrel{}{r}_f𝒟[\stackrel{}{r}(\tau )][\stackrel{}{r}(\tau )]exp[_0^\beta \{\frac{1}{2}\frac{d\stackrel{}{r}}{d\tau }.\frac{d\stackrel{}{r}}{d\tau }d\tau \}]}{d\stackrel{}{r}_f𝒟[\stackrel{}{r}(\tau )]exp[_0^\beta \{\frac{1}{2}\frac{d\stackrel{}{r}}{d\tau }.\frac{d\stackrel{}{r}}{d\tau }d\tau \}]}$$
(1)
In Eq. (1) the symbol $`𝒟[\stackrel{}{r}(\tau )]`$ denotes a functional integral over all paths which start at the origin and end at $`\stackrel{}{r}_f`$. Finally, the endpoint $`\stackrel{}{r}_f`$ is also integrated over. (We set the diffusion constant equal to half.) $`\beta `$ is the time for which the diffusion has occured . Let $`𝒜[\stackrel{}{r}(\tau )]`$ be the algebraic area enclosed by the path $`\stackrel{}{r}(\tau )`$. Clearly, the normalized probability distribution of areas $`P(A)`$ is given by
$$P(A):=<\delta (𝒜[\stackrel{}{r}(\tau )]A)>_𝒲.$$
(2)
The expectation value $`\overline{\varphi }`$ of any function $`\varphi (A)`$ of the area is given by $`P(A)\varphi (A)𝑑A`$. As is usual in probability theory, we focus on the generating function $`\stackrel{~}{P}(B)`$ of the distribution $`P(A)`$ :
$$\stackrel{~}{P}(B):=\overline{e^{iBA}}=P(A)e^{iBA}𝑑A,$$
(3)
which is simply the Fourier transform of $`P(A)`$. For reasons that will soon be clear, we use the symbol $`B`$ for the Fourier transform variable. The distribution $`P(A)`$ can be recovered from its generating function by an inverse Fourier transform. From Eqs.(2) and (3) above we find
$$\stackrel{~}{P}(B)=<e^{iB𝒜}>_𝒲$$
(4)
Let us introduce a fictitious vector potential $`\stackrel{}{A}.d\stackrel{}{r}=(B/2)(xdyydx)`$. Clearly $`\stackrel{}{A}.d\stackrel{}{r}=B𝒜`$, where the integral is over the closed circuit consisting of the Brownian path closed by the geodesic. $`\stackrel{}{A}`$ has been chosen so that its radial component $`\stackrel{}{A}.\stackrel{}{r}`$ vanishes and consequently the contribution of $`\stackrel{}{A}.d\stackrel{}{r}`$ along the geodesic segment from $`\stackrel{}{r}_f`$ to $`\stackrel{}{0}`$ vanishes. It follows that $`B𝒜`$ can be expressed as
$$B𝒜[\stackrel{}{r}(\tau )]=_0^\beta \stackrel{}{A}(\stackrel{}{r}(\tau )).\frac{d\stackrel{}{r}(\tau )}{d\tau }d\tau ,$$
(5)
Eqs.(1), (4) and (5) yield
$$\stackrel{~}{P}(B)=\frac{d\stackrel{}{r}_f𝒟[\stackrel{}{r}(\tau )]exp[_0^\beta \{\frac{1}{2}\frac{d\stackrel{}{r}}{d\tau }.\frac{d\stackrel{}{r}}{d\tau }d\tau \}+i_0^\beta \{\stackrel{}{A}.\frac{d\stackrel{}{r}}{d\tau }d\tau \}]}{d\stackrel{}{r}_f\{𝒟[\stackrel{}{r}(\tau )]exp[_0^\beta \{\frac{1}{2}\frac{d\stackrel{}{r}}{d\tau }.\frac{d\stackrel{}{r}}{d\tau }d\tau \}]}$$
(6)
By inspection of Eq.(6) we arrive at:
$$\stackrel{~}{P}(B)=\frac{Y(B)}{Y(0)},$$
(7)
where $`Y(B)`$ is given by $`Y(B)=𝑑\stackrel{}{r}_fK^B(\stackrel{}{0},\stackrel{}{r}_f)`$, where $`K^B(\stackrel{}{0},\stackrel{}{r}_f)`$ is the quantum amplitude for a particle of unit charge and mass to go from the origin $`\stackrel{}{0}`$ to $`\stackrel{}{r}_f`$ in imaginary time $`\beta `$ in the presence of a homogenous magnetic field. This amplitude can also be expressed as
$$K^B(\stackrel{}{0},\stackrel{}{r}_f)=\underset{\{n\}}{}\mathrm{exp}[\beta E_n]u_n^{}(\stackrel{}{0})u_n(\stackrel{}{r}_f),$$
(8)
where $`u_n(\stackrel{}{r})`$ are a complete set of normalised eigenstates of the Hamiltonian $`\widehat{H}=(1/2)(i\stackrel{}{}\stackrel{}{A})^2`$ and $`E_n`$ are the corresponding eigenvalues. (Throughout this paper we set $`\mathrm{}=c=1`$.) We thus arrive at the expression: $`Y(B)=𝑑\stackrel{}{r}_f_{\{n\}}exp[\beta E_n]u_n^{}(\stackrel{}{0})u_n(\stackrel{}{r}_f)`$.
Now we demonstrate the utility of Eq.(7) by computing the distribution of areas for diffusion on a plane. The function $`Y(B)`$ for a particle of unit mass and charge in a constant magnetic field, is easily computed from the energies $`E_n=(n+\frac{1}{2})B`$ and the eigenfunctions $`u_n(\stackrel{}{r}_f)=\sqrt{B}\mathrm{exp}[(B/4)r_f{}_{}{}^{2}]L_n[(B/2)r_f{}_{}{}^{2}]`$ with $`L_n`$, the $`n^{th}`$ Laguerre polynomial. $`r_f`$ is the modulus of the vector $`\stackrel{}{r}_f`$.
$$Y(B)=\frac{2\pi }{\mathrm{cosh}(\frac{\beta B}{2})}.$$
From (7) we find $`\stackrel{~}{P}(B)=[\mathrm{cosh}(\beta B/2)]^1.`$ Taking the Fourier transform of $`\stackrel{~}{P}(B)`$ by contour integration we get the result $`P(A)=[\beta \mathrm{cosh}(\pi A/\beta )]^1`$ as derived in . This provides a check on Eq.(7) and illustrates its use.
Let us now address the problem posed at the beginning of this paper: What is the distribution $`P(\mathrm{\Omega })`$ of solid angles ($`\mathrm{\Omega }`$) enclosed by a Brownian particle starting at time $`\tau =0`$ at the north pole $`\widehat{r}_n`$ of an unit sphere and ending at any other point $`\widehat{r}_f`$ on the sphere, the final point being joined to the initial one by a geodesic. (As stated earlier this rule breaks down only if the final point is the south pole, which is a zero probability event). Unlike the planar case, $`P(\mathrm{\Omega })`$ is a periodic function with period $`4\pi `$. The generating function $`\stackrel{~}{P}_g`$ of the distribution of solid angles is given by
$$\stackrel{~}{P}_g=_0^{4\pi }𝑑\mathrm{\Omega }P(\mathrm{\Omega })e^{ig\mathrm{\Omega }}$$
(9)
with $`g`$ a half integer. $`P(\mathrm{\Omega })`$ is expressed in terms of $`\stackrel{~}{P}_g`$ by a Fourier series (rather than an integral) with $`g`$ ranging from $`\mathrm{}`$ to $`\mathrm{}`$ in half integer steps:
$$P(\mathrm{\Omega })=\frac{1}{4\pi }\underset{g=\mathrm{}}{\overset{\mathrm{}}{}}e^{ig\mathrm{\Omega }}\stackrel{~}{P}_g.$$
(10)
Relation (7) now takes the form
$$\stackrel{~}{P}_g=\frac{Y_g}{Y_0},$$
(11)
where $`Y_g`$ is given by $`Y_g=𝑑\widehat{r}_f_{\{j,m\}}exp[\beta E_{j,m}^g]u_{j,m}^g(\widehat{r}_n)u_{j,m}^g(\widehat{r}_f)`$ with $`u_{j,m}^g(\widehat{r}_f)`$ the normalised eigenfunctions and $`E_{j,m}^g`$ the energy eigenvalues for a quantum particle of unit charge on a sphere subject to a magnetic field created by a monopole of quantized strength $`g`$ at the center of the sphere. The quantum numbers of the eigenstates are $`j,m`$, where $`j`$ is the total angular momentum quantum number and $`m`$ is its $`z`$ component. We choose the vector potential in the form $`A_\varphi =g(1\mathrm{cos}\theta ),A_\theta =0`$, which is non singular everywhere on the sphere except the south pole. Since $`A_\theta `$ vanishes, (as $`A_r`$ did in the planar case), $`A`$ along the open Brownian path does measure the solid angle as defined by the geodesic rule.
The Hamiltonian for this problem is $`\widehat{H}^g=\frac{1}{2}\frac{1}{\mathrm{sin}\theta }\frac{}{\theta }\mathrm{sin}\theta \frac{}{\theta }+(i\frac{}{\varphi }A^g{}_{\varphi }{}^{})^2`$ in standard spherical co-ordinates. The wave equation can be separated by writing $`\psi (\theta ,\varphi )=exp(ip_\varphi \varphi )R(\theta )`$. From continuity of $`\psi `$, it follows that $`R(0)=0`$ if $`p_\varphi 0`$. Therefore, only those eigenstates for which $`p_\varphi =0`$, contribute to $`Y^g`$. These normalised eigenstates are labelled by $`j`$, which ranges from $`|g|`$ to $`\mathrm{}`$ in integer steps.
$$R^g{}_{j}{}^{}(z)=\sqrt{\frac{2j+1}{2}}\frac{1}{2^g}(1+z)^gP^{0,2g}{}_{jg}{}^{}(z),$$
(12)
where $`z=\mathrm{cos}\theta `$ and $`P^{a,b}_n`$ are the Jacobi Polynomials . The energy levels of this system are : $`E^g{}_{j}{}^{}=\frac{j(j+1)g^2}{2}`$.Using the integer $`n=j|g|`$ to label the states (in place of $`j`$) we find
$$\stackrel{~}{P}^g=\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\mathrm{exp}([n(n+1)+g(2n+1)]\beta /2)\frac{(2n+2g+1)g}{(g+n)(g+n+1)},$$
(13)
for $`g>0`$ and $`\stackrel{~}{P}^0=1`$. From (10) we arrive at
$$P(\mathrm{\Omega })=(1/4\pi )(1+2\underset{g=1/2}{\overset{\mathrm{}}{}}\mathrm{cos}(\mathrm{\Omega }g)\stackrel{~}{P}_g),$$
(14)
where $`g`$ increases in half integer steps. The function (14) is plotted numerically for various values of $`\beta `$ in Fig. $`1`$. The qualitative nature of these plots is easily understood. For small values of $`\beta `$ the particle tends to make small excursions and its path encloses solid angles close to $`0`$ or $`4\pi `$ and consequently the plots are peaked around these two values. As the available time $`\beta `$ increases, other values of $`\mathrm{\Omega }`$ are also probable and the peaks tend to spread and the curves to flatten out. Finally in the limit of $`\beta \mathrm{}`$ the particle has enough time to enclose all possible solid angles with equal probability. These plots give the answer to the question that was raised in the beginning of the paper.
The analytical solution given above was checked against the results of computer experiments. The diffusion process was simulated on a Silicon Graphics workstation by Monte Carlo methods. The simulation methodology was adapted from (where more details are given) and is briefly as follows. We start the simulation with $`10^5`$ molecules located at the north pole of a sphere with all molecular dipoles aligned parallel to the spherical surface (perpendicular to the position vectors). Let $`\widehat{V}_0`$ represent the dipoles at zero time, oriented along the laboratory x-axis which is perpendicular to the position vector $`\widehat{R}_0`$ along the laboratory z-axis. These molecules are then subjected to random “kicks” and allowed to diffuse over the sphere independent of each other. The diffusion of the dipoles on the spherical surface is effected by rotation of the position vector $`\widehat{R}`$ and the dipole vector $`\widehat{V}`$ by the same three dimensional rotation matrix about a randomly chosen vector $`\widehat{n}`$ normal to the radial position vector $`\widehat{R}`$. In the simulation, the diffusion constant $`D`$ is chosen to be $`10^3`$ per time step and the angle of rotation per time step (one iteration) is obtained from the corresponding probability equation for planar diffusion, which is valid for very small diffusion lengths. The orientation of the molecule $`\widehat{V}`$ during diffusion is used as an index of the solid angle swept out by the Brownian path. We numerically evaluate this for each molecule applying the geodesic rule. Let $`\widehat{V}_\beta `$ and $`\widehat{R}_\beta `$ represent the dipole and the radial vectors at time $`\beta `$. To determine the orientation of $`\widehat{V}_\beta `$, $`\widehat{V}_\beta `$ is transformed to $`\widehat{V^{}}_\beta `$ using a three dimensional rotation matrix with an angle of rotation (less than $`\pi `$) determined by $`cos^1(\widehat{R}_\beta .\widehat{R}_0)`$. The angle between $`\widehat{V^{}}_\beta `$ and $`\widehat{V}_0`$ is the solid angle enclosed by the Brownian path of the diffusing molecule at time $`\beta `$. Similarly, the solid angles are calculated for all the $`10^5`$ dipoles at different times. These are sorted into 360 different bins, each corresponding to one degree increment between $`0`$ and $`2\pi `$ and are plotted in Figure 2.
Since the orientation of the molecule is only determined modulo $`2\pi `$ in the simulation, $`\mathrm{\Omega }`$ is only measured modulo $`2\pi `$ and not modulo $`4\pi `$. So in effect the function computed by the simulation is $`Q(\mathrm{\Omega })=P(\mathrm{\Omega })+P(\mathrm{\Omega }+2\pi )`$. From (14) it is clear that the half integer values of $`g`$ cancel out and the integer values are doubled. Our theoretical expression for $`Q(\mathrm{\Omega })`$ is therefore
$$Q(\mathrm{\Omega })=(1/2\pi )(1+2\underset{g=1}{\overset{\mathrm{}}{}}\mathrm{cos}(\mathrm{\Omega }g)\stackrel{~}{P}_g),$$
(15)
where $`g`$ now increases in integer steps. The results of the numerical simulation and the function $`Q(\mathrm{\Omega })`$ (with suitable cutoffs on $`g`$ and $`n`$) given above are plotted in Figure 2 for different values of $`\beta `$. The agreement between the simulations and the theory is excellent and the fractional deviation is of order $`1/\sqrt{N}`$ where $`N`$ is the number of molecules in a bin.
In computing the distribution of solid angles, we have chosen all the molecules to be initially at the north pole ($`\widehat{r}_i=\widehat{r}_n`$) and integrated over the final point $`\widehat{r}_f`$ with a uniform weight, to match with the procedure followed in our simulation. In an experimental situation, where the diffusing particles are excited and observed with external probes, it may be necessary to consider a weighted average over the initial and final positions of the diffusing particles . Our analysis is easily adapted to take this into account. One just writes
$$Y^g=𝑑\widehat{r}_fw_f(\widehat{r}_f)𝑑\widehat{r}_iw_i(\widehat{r}_i)K^g(\widehat{r}_i,\widehat{r}_f).$$
(16)
This leads to $`Y^g=_nw^g{}_{}{}^{}{}_{i}{}^{}w_{}^{g}{}_{f}{}^{}\mathrm{exp}[\beta E_n^g/2]`$, where $`w^g{}_{n}{}^{}=d\widehat{r}w(\widehat{r})u_n^g(\widehat{r})`$ and the sum is over all eigenstates. In this paper we have set $`w_i(\widehat{r}_i)=\delta ^2(\widehat{r}_i\widehat{r}_n)`$ and $`w_f(\widehat{r}_f)=1`$.
In the problem of diffusion of fluorescent molecules, the tangential component of the orientation of the molecule determines the solid angle enclosed by the diffusing particle only modulo $`2\pi `$. This is because, in parallel transport of a vector, the vector rotates by an angle equal to the solid angle enclosed by the curve. As Pancharatnam showed , a quantum two state system transported on the Poincaré sphere picks up a geometric phase equal to half the enclosed solid angle. This measures the solid angle modulo $`4\pi `$. The distribution $`P(\mathrm{\Omega })`$ computed in the body of the text describes the distribution of geometric phases. However, since our simulations were adapted from , the distribution measured in the simulation is $`Q(\mathrm{\Omega })`$ and not $`P(\mathrm{\Omega })`$.
In , the generating function for the distribution of solid angles for closed Brownian paths involved only the energy eigenvalues of the Hamiltonian. In contrast, in the present case of open Brownian paths (closed by the geodesic rule), one needs to use both the eigenvalues as well as the eigenfunctions of the Hamiltonian. This is the main difference between and the present theoretical analysis.
We have computed the distribution of solid angles for Brownian motion on the sphere. This calculation has also been checked against computer experiments. In simple geometries like the surface of a sphere, an analytical treatment is possible. In more complicated situations, one is forced to rely on computer simulations. Monte Carlo simulations can be used in two distinct roles. In the present problem, where an analytical treatment is possible, simulations can be viewed as a computer experiment against which the theory can be tested. Such experiments are much easier to control than real experiments. In cases where an analytical treatment is not possible, the simulations serve as a substitute for the theory, to be checked against real experiments. The agreement between theory and computer experiments for Brownian motion on the sphere also serves as a check on the computer simulations so that these simulations can be confidently used in geometries which cannot be treated analytically.
Acknowledgements: It is a pleasure to thank R. Nityananda for discussions on this subject and S. Ramasubramanian for drawing our attention to Ref. and Y. Hatwalne for a critical reading of the manuscript. JS and SS thank N. Kumar for raising the question of the distribution of solid angles on a sphere. MMGK thanks N. Periasamy and S.W. Englander for their encouragement. We thank Alain Comtet for drawing our attention to some relevant references.
## Figure Captions
Fig. 1: The solid angle distributions $`P(\mathrm{\Omega })`$ computed from Eq. 14 at five different values of $`\beta `$: 0.5, 1, 2, 5 and 10. The flatter curves correspond to higher $`\beta `$.
Fig. 2: Comparision between solid angle distributions $`Q(\mathrm{\Omega })`$ from Monte Carlo simulations (red line) and the theory (Eq. 15) (blue line). The plot shows the distributions at five different values of $`\beta `$: 0.5, 1, 2, 5 and 10.
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# Finite Temperature Collapse of a Bose Gas with Attractive Interactions
## I introduction
The trapped Bose gas with attractive interactions is a novel physical system. At high densities such clouds are mechanically unstable; however at low densities they can be stabilized by quantum mechanical and entropic effects. Stability and collapse has been observed in clouds of degenerate <sup>7</sup>Li and <sup>85</sup>Rb . The collapse of the Bose gas shares many features of the gravitational collapse of cold interstellar hydrogen, described by the Jeans instability ; related instabilities occur in supercooled vapors. Theoretical studies of the attractive Bose gas, typically numerical, have been limited to zero or very low temperature . Here we give a simple analytic description of the region of stability and threshold for collapse valid from zero temperature to well above the Bose condensation transition, and thus provide a consistent global picture of the instability. As we see, at finite temperature, the phase diagram of the Bose gas includes regions where the non-condensed particles play a significant role in the collapse .
Our results are summarized in the phase diagram in Fig. 1 which shows three regions: normal, Bose condensed, and collapsed. This third region is not readily accessible experimentally; the system becomes unstable at the boundary (the solid line in the figure). In this figure the condensation and collapse lines actually meet at a finite angle. The gross features of the phase diagram can be understood qualitatively with dimensional arguments. At low temperatures the only stabilizing force is the zero-point motion of the atoms. This “quantum pressure” has a characteristic energy per particle of $`E_Q\mathrm{}^2V^{2/3}/m`$, where $`\mathrm{}`$ is Planck’s constant, $`V`$ is the volume in which the cloud is confined, and $`m`$ is the mass of an atom. The attractive interactions which drive the collapse are associated with an energy $`U=\mathrm{}^2a_sn/m`$, where $`n`$ is the density, and $`a_s`$ ($`=1.45`$nm for <sup>7</sup>Li and $`20`$nm for <sup>85</sup>Rb) is the s-wave scattering length. Comparison of $`U`$ with $`E_Q`$ suggests that for low temperatures the density $`n`$ at collapse should be $`(V^{2/3}a_s)^1`$, independent of temperature. At high temperatures, the stabilizing force is thermal pressure, $`P=TS/V`$, which is characterized by the thermal energy $`E_Tk_BT`$. Here comparison of $`U`$ with $`E_T`$ suggests that at high temperatures the density of collapse should be $`mk_BT/\mathrm{}^2a_s`$, linear in temperature. The crossover between the quantum and classical behaviors ($`nT^0`$ and $`nT`$) occurs near $`T_c`$. The collapse, as described here, is a phenomena in which the cloud as a whole participates, not just the condensate.
Our basic approach is to identify the collapse with an instability in the lowest energy mode of the system (the breathing mode in a spherically symmetric cloud). As the system goes from stable to unstable, the frequency of this mode goes from real to complex, passing through zero when the instability sets in. The density response function, $`\chi (k,\omega )`$, which measures the response of the cloud to a probe at wave vector $`k`$ and frequency $`\omega `$, diverges at the resonant frequencies of the cloud. Therefore, by virtue of the vanishing frequency of the lowest energy mode, the collapse is characterized by
$$\chi (k=2\pi /L,\omega =0)^1=0.$$
(1)
Here $`k=2\pi /L`$ is the wave-vector of the unstable mode, whose wavelength $`L`$ should be of order the size of the system. Equation (1) implicitly determines the line of collapse in the temperature-density plane.
To evaluate the response function analytically, we use a local-density approximation, replacing the response function $`\chi (k,\omega )`$ of the inhomogeneous cloud by that of a gas with uniform density $`n`$. The response of the uniform gas is evaluated at the same frequency and wave-vector as for the inhomogeneous system, and $`n`$ is given by the central density of the atomic cloud. The local-density approximation should be valid for temperatures large enough that the thermal wavelength, $`\mathrm{\Lambda }_T=(2\pi \mathrm{}^2/mk_BT)^{1/2}`$, is much smaller than the size of the trap. In all experiments to date, this condition is satisfied, and we treat $`k\mathrm{\Lambda }_T`$ as a small parameter in our calculation. With this approximation, we calculate the line of collapse using the well developed theory of $`\chi (k,\omega )`$ of a uniform gas (reviewed in ).
In the experiments on <sup>7</sup>Li, the atoms are held in a magnetic trap with a harmonic confining potential $`V(r)=\frac{1}{2}m\omega ^2r^2`$, with $`\omega 2\pi \times 145s^1`$. This potential gives the cloud a roughly Gaussian density profile (see Sec. IV). The temperature $`T`$ is typically 50 times the trap energy $`\mathrm{}\omega 7`$nK, so the parameter $`k\mathrm{\Lambda }_T`$, $`\sqrt{\mathrm{}\omega /k_BT}0.14`$, is small. The experiments on <sup>85</sup>Rb use softer traps, $`\mathrm{}\omega 0.6`$nK, and colder temperatures $`T15`$nK, so that $`k\mathrm{\Lambda }_T0.2`$.
The goal of our approach is to provide a framework for investigating the interplay of condensation and collapse. Although our results are not as accurate as can be obtained numerically (comparing with previous numerical work , we find that are results are always well within a factor of two of those calculated using more sophisticated models), the conceptual and computational advantages of working with a uniform geometry far outway any loss in accuracy. Due to their simplicity, the arguments used here provide an essential tool to choosing which parameter ranges to investigate in future experiments and computations.
The tools introduced to discuss the collapse of an attractive gas can also be used to describe other instabilities in trapped Bose gases. In Section V and VI, we apply these methods to the problem of BCS type pairing between bosons, and towards domain formation in spinor condensates.
## II Simple Limits
### A Zero Temperature
To illustrate our approach we first consider the stability of a zero temperature Bose condensate. The excitation spectrum of a uniform gas is :
$$\omega ^2=\left(\frac{k^2}{2m}\right)^2+gn\frac{k^2}{m}E_k^2,$$
(2)
where $`g=4\pi \mathrm{}^2a_s/m`$. In the attractive case, $`g<0`$, all long wavelength modes with $`k^2/2m<2|g|n`$ have imaginary frequencies and are unstable. A system of finite size $`L`$ only has modes with $`k>2\pi /L`$, and for larger $`k`$ has an excitation spectrum similar to Eq. (2). If $`|g|n<\mathrm{}^2\pi ^2/mL^2`$, the unstable modes are inaccessible and the attractive Bose gas is stable.
This information is included in the density response function $`\chi (k,\omega )`$ of the dilute zero-temperature gas ,
$$\chi (k,\omega )=\frac{2n\epsilon _k}{\omega ^2E_k^2},$$
(3)
where $`\epsilon _k=k^2/2m.`$ The poles of $`\chi `$ are at the excitation energies, $`\pm E_k`$. In particular, $`\chi (k=2\pi /L,\omega =0)`$ diverges when $`|g|n=\mathrm{}^2\pi ^2/mL^2`$.
### B High Temperature ($`TT_c`$)
We further illustrate our procedure by calculating the stability of an attractive Bose gas at temperature much larger than $`T_c`$, where thermal pressure is the predominant stabilizing force. Quantum effects are negligible in this limit, and the line of collapse is simply the spinodal line of the classical liquid-gas phase transition , as characterized by Mermin . We neglect finite size effects, and look for an instability in the uniform gas at zero wavevector, $`k=0`$, corresponding to finding where $`\chi (k=0,\omega =0)`$ diverges. The susceptibility $`\chi (0,0)=n/\mu `$ (where $`\mu `$ is the chemical potential), is proportional to the compressability of the system, which diverges when the gas becomes unstable.
At high temperature we work in the Hartree-Fock approximation, where the density is given by the self-consistent solution of
$$n=\frac{d^3p}{(2\pi \mathrm{})^3}\frac{1}{e^{\beta (ϵ_p\mu )}1},$$
(4)
with Hartree-Fock quasiparticle energies $`ϵ_p=\epsilon _p+2gn`$; here $`\beta =1/k_BT`$. In the classical limit ($`\beta \mu 1`$), $`n=e^{\beta (\mu 2gn)}/\mathrm{\Lambda }_T^3`$. The response $`\chi (0,0)`$ has the structure of the random phase approximation (RPA),
$$\chi (k,\omega )=\frac{\chi _0(k,\omega )}{12g\chi _0(k,\omega )},$$
(5)
where $`\chi _0(0,0)=\left(n/\mu \right)_ϵ`$ (where the $`ϵ`$ are held fixed) is the “bare” response. In the classical limit $`\chi _0(0,0)=\beta n`$. Since $`\chi _0(0,0)`$ is negative, the repulsive system ($`g>0`$) is stable. However, for attractive interactions $`g<0`$), the denominator of Eq. (5) vanishes when $`2g\chi _0=1`$, which in the classical limit occurs when $`2|g|n=k_BT`$.
The above calculation is only valid well above $`T_c`$. When $`|\mu |\mathrm{}^2/mL^2`$, finite size effects start to become important, and a more sophisticated approach is needed. If one blindly used the above result near $`T_c`$ one would erroneously find that the instability towards collapse prevents Bose condensation from occurring. This difficulty can be avoided by working with the finite wave-vector response $`\chi (k=2\pi /L,\omega =0)`$, to which we now turn.
## III Density Response Function
The approximation we shall use for $`\chi `$ is to treat the response of the gas in the RPA, with the simplifying assumption that the bare response of the condensed and non-condensed particles is taken to be the response of a non-interacting system. This approach, employed by Szépfalusy and Kondor in studying the critical behavior of collective modes of a Bose gas, and later modified by Minguzzi and Tosi to include exchange, is simple to evaluate analytically, and is valid both above and below $`T_c`$. It generates an excitation spectrum which is conserving and gapless . At zero temperature it yields the Bogoliubov spectrum, Eq. (2), and above $`T_c`$ it becomes the standard RPA with exchange.
The susceptibility in this approximation has the form,
$$\chi (k,\omega )=\frac{\chi _0^c+\chi _0^n+g\chi _0^c\chi _0^n}{(1g\chi _0^c)(12g\chi _0^n)4g^2\chi _0^c\chi _0^n},$$
(6)
where $`\chi _0^c`$ and $`\chi _0^n`$ are the condensate and non-condensed particle contributions to the response of the non-interacting cloud,
$`\chi _0^n(k,\omega )`$ $`=`$ $`{\displaystyle \frac{d^3q}{(2\pi )^3}\frac{f(qk/2)f(q+k/2)}{\omega (\epsilon _{q+k/2}\epsilon _{qk/2})}},`$ (8)
$`\chi _0^c(k,\omega )`$ $`=`$ $`{\displaystyle \frac{n_0}{\omega \epsilon _k}}{\displaystyle \frac{n_0}{\omega +\epsilon _k}}.`$ (9)
Here $`n_0`$ is the condensate density, the $`\epsilon _k=k^2/2m`$ are the free particle kinetic energies as before, and the Bose factors $`f(k)`$ are given by $`(e^{\beta (\epsilon _k\mu )}1)^1`$. In Appendix A we briefly review the derivation of this response function. At zero temperature $`\chi _0^n=0`$ and the susceptibility reduces Eq. (3), while above $`T_c`$, $`\chi _0^c=0`$, and $`\chi `$ reduces to Eq. (5). Figure 2 shows the class of diagrams summed in this approximation.
Expanding $`\chi _0^n`$ in the small parameter $`k\mathrm{\Lambda }_T`$ (see details in Appendix B), we derive for $`T>T_c`$,
$$g\chi _0^n(k,0)=2\frac{a_s}{\mathrm{\Lambda }_T}\left[\frac{4\pi }{k\mathrm{\Lambda }_T}\mathrm{arctan}\left|ϵ_k/4\mu \right|^{1/2}+g_{1/2}(e^{\beta \mu })\left|\pi /\beta \mu \right|^{1/2}+𝒪(k\mathrm{\Lambda }_T)\right],$$
(10)
where $`g_\nu (z)_jz^j/j^\nu `$ is the polylogarithm function. For chemical potential $`\mu `$ much larger in magnitude than $`k_BT`$, the system is classical, and Eq. (10) reduces to $`g\chi _0^n=\beta gn`$, as in the Hartree-Fock approach, Sec. II B. Below $`T_c`$ the chemical potential of the non-interacting system vanishes and the response functions are:
$`g\chi _0^n(k,0)`$ $`=`$ $`{\displaystyle \frac{4\pi ^2a_s}{k\mathrm{\Lambda }_T^2}}+𝒪((k\mathrm{\Lambda }_T)^0),`$ (12)
$`g\chi _0^c(k,0)`$ $`=`$ $`16\pi {\displaystyle \frac{a_sn_0}{k^2}}.`$ (13)
Using these expressions we calculate the spinodal line separating the stable and unstable regions of Fig. 1 by setting $`k=2\pi /L`$ and solving the equation
$$\chi ^11g(\chi _0^c+2\chi _0^n)2g^2\chi _0^c\chi _0^n=0,$$
(14)
which gives the line of collapse as a function of $`\mu `$ and $`T`$ (for $`T>T_c`$) or as a function of $`n_0`$ and $`T`$ (for $`T<T_c`$). We use the following relations to plot the instability on the $`nT`$ phase diagram,
$$n=\{\begin{array}{cc}\mathrm{\Lambda }_T^3g_{3/2}(e^{\beta \mu }),\hfill & T>T_c\hfill \\ n_0+\mathrm{\Lambda }_T^3\zeta (3/2),\hfill & T<T_c.\hfill \end{array}$$
(15)
Equations (12) and (13) indicate that below $`T_c`$ the noncondensate response $`\chi _0^n`$ scales as $`k^1L`$, while the condensate response $`\chi _0^c`$ scales as $`k^2L^2`$. For realistic parameters, $`L`$ is the largest length in the problem, so that the condensate dominates the instability except when $`n_0`$ is much smaller than $`n`$. Since the condensate is very localized, even a few particles in the lowest mode make $`n_0`$ locally much greater than the density of noncondensed particles. In Fig. 3 we show how the line of instability depends on the size of the system, $`L`$.
From the above equations we calculate the maximum stable value of the condensate density $`n_0`$. The line of collapse crosses the line of condensation at a temperature $`T^{}=\mathrm{}^2/2mk_BL|a|`$. Above this temperature no condensation can occur. For $`T<T^{}`$ the collapse limits the density of condensed particles to
$$(n_0)_{\mathrm{max}}=\frac{\pi }{4L^3}\left(\frac{L}{|a|}\right)\left(\frac{T^{}T}{T^{}+T}\right);$$
(16)
$`(n_0)_{\mathrm{max}}`$ decreases monotonically with temperature, from the value $`(n_0)_{\mathrm{max}}^{\mathrm{T}=0}=\pi /4L^2|a|`$, eventually vanishing at $`T=T^{}`$. Using parameters from experiments , we find for the Rice <sup>7</sup>Li trap, $`(N_0)_{\mathrm{max}}^{\mathrm{T}=0}=L^3(n_0)_{\mathrm{max}}^{\mathrm{T}=0}=1700`$, and $`T^{}`$ = 7.5 $`\mu `$K, while for the JILA <sup>85</sup>Rb trap, $`(N_0)_{\mathrm{max}}^{\mathrm{T}=0}=120`$, and $`T^{}`$ = 46 nK, The maximum number of particles vs. temperature for the two experiments are plotted in Fig. 4; these results are consistent with the experiments, and agree quite well with numerical mean-field calculations . In particular, our curve $`(N_0)_{\mathrm{max}}(T)`$ for lithium has a slope of $`1/2.2`$ nK at $`T=0`$, which lies between the calculated slopes of Davis et al. and Houbiers et al. . Although $`(n_0)_{\mathrm{max}}`$ decreases with temperature, the non-condensed density $`n^{}`$ increases ($`n^{}=\mathrm{\Lambda }^3\zeta (3/2)`$). Thus the total density at collapse need not be monotonic with temperature (cf. the low temperature region of Fig. 1b).
Future condensate experiments at higher temperatures and densities should be able to study the structure in Eq. (16), and map out the phase diagram in Fig. 1. The rubidium experiments are performed at temperatures near $`T^{}`$, where the spinodal line intersects $`T_c`$, and in principle should be able to explore the crossover between the quantum mechanical and classical behavior of the instability. The lithium experiments are much further away from exploring this regime, and in the current geometry, inelastic processes make such an investigation impractical . Since $`T^{}`$ is proportional to $`1/L`$, a softer trap could be used to bring this crossover down to lower temperatures where these difficulties are less severe (see Fig. 3). More precise numerical studies at higher temperatures are needed to guide these experiments.
## IV Modeling the harmonic trap
Most experimental and theoretical results are reported in terms of numbers of particles instead of density. By appropriately modeling the density distribution of a harmonically trapped gas, we can present our conclusions in such a form. Once the interactions are strong enough to modify the density distribution significantly, the system undergoes collapse; thus we can take the density distribution to be that of non-interacting particles. For $`k_BT\mathrm{}\omega `$, the density profile is well-approximated by
$$n(r)=\frac{d^3p}{(2\pi )^3}\frac{1}{e^{\beta (\epsilon _p+V(r)\mu )}1}+n_0e^{r^2/d^2},$$
(17)
where $`V(r)=m\omega ^2r^2/2`$ is the confining potential, with characteristic length $`d=(\mathrm{}/m\omega )^{1/2}`$. The density of condensed particles at the center of the trap is $`n_0`$. Above $`T_c`$, $`n_0=0`$, and below $`T_c`$, $`\mu =0`$. Integrating over space, we have
$$N=\{\begin{array}{cc}\left(k_BT/\mathrm{}\omega \right)^3g_3(e^{\beta \mu }),\hfill & T>T_c\hfill \\ \left(k_BT/\mathrm{}\omega \right)^3\zeta (3)+\left(\pi \mathrm{}/m\omega \right)^{3/2}n_0,\hfill & T<T_c.\hfill \end{array}$$
(18)
The instability occurs in the lowest energy mode of the system, the breathing mode, whose wave-vector is proportional to $`1/d`$. In a zero temperature non-interacting gas the breathing mode has a density profile $`\delta \rho (2r^2/d^23)\mathrm{exp}(r^2/d^2)`$, where $`r`$ is the radial coordinate. In momentum space this distribution is peaked at wave-vector $`k=2/d`$. At finite temperature thermal pressure increases the radius of the cloud and the wave vector of the breathing mode becomes smaller. Since the response of the non-condensed cloud is relatively insensitive to the wave-vector, we look for an instability at $`k=2/d`$.
The resulting phase diagram, Fig. 5, is similar to that in Fig. 1. The most significant difference is that the line of collapse follows the condensation line (on the scale of the figure they appear to coincide over a significant temperature range). This behavior can be understood by noting that for trapped particles, condensation results in a huge increase in the central density of the cloud (a standard diagnostic of BEC).
## V Pairing
With minor changes the formalism presented here can be used to investigate the instability towards forming loosely bound dimers, or “pairs,” the Evans-Rashid transition. Such an instability occurs in an electron gas at the superconducting BCS transition , and has been predicted by Houbiers and Stoof to occur in the trapped alkalis. The pairing is signalled by an instability in the T-matrix of the normal phase , which plays the role that the density response function plays in the collapse. Again, we simulate the finite size of the cloud by looking for an instability at $`k=2\pi /L`$, as opposed to $`k=0`$ in a bulk sample. In analogy to Eq. (5), the T-matrix can be written as a ladder sum,
$$𝒯(k,\omega )=\frac{g}{1g\mathrm{\Xi }(k,\omega )}.$$
(19)
In this equation, $`k`$ is the relative momentum of the pair. The instability towards pairing is signalled by $`𝒯\mathrm{}`$, when $`g\mathrm{\Xi }=1`$. To the same level of approximation as Eq. (8), the medium-dependent part of the “pair bubble” $`\mathrm{\Xi }`$ is
$$\mathrm{\Xi }(k,\omega )=\frac{d^3q}{(2\pi )^3}\frac{f(qk/2)+f(q+k/2)}{\omega (\epsilon _{q+k/2}+\epsilon _{qk/2})}.$$
(20)
Setting $`\omega =0`$, and expanding in small $`k\lambda _T`$, we find
$$g\mathrm{\Xi }(k,\omega =0)=4\frac{a_s}{\mathrm{\Lambda }_T}\left[\frac{4\pi }{k\mathrm{\Lambda }_T}\mathrm{arctan}\left(\frac{\left|ϵ_k/4\mu \right|^{1/2}}{1+\left|ϵ_k/4\mu \right|^{1/2}}\right)+g_{1/2}(e^{\beta \mu })\left|\pi /\beta \mu \right|^{1/2}+𝒪(k\mathrm{\Lambda }_T)\right].$$
(21)
Except for the argument of the arctangent, this expression is identical to twice $`g\chi _0^n`$ as given in Eq. (10). Since arctangent is a monotonic function, and its argument here is smaller than in Eq. (10), we see that $`g\mathrm{\Xi }<2g\chi _0^n`$, which implies that the instability towards collapse occurs at a lower density than the pairing instability. Thus we conclude that the pairing transition does not occur in an attractive Bose gas. Interestingly, in the classical limit, the instabilities towards pairing and collapse coincide.
## VI Domain Formation in Spinor Condensates
The approach used here to discuss the collapse of a gas with attractive interactions also describes domain formation in spinor condensates, and gives a qualitative understanding of experiments at MIT in which optically trapped <sup>23</sup>Na is placed in a superposition of two spin states. Although all interactions in this system are repulsive, the two different spin states repel each other more strongly than they repel themselves, resulting in an effective attractive interaction. The collapse discussed earlier becomes, in this case, an instability towards phase separation and domain formation. The equilibrium domain structure is described in . Here we focus on the formation of metastable domains.
The ground state of sodium has hyperfine spin $`F=1`$. In the experiments the system is prepared so that only the states $`|1=|F=1,m_F=1`$ and $`|0=|F=1,m_F=0`$ enter the dynamics. The effective Hamiltonian is then
$$H=d^3r\frac{\psi _i^{}\psi _i}{2m}+V(r)\psi _i^{}\psi _i+\frac{g_{ij}}{2}\psi _i^{}\psi _j^{}\psi _j\psi _i.$$
(22)
where $`\psi _i`$ ($`i=0,1`$) is the particle destruction operator for the state $`|i`$; summation over repeated indices is assumed. The effective interactions, $`g_{ij}=4\pi \mathrm{}^2a_{ij}/m`$, are related to the scattering amplitudes $`a_{F=0}`$ and $`a_{F=2}`$, corresponding to scattering in the singlet ($`F_1+F_2=0`$) and quintuplet ($`F_1+F_2=2`$) channel, by :
$`\stackrel{~}{a}a_{11}=a_{01}=a_{10}=a_{F=2},`$ (24)
$`\delta aa_{11}a_{00}=(a_{F=2}a_{F=0})/3.`$ (25)
Numerically, $`\stackrel{~}{a}=2.75`$nm and $`\delta a=0.19`$nm. We introduce $`\stackrel{~}{g}=4\pi \mathrm{}^2\stackrel{~}{a}/m`$ and $`\delta g=4\pi \mathrm{}^2\delta a/m`$. In the mean field approximation, the interaction in Eq. (22) becomes a function of $`n_{m=0}`$ and n, the density of particles in the state $`|0`$ and the total density, respectively:
$$H_{int}=d^3r\left(\frac{\stackrel{~}{g}}{2}n^2\frac{\delta g}{2}n_{m=0}^2\right),$$
(26)
which shows explicitly the effective attractive interaction. Initially the condensate is static with density $`n=10^{14}\mathrm{cm}^3`$, and all particles in state $`|1`$. A radio-frequency pulse places half the atoms in the $`|0`$ state without changing the density profile. Subsequently the two states phase separate and form domains from 10 to 50 $`\mu `$m thick. The trap plays no role here, so we can neglect $`V(r)`$ in Eq. (22) and consider a uniform cloud.
Linearizing the equations of motion with an equal density of particles in each state, we find two branches of excitations corresponding to density and spin waves ,
$`\omega _{ph}^2`$ $`=`$ $`\left({\displaystyle \frac{k^2}{2m}}\right)^2+\stackrel{~}{g}n{\displaystyle \frac{k^2}{m}}+𝒪(\delta g),`$ (28)
$`\omega _{sp}^2`$ $`=`$ $`\left({\displaystyle \frac{k^2}{2m}}\right)^2\delta gn{\displaystyle \frac{k^2}{4m}}+𝒪((\delta g)^2).`$ (29)
Since $`\delta g>0`$ spin excitations with imaginary frequencies appear. The mode with the largest imaginary frequency grows most rapidly, and the width of the domains formed should be comparable to the wavelength $`\lambda `$ of this mode. By minimizing Eq. (29) we find $`\lambda =\sqrt{2\pi /n\delta a}=10\mu `$m, in rough agreement with the observed domain size.
## VII Acknowledgements
The authors are grateful to the Ecole Normale Supérieure in Paris, and the Aspen Center for Physics, where this work was carried out. We owe special thanks to Eugene Zaremba and Dan Sheehy for critical comments. We are particularly indebted to Henk Stoof for his insightful recommendations and stimulating discussions, including raising the question of the relation between the instabilities towards pairing and collapse. This research was supported in part by the Canadian Natural Sciences and Engineering Research Council, and National Science Foundation Grant No. PHY98-00978, and facilitated by the Cooperative Agreement between the University of Illinois at Urbana-Champaign and the Centre National de la Recherche Scientifique.
## A Review of the RPA
Here we give a brief derivation of the response function $`\chi (k,\omega )`$, Eq. (6), used in this paper. Generically, the response of a gas, $`\chi =\delta n/\delta U`$, is the direct response to the perturbation $`\chi _0`$ plus the response to the mean field generated by the disturbed atoms. For a normal gas in the Hartree approximation, $`\delta n=\chi _0\delta U+\chi _0g\delta n`$, while including exchange gives $`\delta n=\chi _0\delta U+\chi _02g\delta n`$, as in Eq. (5). The RPA amounts to making a simple particular approximation to the polarization part $`\chi _0`$. For our purposes it suffices to take $`\chi _0=\chi _0^n`$, the response of an ideal gas, Eq. (8).
Generalizing the Hartree approximation to the condensed gas simply requires replacing $`\chi _0`$ with the response of the non-condensed particles $`\chi _0^n`$ plus the response of the condensate $`\chi _0^c`$, Eq. (9).
Including exchange in the degenerate gas requires some work, since exchange only occurs in interactions involving non-condensed atoms. A simple technique, demonstrated by Minguzzi and Tosi is to look separately at the change in the density of condensed and non-condensed atoms, $`\delta n_0`$ and $`\delta \stackrel{~}{n}`$. Within Hartree-Fock these changes are related to the applied perturbation $`\delta U`$ by
$$\left(\begin{array}{c}\delta n_0\\ \delta \stackrel{~}{n}\end{array}\right)=\left(\begin{array}{c}\chi _0^c\\ \chi _0^n\end{array}\right)\delta U+\left(\begin{array}{cc}\chi _0^c& 2\chi _0^c\\ 2\chi _0^r& 2\chi _0^r\end{array}\right)\left(\begin{array}{c}g\delta n_0\\ g\delta \stackrel{~}{n}\end{array}\right),$$
(A1)
which gives the relationship
$$\chi =\frac{\delta n_0}{\delta U}+\frac{\delta \stackrel{~}{n}}{\delta U}=\frac{\chi _0^c+\chi _0^n+g\chi _0^c\chi _0^n}{(1g\chi _0^c)(12g\chi _0^n)4g^2\chi _0^c\chi _0^n}.$$
(A2)
Diagrammatic expressions for these different approximations are shown in Fig. 2.
## B Asymptotic Expansions
In this Appendix we derive the $`k\mathrm{\Lambda }_T0`$ asymptotic expansions for the functions $`\chi _0^n`$ and $`\mathrm{\Xi }`$, defined by Eqs. (8) and (20). These expansions are constructed to be valid for all $`\mu `$. We begin by breaking $`\chi _0^n`$ into two terms, one containing $`f(qk/2)`$ and one containing $`f(q+k/2)`$. After shifting $`q`$ by $`\pm k/2`$ and integrating out the angular variables, we have
$$\chi _0^n(k,\omega )=\frac{m}{(2\pi )^2k}_0^{\mathrm{}}𝑑qq^2f(q)\mathrm{log}\left[\frac{(\overline{p}+k/2q)(\overline{p}k/2+q)}{(\overline{p}+k/2+q)(\overline{p}k/2q)}\right],$$
(B1)
with $`\overline{p}=m\omega /k`$. We extract the important structure by rewriting the logarithm as an integral of the form $`𝑑x/x`$. Scaling all lengths by a multiple of the thermal wavelength, we find
$`\chi _0^n(k,\omega )`$ $`=`$ $`{\displaystyle \frac{m}{\pi k\mathrm{\Lambda }_T^2}}{\displaystyle _z_{}^{z_+}}𝑑zI(z)`$ (B2)
$`I(z)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑x{\displaystyle \frac{x}{xz}}{\displaystyle \frac{1}{e^{x^2\beta \mu }1}},`$ (B3)
where $`z_\pm =\omega /2\sqrt{\epsilon _kT}\pm k\mathrm{\Lambda }_T/4\sqrt{\pi }`$. The integral $`I(z)`$ has been characterized by Szépfalusy and Kondor . In particular, by expanding the distribution function in terms of Matsubara frequencies, one arrives at the $`z0`$ asymptotic expansion,
$`I(z)`$ $`=`$ $`\sqrt{\pi }g_{1/2}(e^{\beta \mu }){\displaystyle \frac{i\pi z}{2}}+i\pi \left({\displaystyle \frac{z}{z+i\sqrt{\beta \mu }}}\right)`$ (B5)
$`2\pi {\displaystyle \underset{j=1}{\overset{\mathrm{}}{}}}z^j{\displaystyle \underset{\nu =1}{\overset{\mathrm{}}{}}}\mathrm{}\left[(\beta \mu +2\pi i\nu )^{(j+1)/2}\right],`$
where $`\mathrm{}(z)`$ is the real part of $`z`$. Integration of the leading terms gives Eq. (10).
Following a similar procedure of integrating out the angular variables we write $`\mathrm{\Xi }`$ as
$$\mathrm{\Xi }(k,\omega )=\frac{2m}{\pi k\mathrm{\Lambda }_T^2}_{\overline{z}_{}}^{\overline{z}_+}𝑑zI(z),$$
(B6)
with $`\overline{z}_\pm =(k\mathrm{\Lambda }/4\sqrt{\pi })[i\sqrt{1+2\omega /\epsilon _k}\pm 1]`$. Using the expansion in Eq. (B5) gives Eq. (21).
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# Thermal Fluctuations of Elastic Filaments With Spontaneous Curvature and Torsion
## Abstract
We study the effects of thermal fluctuations on thin elastic filaments with spontaneous curvature and torsion. We derive analytical expressions for the orientational correlation functions and for the persistence length of helices, and find that this length varies non–monotonically with the strength of thermal fluctuations. In the weak fluctuation regime, the persistence length of a spontaneously twisted helix has three resonance peaks as function of the twist rate. In the limit of strong fluctuations, all memory of the helical shape is lost.
Recent advances in the art of micro–manipulation of molecules led to many experimental studies of the elasticity of biomolecules such as DNA , chromatin, proteins, and rodlike protein assemblies. The tacit assumption behind many of the theories is that the elasticity of these biomolecules is of entropic origin and, consequently, they are modeled as random walks. An alternative approach to the modeling of such systems is based on the assumption that the origin of elasticity is energy rather than entropy – there exists a lowest energy equilibrium configuration with associated spontaneous curvature , deviations from which give rise to restoring forces. While such an approach is a straightforward extension of the usual theory of elasticity of thin rods, the description of arbitrary spontaneous curvature and twist involves rather complicated differential geometry and most DNA–related studies considered only fluctuations around the straight rod configuration (see, however, references and ). Following recent studies on the elasticity and stability of thin rods with arbitrary spontaneous curvature and torsion, in this work we investigate the effect of thermal fluctuations on the statistical properties of such filaments. We derive the differential equations for the orientational correlation functions of the vectors pointing along the principal axes of the filament, and use them to calculate the correlators and the effective persistence length of an untwisted helix. Analytical expressions for the persistence length of a spontaneously twisted helix are obtained and it is found that this length varies non–monotonically with the amplitude of fluctuations and exhibits resonant–like dependence on the rate of twist. We would like to emphasize that although the present work is motivated by recent studies of biomolecules, its aim is to construct a theoretical framework for the description of fluctuating string–like objects that goes beyond current models of polymer physics, and we do not attempt here to model particular experiments involving single–molecule manipulation.
A filament of small but finite and, in general, non–circular cross–section, is modeled as an inextensible but deformable physical line parametrized by a contour length $`s`$ ($`0sL`$ where $`L`$ is the length of the filament). To each point $`s`$ one attaches a triad of unit vectors $`𝐭(s)`$ whose component $`𝐭_3`$ is the tangent vector to the curve at $`s`$, and the vectors $`𝐭_1(s)`$and $`𝐭_2(s)`$ are directed along the axes of symmetry of the cross–section. Note that $`𝐭(s)`$, together with the inextensibility condition $`d𝐱/ds=𝐭_3`$, gives a complete description of the space curve $`𝐱(s)`$, as well as of the twist of the cross–section about this curve. The rotation of the triad $`𝐭`$ as one moves along the curve is determined by the generalized Frenet equations
$$\frac{d𝐭_i(s)}{ds}=\underset{j}{}\mathrm{\Omega }_{ij}(s)𝐭_j(s)$$
(1)
Here $`\mathrm{\Omega }_{ij}(s)=_k\epsilon _{ijk}\omega _k(s)`$, $`\epsilon _{ijk}`$ is the antisymmetric tensor and $`\left\{\omega _k\right\}`$ are generalized torsions.
The theory of elasticity of thin rods is based on the notion that there exists a stress–free reference configuration defined by the set of spontaneous curvatures and torsions $`\left\{\omega _{0k}\right\}`$. The set $`\left\{\omega _{0k}\right\}`$ together with Eqs. (1) completely determines the equilibrium shape of the filament. Neglecting excluded–volume effects and other non–elastic interactions, the elastic energy associated with some actual configuration $`\left\{\omega _k\right\}`$ of the filament is a quadratic form in the deviations $`\delta \omega _k(s)=\omega _k(s)\omega _{0k}(s)`$
$$U_{el}\left\{\delta \omega _k\right\}=\frac{kT}{2}_0^L𝑑s\underset{k}{}a_k\delta \omega _k^2$$
(2)
where the bare persistence lengths $`a_1=E_1I_1/kT,`$ $`a_2=E_1I_2/kT`$ and $`a_3=E_2(I_1+I_2)/kT`$ ($`T`$ is the temperature and$`k`$ is the Boltzmann constant) are expressed in terms of the elastic moduli $`E_i`$ and the moments of inertia $`I_i`$ about the axes of symmetry of the cross–section. The only limitation on the validity of Eq. (2) is that deformations are small on microscopic length scales, of order of the thickness of the filament. The elastic energy $`U_{el}\left\{\delta \omega _k\right\}`$ determines the statistical weight of the configuration $`\left\{\omega _k\right\}`$. Calculating the corresponding Gaussian path integrals we find that $`\delta \omega _i(s)=0`$ and
$$\delta \omega _i(s)\delta \omega _j(s^{})=a_i^1\delta _{ij}\delta (ss^{})$$
(3)
We conclude that fluctuations of generalized torsions at two different points along the filament contour are uncorrelated, and that the amplitude of fluctuations is inversely proportional to the corresponding bare persistence length.
The statistical properties of fluctuating filaments are determined by the orientational correlation functions, $`𝐭_i(s)𝐭_j(s^{})`$. In order to derive a differential equation for this correlators, we calculate the variation of $`𝐭_i(s)`$ under the substitution $`ss+\mathrm{\Delta }s`$. Integrating Eq. (1) yields in matrix notation (for small $`\mathrm{\Delta }s`$):
$`𝐭(s+\mathrm{\Delta }s)`$ $`=`$ $`\{\mathrm{𝟏}{\displaystyle _s^{s+\mathrm{\Delta }s}}ds_1𝛀(s_1)+{\displaystyle \frac{1}{2}}{\displaystyle _s^{s+\mathrm{\Delta }s}}ds_1{\displaystyle _s^{s+\mathrm{\Delta }s}}ds_2𝛀(s_1)𝛀(s_2)+`$ (4)
$`{\displaystyle \frac{1}{2}}{\displaystyle _s^{s+\mathrm{\Delta }s}}ds_1{\displaystyle _s^{s_1}}ds_2[𝛀(s_1)𝛀(s_2)𝛀(s_2)𝛀(s_1)]+\mathrm{}\}𝐭(s)`$
where the last term appears because of noncommutativity of matrices $`𝛀(s_1)`$ and $`𝛀(s_2)`$ for different $`s_1`$ and $`s_2`$. We multiply the above expression by $`𝐭_j(s^{})`$, average the result and note that for $`s+\mathrm{\Delta }s>s>s^{}`$ the fluctuations $`\delta \omega _i(s_1)`$ and $`\delta \omega _j(s_2)`$ are uncorrelated with the fluctuations of $`𝐭_i(s)`$ and $`𝐭_j(s^{})`$. This implies that averages of products of $`𝛀`$’s and $`𝐭`$’s factorize into products of the averages of $`𝛀`$’s and those of $`𝐭`$’s. Since the averages of the terms in the square brackets in Eq. (4) depend only on $`|s_1s_2|`$ and their difference vanishes, in the limit $`\mathrm{\Delta }s0`$ this yields
$$\frac{}{s}𝐭_i(s)𝐭_j(s^{})=\underset{k}{}\mathrm{\Lambda }_{ik}(s)𝐭_k(s)𝐭_j(s^{})$$
(5)
where
$$\mathrm{\Lambda }_{ik}=\gamma _i\delta _{ik}+\underset{l}{}\epsilon _{ikl}\omega _{0l}\text{ with }\gamma _i=\underset{k}{}\frac{1}{2a_k}\frac{1}{2a_i}$$
(6)
The above equations describe the fluctuations of filaments of arbitrary shape and flexibility and in the following this general formalism is applied to helical filaments.
Consider a helix without spontaneous twist, such that the generalized spontaneous curvatures and torsions $`\left\{\omega _{0k}\right\}`$ are independent of position $`s`$ along the contour. In order to visualize the stress–free configuration of such a filament, it is convenient to introduce the conventional Frenet triad of unit vectors which consists of the tangent, normal and binormal to the space curve spanned by the centerline, supplemented by a constant twist angle $`\alpha _0`$ which describes the rotation of the cross–section about this line. The relation between the two triads is given by $`\omega _{01}=\kappa _0\mathrm{cos}\alpha _0,`$ $`\omega _{02}=\kappa _0\mathrm{sin}\alpha _0`$ and $`\omega _{03}=\tau _0`$, where $`\kappa _0`$ and $`\tau _0`$ are the constant curvature and torsion of the space curve. The rate of rotation of the centerline about the long axis of the helix is $`\omega _0=(\kappa _0^2+\tau _0^2)^{1/2}`$ , the helical pitch is $`2\pi \tau _0/\omega _0^2`$ and the radius is $`2\pi \kappa _0/\omega _0^2`$. For constant $`\{\kappa _0,\tau _0,\alpha _0\}`$, $`𝚲`$ is a constant matrix and $`𝐭_i(s)𝐭_j(s^{})`$ is given by $`ij`$th element of the matrix $`\mathrm{exp}\left[𝚲\left(ss^{}\right)\right]`$. The eigenvalues of the matrix $`𝚲`$ can be obtained analytically by solving for the roots of a characteristic cubic polynomial but the resulting expressions are cumbersome and will be presented in a longer report. Here we discuss only two limiting cases:
Weak fluctuations, $`_i\gamma _i\omega _0`$ In this case
$$𝐭_i(s)𝐭_i(0)=\left(\omega _{0i}^2/\omega _0^2\right)e^{s/l}+\left(1\omega _{0i}^2/\omega _0^2\right)\mathrm{cos}(\omega _0s)e^{s/2ls/2l_\varphi }$$
(7)
with the decay lengths $`l`$ and $`l_\varphi `$ defined by $`l^1=_k\gamma _k\omega _{0k}^2/\omega _0^2`$ and $`l_\varphi ^1=_ka_k^1\omega _{0k}^2/\omega _0^2`$. The physical meaning of this correlator becomes clear by switching off thermal fluctuations ($`\gamma _k=0`$). The first term on the rhs of this equation expresses the fact that the projection of any vector $`𝐭_i`$ of the triad on the symmetry axis of the helix is constant, with magnitude $`\omega _{0i}/\omega _0`$. The projections on the plane normal to this axis with magnitudes $`(1\omega _{0i}^2/\omega _0^2)^{1/2},`$ rotate with angular rate $`\omega _0`$. In the presence of weak fluctuations, the axis of symmetry of the helix becomes a random walk and the loss of correlations of its projections along the axes of the triad is described by the factor $`\mathrm{exp}(s/l)`$. In the second term of Eq. (7), $`\mathrm{exp}(s/2l)`$ describes the loss of correlations of the orientation of the plane normal to the axis of the helix. The angular persistence length $`l_\varphi `$ results from averaging over the random angle of rotation ($`\varphi `$) with respect to the axis of the helix, $`\mathrm{cos}[\omega _0s+\varphi (s)\varphi (0)]=\mathrm{cos}\left(\omega _0s\right)\mathrm{exp}\left(s/2l_\varphi \right),`$ with $`\left[\varphi (s)\varphi (0)\right]^2=s/l_\varphi `$.
Strong fluctuations, $`_i\gamma _i\omega _0`$. In this limit
$$𝐭_i(s)𝐭_j(0)=e^{\gamma _i\left(ss^{}\right)}\delta _{ij}$$
(8)
i.e., the correlators depend only on the bare persistence lengths and memory about the orientation of the vector $`𝐭_i`$ decays over contour distance $`\gamma _i^1`$. Strong fluctuations destroy all phase coherence and all correlations between different vectors of the triad and lead to complete “melting” of the helix.
We now proceed to calculate the effective persistence length $`l_p`$ which controls both the thermal fluctuations of a filament and its elastic response to external forces. It is defined as the ratio of the rms end–to–end separation $`R^2`$ and the contour length of the filament $`L`$, in the limit $`L\mathrm{}`$. The end–to–end vector is defined as $`𝐑=_0^L𝐭_3(s)𝑑s`$ and thus
$$l_p=\underset{L\mathrm{}}{lim}\frac{2}{L}_0^L𝑑s_0^s𝑑s^{}𝐭_3(s)𝐭_3(s^{})$$
(9)
A simple calculation yields a result valid for arbitrarily strong fluctuations:
$$l_p=2\frac{\tau _0^2+\gamma _1\gamma _2}{\kappa _0^2\left(\gamma _1\mathrm{cos}^2\alpha _0+\gamma _2\mathrm{sin}^2\alpha _0\right)+\tau _0^2\gamma _3+\gamma _1\gamma _2\gamma _3}$$
(10)
For non–vanishing curvature and torsion, this expression diverges in the weak fluctuation limit $`\gamma _i0`$ and the shape of the filament is nearly unaffected by fluctuations. Non–monotonic behavior is observed for “plate–like” helices, with large radius to pitch ratio, $`\kappa _0/\tau _0`$. For $`\gamma _i0`$, the effective persistence length approaches zero (recall that $`L\mathrm{}`$ in Eq. (9)). Thermal fluctuations expand the helix by releasing stored length and initially increase the persistence length. Eventually, in the limit of strong fluctuations, the persistence length vanishes again (as $`\gamma _3^1)`$ because of complete randomization of the filament. The sensitivity to the constant angle of twist $`\alpha _0`$ increases with radius to pitch ratio. In the opposite limit of “rod–like” helices $`\kappa _00,`$ the effective persistence length approaches $`2/\gamma _3`$ and becomes a function of $`a_1`$ and $`a_2`$ only. Indeed, since straight inextensible rods do not have stored length, their end–to–end distance and persistence length are determined by random bending and torsional fluctuations only, and are independent of twist.
The preceding analysis can be extended to fluctuating filaments with twisted stress–free states characterized by constant curvature $`\kappa _0,`$ torsion $`\tau _0`$ and rate of twist of the cross–section about the centerline, $`d\alpha _0/ds`$. The relation between generalized torsions $`\left\{\omega _{0k}(s)\right\}`$ and Frenet parameters $`\{\kappa _0,\tau _0,\alpha _0(s)\}`$ is given by $`\omega _{01}=\kappa _0\mathrm{cos}\alpha _0,`$ $`\omega _{02}=\kappa _0\mathrm{sin}\alpha _0`$ and $`\omega _{03}=\tau _0+d\alpha _0/ds`$. The calculation of the persistence length involves the solution of Eq. (5) with periodic coefficients. Details of the analytical solution will be given elsewhere. For filaments with circular cross–section $`a_1=a_2`$, the persistence length is independent of twist. In Fig. 1 we present a plot of the persistence length given in units of the helical pitch $`l^{}=l_p\omega _0^2/2\pi \tau _0,`$ on the dimensionless rate of twist $`w^{}=2\omega _0^1d\alpha _0/ds`$, for a “plate–like” helix with large radius to pitch ratio $`\kappa _0/\tau _0`$ and ribbon–like cross–section, $`a_1a_2`$. Curve 1 corresponds to the case of weak fluctuations, $`\gamma _i\omega _0.`$ Throughout most of the range, the persistence length is independent of the rate of twist but a sharp peak appears at $`d\alpha _0/ds=0`$ (see insert), accompanied by two smaller peaks at $`d\alpha _0/ds=\pm \omega _0/2`$. Note that while in the limit of vanishing pitch, a ribbon–like untwisted helix degenerates into a normal ring, the cross–section of a twisted helix with $`d\alpha _0/ds=\pm \omega _0/2`$ rotates by $`\pm \pi `$ and the helix degenerates into a Mőbius ring. As the amplitude of fluctuations increases, the central peak transforms into a narrow dip and the two Mőbius peaks become broad minima (curve 2). Further increase of fluctuations leads to the disappearance of the Mőbius dips and the central dip becomes broad and shallow (curve 3). Finally when $`\gamma _i\omega _0`$, all dependence of the persistence length on the spontaneous twist disappears (curve 4). Note that, as expected from the discussion following Eq. (10), the persistence length depends non–monotonically on the amplitude of thermal fluctuations.
In order to understand the origin of the Mőbius resonances we note that while the effective persistence length is a property of the space curve given by the Frenet triad, the microscopic Brownian motion of the filament arises as the result of random forces that act on its cross–section and therefore are given in the frame associated with the principal axes of the filament. Since the two frames are related by a rotation of the cross–section by an angle $`\alpha _0(s)`$, the random force in the Frenet frame is modulated by linear combination of $`\mathrm{sin}\alpha _0(s)`$ and $`\mathrm{cos}\alpha _0(s)`$. This gives a deterministic contribution to the persistence length which, to lowest order in the force, is proportional to the mean square amplitude of the random force and therefore varies sinosoidally with $`\pm 2\alpha _0(s)`$. The observed resonances occur whenever the natural rate of rotation of the helix $`\omega _0`$ coincides with the rate of variation of this deterministic contribution of the random force, $`\pm 2d\alpha _0/ds`$.
In summary, we presented a statistical mechanical description of thermally fluctuating elastic filaments of arbitrary shape and flexibility. We would like to emphasize that the only limitation on the magnitude of fluctuations is that they are small on microscopic length scales, and that our model describes arbitrary deviations of a long filament from its equilibrium shape. The general formalism was applied to helical filaments with and without twist. Strong thermal fluctuations lead to melting of the helix, accompanied by a complete loss of helical correlations. In general, the persistence length is a non–monotonic function of the amplitude of thermal fluctuations. Although through most of its range twist has a minor effect on the persistence length, resonant peaks and dips are observed whenever the rate of twist approaches zero or equals in absolute magnitude to half the rate of rotation of the helix.
Acknowledgment We would like to thank A. Drozdov for illuminating discussions. YR acknowledges support by a grant from the Israel Science Foundation.
Figure captions
Plot of the dimensionless persistence length $`l^{}`$ as a function of the dimensionless rate of twist $`w^{}`$ for a helical filament with spontaneous curvature $`\kappa _0=1`$, and torsion $`\tau _0=0.01`$ (in arbitrary units). The different curves correspond to different bare persistence lengths: (1) $`a_1=100,`$ $`a_2=a_3=5000`$, (2) $`a_1=1,`$ $`a_2=a_3=100`$, (3) $`a_1=0.1,`$ $`a_2=a_3=10`$, (4) $`a_1=0.01,`$ $`a_2=a_3=10.`$ A magnified view of the region of small twist rates is shown in the insert.
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# IR-TF Relation in the Zone of Avoidance with 2MASS
## 1. Introduction. The Tully-Fisher Relation
In 1977, Tully and Fisher discovered that there is a good correlation between maximum rotational velocity, usually given by the global neutral hydrogen line profile width $`\mathrm{\Delta }V`$ (corrected for inclination) and the absolute magnitude $`\mathrm{M}`$ for late-type galaxies (Tully & Fisher 1977).
Observationally, the correlation (luminosity vs. maximum velocity) gives $`L\mathrm{\Delta }V^\alpha `$ with $`\alpha `$ in the range of 2.5-3 at optical wavelengths. There is a general trend that at longer wavelengths, the slope steepens. At infrared wavelengths, the slope ($`2.5\times \alpha `$ in magnitudes) climbs to about $`10`$ (or $`\alpha 4`$) which is close enough to the Faber-Jackson relation $`L\sigma _v^4`$ to suspect a possible common physical origin.
The origin of the TF relation lies (in the very roughest terms) in the virial theorem. Assuming a constant mass-to-light ratio ($`M/L`$) and constant surface brightness (I), one can rewrite the virial theorem as
$$LV^4[(M/L)^2I]^1$$
(1)
The zero-point is directly related to basic physical properties as can be seen if rewritten as $`2.5\times \mathrm{log}[I(M/L)^2]`$ or $`\mathrm{log}[(M/L)\mathrm{\Sigma }]`$. $`\mathrm{\Sigma }`$ is the surface mass density which is a strong function of the angular momentum ($`\mathrm{\Sigma }M^7/J^4`$) whereas the mass-to-light ratio gives information on the stellar content. This can be used to put strong constraints on galaxy formation theory (see Steinmetz & Navarro 1999; Koda et al. 2000).
Zwaan et al. 1995 also have shown that low surface brightness late-type galaxies follow the same B-band TF relation as the high surface brightness galaxies, even though their mass-to-light ratios are very different as well as the ratio of dark to luminous matter (de Blok & McGaugh 1997), but this is still somewhat controversial (see O’Neil et al. 2000).
The HST Key Project group on the extragalactic distance scale has used the TF relation to infer the Hubble constant ($`\mathrm{H}_0`$) (see Sakai et al. 2000 for a recent discussion). Another application is to use the TF relation to measure the peculiar velocity field of galaxies. This provides an estimate of the mass density on large scales and of the density of the universe through $`\mathrm{\Omega }^{0.6}/b`$.
Using near-IR for TF studies has two advantages: (1) it is sensitive to old red stars and less sensitive to young blue stars. In other words, it is fairly independent of current star formation, so it provides a measure of current stellar content that is, presumably, more closely tied to the total mass of the galaxy and hence its kinematics; (2) uncertainties due to interstellar absorption in both our own Galaxy and in the observed galaxy are greatly reduced.
This project shows the feasibility of using 2MASS (Huchra et al. these proceedings; Jarrett et al. 2000) and H I data to study large scale structure. We studied the Pisces-Perseus (PP) supercluster region and the Great Attractor (GA) region in the Zone of Avoidance. In particular, we ask the following questions regarding 2MASS:
(i) Which band (J, H or K) generates the best TF relation?
(ii) What magnitude measurement gives the tightest TF relation?
(iii) Can we use 2MASS data at low Galactic latitude?
and regarding the TF relation:
(iv) What is the highest axial ratio $`b/a`$ that we can use?
(v) Is there any internal extinction?
The remainder of this paper is organized as follows. In section 2, we present our data sets. The method is described in section 3. Our results are presented in section 4 and summarized in section 5.
## 2. Data Sets
We construct two data sets. The first (PP) data set consists of H I data selected from the Huchtmeier & Richter catalog (Huchtmeier & Richter 1989). Because we found substantial scatter was introduced when using line widths measured in different ways, we select only papers with Giovanelli and collaborators. The data was also restricted to the Arecibo telescope and measurements made at 50% of the peak. The sample contains $`2700`$ objects. We cross-matched 2MASS data with the Giovanelli subsample. There were 502 matches as of spring 1999 <sup>1</sup><sup>1</sup>1As of January 2000, there are about 3 times as many matches, but the analysis is not completed.. This sample covers the sky from $`0^\mathrm{h}`$ to $`5^\mathrm{h}`$ in right ascension and $`15^{}`$ to $`35^{}`$ in declination which corresponds to the Pisces-Perseus supercluster region.
We also used H I data from Giovanelli et al. 1997 from their extensive survey of the I-band TF in clusters. We again cross-matched this subsample with the 2MASS database and obtained 300 matches from 22 different clusters (hereafter cluster sample). We took into account the different offsets from the Hubble flow (column (4) in Table 2 of Giovanelli et al. 1997) for each cluster. This sample also contains inclination information both from I-band photometry and from 2MASS which enabled us to perform checks on the 2MASS estimate of the axial ratio.
## 3. Method and Error Budget
### 3.1. Sub-Sample
In most cases, we used a subsample of galaxies chosen to have the following properties: (i) axial ratio $`b/a`$ less than 0.5 (the axial ratio used was the “super coadd” axial ratio, keyword sup\_ba in the 2MASS database, which is determined from the coadded J, H and K images); (ii) H I flux minimum of 1 Jy km s<sup>-1</sup> in order to remove low SNR H I data; (iii) magnitude cutoff of K=13 mag (in the PP sample, we had a cutoff on the magnitude errors ($`\sigma =0.1`$) which corresponds to a cutoff in magnitude of $`13`$ determined from the dependence of the magnitude errors with magnitude); (iv) excluding S0’s and ellipticals (in general, the $`b/a`$ constraint and H I detection are sufficient to ensure that no ellipticals are in the subsample), and (v) redshifts between $`z_{min}`$ and $`z_{max}`$ to examine galaxies in different distance ranges.
Figure 1 shows the TF relation for the cluster sample.
### 3.2. Corrections
In order to improve the quality of the TF relation, a number of corrections are usually made for turbulent motion, instrumental broadening, and extinction. These corrections have a number of adjustable parameters that we investigate here to determine which values or methods work best with the 2MASS data.
We have corrected for (i) instrumental broadening following Bottinelli et al. (1990) and Giovanelli et al. (1997); (ii) turbulent motion according to the scheme proposed by Tully & Fouqué 1985 (this correction is small and did not improve the scatter of the TF relation); (iii) inclination, i.e. $`v_{rot}=\frac{W_R}{2\mathrm{sin}i}`$ where $`W_R`$ is the observed velocity width $`W_{50}`$ corrected for instrumental broadening and $`i`$ is the disk inclination based on the axial ratio $`b/a`$ from 2MASS, and (iv) photometric corrections such as the Galactic extinction (from COBE/ DIRBE maps, Schlegel et al. 1998) and internal extinction $`\mathrm{\Delta }m=\gamma \mathrm{log}b/a`$ where $`\gamma `$ characterizes the extinction for a highly inclined galaxy. In general, $`\gamma `$ can be a function of galaxy luminosity or velocity width (see Sakai et al. 2000 and references therein), but we treated it as a constant.
### 3.3. Parameterization and Error Budget
We parameterize the TF relation in the same way as Giovanelli et al. (1997), $`y=\mathrm{a}+\mathrm{b}x`$ where $`y`$ is the absolute magnitude and $`x`$ is the logarithm of the rotation velocity, i.e.:
$`x`$ $`=`$ $`\mathrm{log}\left[{\displaystyle \frac{W_R}{2\mathrm{sin}i}}\right]2.5`$ (2)
$`y`$ $`=`$ $`(m_{obs}A\mathrm{\Delta }mk_z5\mathrm{log}{\displaystyle \frac{cz}{100\mathrm{h}_{100}}}25)5\mathrm{log}\mathrm{h}_{100}`$ (3)
$`=`$ $`M5\mathrm{log}\mathrm{h}_{100}`$
where h<sub>100</sub> is the Hubble constant in units of 100 km s<sup>-1</sup> Mpc<sup>-1</sup>. The last term $`5\mathrm{log}\mathrm{h}_{100}`$ makes $`y`$ independent of the Hubble constant so it is easier to compare results with other authors. $`k_z`$ is the $`k`$-correction and was generally ignored here since the sample contains only nearby galaxies ($`cz<10,000`$ km s<sup>-1</sup>).
A standard $`\chi ^2`$ fit to the data was performed using errors both along the $`x`$ and $`y`$ axes, $`\sigma _x`$ and $`\sigma _y`$ (bivariate fit). $`\sigma _x`$ includes velocity width uncertainty as well as inclination errors, which dominate the overall error. $`\sigma _y`$ includes magnitude uncertainty from 2MASS, distance uncertainty, and an added intrinsic scatter $`\sigma _i`$ which takes into account any “cosmological” scatter inherent in the TF relation. The $`1\sigma `$ uncertainty on the parameters (a and b) is given by the projection of the $`1\sigma `$ contour semi-major axis onto the parameter axis (Press et al. 1992).
## 4. Results and Discussion
2MASS produces many different magnitude measurements for each object (Jarrett et al. 2000). The different types of magnitude that we used are (i) the 20th mag/arcsec<sup>2</sup> isophotal magnitude (i20 in the 2MASS database <sup>2</sup><sup>2</sup>2 In order to find the proper keyword in the 2MASS database, one needs to add the prefix j\_m\_, h\_m\_ or k\_m\_ and the suffix e (c) for elliptical (circular) aperture. ); (ii) the 21st mag/arcsec<sup>2</sup> isophotal magnitude (i21); (iii) the 20th mag/arcsec<sup>2</sup> K fiducial <sup>3</sup><sup>3</sup>32MASS has two kinds of magnitude type, “individual” (i20 and j21 here) and “fiducial” (k20f, j21f and kf here) . “Fiducial” means the aperture used for the photometry was selected at one band and applied to the other two. magnitude (k20f); (iv) the 21st mag/arcsec<sup>2</sup> J fiducial magnitude (j21f), and (v) the Kron (Kron 1980) K fiducial (kf) magnitude which measures the flux within an aperture 2$`r_1`$ where $`r_1`$ is the first moment of the light distribution. This type of magnitude is thought to be less sensitive to observing conditions and is thought to be a “total” measure of the integrated flux (Koo 1986).
First, which band J, H or K produces the least scatter in the TF relation? Figure 2 shows the minimum value $`\chi _{min}^2`$ of the bivariate fit for each of the magnitude types. Figure 2 indicates that (1) the 20th mag/arcsec<sup>2</sup> K fiducial magnitude and Kron K fiducial produces the smallest scatter; (2) in general K gives a tighter TF relation than J or H, considering that H-band photometry of 2MASS is subject to uncertainty due to air glow; and (3) elliptical apertures tend to give a better TF relation than circular ones, although this was not consistently true for all datasets.
Since the different magnitude types have different measurement errors, the $`\chi _{min}^2`$ could look artificially low if the error estimates were inflated. However, we found no effects of this sort. Magnitude types with the largest mean measurement error did not produce the smallest $`\chi ^2`$.
Figure 3 shows the fitted values of the slope, a, and the zero-point, b, for each of the magnitude types assuming an intrinsic scatter $`\sigma _i`$ of 0.2 for the cluster sub-sample and using the 2MASS axial ratios for the inclination. First note that the zero-point of the TF relation is well determined and does not depend on the magnitude used. Second, the slope is more dependent on the magnitude type. However, the “fiducial” magnitudes are well within 1-$`\sigma `$ of each other. Third, the slope tends to be steeper as the wavelength increases. This is often called the color-TF relation. The same conclusions can be reached using either the 2MASS or I-band inclinations and different values of $`\sigma _i`$.
Figure 4 shows the value of the reduced $`\chi ^2`$ as a function of $`\gamma `$ for the PP sample using the magnitude k20fc and an intrinsic scatter of 0.2 mag. We find a minimum for each band: $`0.5`$ mag at J, $`0.3`$ at H and $`0.1`$ at K, although the amount of extinction at K is consistent with none. These curves are consistent with our results from the cluster sample, however, they are not fully reproduced when using the inclination from the I-band photometry. Our results are consistent with a reddening of E(B-V) of 0.5 mag which gives a $`A_\lambda `$ of 0.45 mag, 0.28 and 0.18 for the three bands.
We did not see any correlation with type for the cluster data set while the PP data set shows a slight correlation. A linear fit gives a non-zero slope ($`0.06`$) with 1-$`\sigma `$ uncertainty of 0.03 for the PP data set, such that both results are consistent with no type dependence (at the 95% confidence level).
We have not tested for dependence of the TF relation on surface brightness because the current 2MASS galaxy processor is not very sensitive to low surface brightness disks. Attempts are underway to extract them in the future.
## 5. Conclusions
From our two samples (PP and cluster data sets), we found that using K-band is generally preferable over the other bands. The K-fiducial 20th mag/arcsec<sup>2</sup> (and the K-fiducial Kron, although in the PP sample, the Kron K-fiducial does not produce such a low $`\chi ^2`$) seems to produce the least scatter. This is consistent with the K-fiducial isophotal elliptical aperture magnitude being the most robust photometric measurement (Jarrett et al. 2000). At low Galactic latitude, preliminary results suggests one would need optical axial ratios because the high star density can mislead the 2MASS ellipse-fitting program. The zero-point is well determined and is even independent of the magnitude types and of the sample. The TF slope shows more scatter but it steepens with wavelength. An interesting result is that we find internal extinction at each band, although the amount of extinction at K is consistent with none. We find no type dependence in our samples but we cannot yet conclude whether there is a surface brightness dependence.
Our results are summarized below for the PP sample (equations 4 & 5) and the cluster sample (equations 6 & 7) using an intrinsic scatter of $`\sigma _i=0.2`$ and assuming $`\gamma =0.1`$ mag of internal extinction.
$`M_K5\mathrm{log}\mathrm{h}_{100}`$ $`=`$ $`(9.90\pm 0.35)(\mathrm{log}v_{rot}2.5)+(24.70\pm 0.10)`$ (4)
$`M_K5\mathrm{log}\mathrm{h}_{100}`$ $`=`$ $`(9.64\pm 0.35)(\mathrm{log}v_{rot}2.5)+(24.76\pm 0.09)`$ (5)
$`M_K5\mathrm{log}\mathrm{h}_{100}`$ $`=`$ $`(9.63\pm 0.27)(\mathrm{log}v_{rot}2.5)+(24.74\pm 0.08)`$ (6)
$`M_K5\mathrm{log}\mathrm{h}_{100}`$ $`=`$ $`(9.58\pm 0.29)(\mathrm{log}v_{rot}2.5)+(24.86\pm 0.10)`$ (7)
Equations 4 & 6 are for the k20fc magnitude, while equations 5 & 7 are for the k20fe magnitude. Both samples give very similar results.
As part of future work to study large scale flows in the ZOA, we have recently collected H I data at Arecibo for low Galactic latitude galaxies. Of 169 sources at $`|b|<10^{}`$, 72 were detected; of 147 sources at $`10{}_{}{}^{}<|b|<20^{}`$, 75 were detected. The galaxies were identified using the 2MASS galaxy processor, which holds promise for identifying a large sample of galaxies within the ZOA.
Acknowledgments N. Bouché is the recipient of a Five College Astronomy Graduate Fellowship supported by the Mary Dailey Irvine Fund which allowed him to attend the meeting.
## References
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de Blok, W.J.G., & McGaugh, S.S. 1997, MNRAS, 290, 533
Giovanelli, R., Haynes, M.P., Herter, T., & Vogt, N. 1997, AJ, 113, 53
Huchtmeier, W.K., & Richter, O.G. 1989, A&A, 210, 1
Jarrett, T.H., Chester, T., Cutri, R., Schneider, S., Skrutskie, M., & Huchra, J.P. 2000, AJ, 119, 2498
Koda, J., Sofue, Y., & Wada, K. 2000, ApJ, 531, L17
Koo, D.C. 1986, ApJ, 311, 651
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O’Neil, K., Bothun, G.D., & Schombert, J. 2000, AJ, 119, 1360
Press, W.H., Teukolsky, S.A., Vetterling, W.T., & Flannery, B. 1992, Numerical Recipes in C, (Cambridge: Cambridge University Press)
Sakai, S., Mould, J.R., Hughes, S.M.G., et al. 2000, ApJ, 529, 698
Schlegel, D.J., Finkbeiner, D.P., & Davis, M. 1998, ApJ, 500, 525
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# Long-range interaction of two metastable rare-gas atoms
## Abstract
We present semiempirical calculations of long-range van der Waals interactions for two interacting metastable rare-gas atoms Ne through Xe. Dispersion coefficients $`C_6`$ are obtained for homonuclear molecular potentials asymptotically connecting to the $`ns(3/2)_2+ns(3/2)_2`$ atomic states. The estimated uncertainty of the calculated $`C_6`$ dispersion coefficients is 4%.
Motivated by cold-collision studies of metastable rare-gas atoms and prospects of achieving Bose-Einstein Condensation in these systems , we present calculations of long-range dispersion (van der Waals) coefficients for two atoms interacting in the $`ns(3/2)_2`$ atomic states ($`n=3`$ for Ne, $`n=4`$ for Ar, $`n=5`$ for Kr, and $`n=6`$ for Xe). The metastable states have long lifetimes, 43 sec for Xe , decaying to the ground $`{}_{}{}^{1}S_{0}^{}`$ state by a weak magnetic-quadrupole transition. With such a long lifetime the metastable atom behaves as an effective ground state in experiments. Compared to alkali-metal systems, an attractive feature of the noble gas atoms is the availability of isotopes with zero nuclear spin. The lack of hyperfine structure leads to a substantial simplification of molecular potentials, though some complexity arises due to the nonvanishing total electron angular momentum ($`J`$=2) of the metastable state. The anisotropy leads to fifteen distinct long-range molecular states connecting to the $`ns(3/2)_2+ns(3/2)_2`$ asymptotic configuration.
Our theoretical treatment of long-range interactions is similar to recent high-precision calculations of van der Waals coefficients for alkali-metal atoms . By using many-body methods and accurate experimental matrix elements for the principal transitions, leading dispersion coefficients $`C_6`$ were determined to an accuracy better than 1% for Na, K, and Rb, and of 1% for Cs and 1.5% for Fr. The semiempirical values of $`C_6`$ coefficients for metastable noble-gas atoms obtained here have an estimated uncertainty of 4%. The approach relies on the determination of dynamic polarizability functions. To construct the polarizabilities we combine experimental lifetime and energy data of the excited states with accurate semiempirical dynamic polarizabilities of the ground states of noble-gas atoms . The theoretical lifetimes and branching ratios are adjusted to reproduce the measured static polarizabilities , which are known with a 2% uncertainty. We estimated the additional small contributions within the Dirac-Hartree-Fock framework. The resulting polarizabilities satisfy the Thomas-Reiche-Kuhn oscillator strength sum rule.
The Racah notation for atomic levels is used. The particle-hole states are labeled as $`n\mathrm{}(K)_J`$ or $`n\mathrm{}^{}(K)_J`$, where $`n`$ and $`\mathrm{}`$ are the principal and the orbital angular momentum quantum numbers of the valence electron and $`𝐊=𝐉_c+\mathrm{}`$, where $`J_c`$ is the angular momentum of the core. The primed configurations converge to a Rydberg series limit with a hole in the $`(n1)p_{1/2}`$ state, and the unprimed to a hole in the $`(n1)p_{3/2}`$ state. The manifold of the lowest $`ns`$ valence states has four fine-structure states $`ns^{}(1/2)_{0,1}`$ and $`ns(3/2)_{1,2}`$, and the lowest $`np`$ manifold consists of ten states. We investigate here the molecular potentials asymptotically connecting to the $`ns(3/2)_2`$ atomic states.
We calculate the long-range molecular potentials in the framework of Rayleigh-Schroedinger perturbation theory. The basis functions are defined as products of atomic wavefunctions
$$|M_1M_2;\mathrm{\Omega }=|ns(3/2)_2M_1_1|ns(3/2)_2M_2_2,$$
(1)
where the index 1(2) describes the wavefunction located on the center 1(2) and $`\mathrm{\Omega }=M_1+M_2`$, $`M_{1,2}`$ being projections of the atomic total angular momentum on the internuclear axis. Due to the axial symmetry of a dimer $`\mathrm{\Omega }`$ is a conserved quantum number. It takes values ranging from zero to four. The two-atom basis (1) is degenerate and the correct molecular wavefunctions are obtained by diagonalizing the molecular Hamiltonian
$$\widehat{H}=\widehat{H}_1+\widehat{H}_2+\widehat{V}(R).$$
(2)
In expression (2) $`\widehat{H}_k`$ represent the Hamiltonians of the two non-interacting atoms, and $`\widehat{V}(R)`$ is the interaction potential at an internuclear distance $`R`$. The energy of the $`ns(3/2)_2`$ metastable state is designated as $`^{}`$. Then in the model space (1)
$`\left(\widehat{H}_1+\widehat{H}_2\right)|M_1M_2;\mathrm{\Omega }=2^{}|M_1M_2;\mathrm{\Omega }.`$
The residual electrostatic potential $`\widehat{V}(R)`$ is defined as the full Coulomb interaction energy in the dimer excluding interactions of the atomic electrons with their parent nuclei.
The multipole interactions ($`L=1`$ for dipole-dipole, and $`L=2`$ for quadrupole-quadrupole interactions) are given by
$$V_{LL}(R)=\frac{1}{R^{2L+1}}\underset{\mu =L}{\overset{L}{}}\frac{(2L)!}{(L\mu )!(L+\mu )!}\left(T_\mu ^{(L)}\right)_1\left(T_\mu ^{(L)}\right)_2,$$
(3)
with the multipole spherical tensors
$$T_\mu ^{(L)}=|e|\underset{i}{}r_i^LC_\mu ^{(L)}(\widehat{𝐫}_i),$$
(4)
where the summation is over atomic electrons, $`𝐫_i`$ is the position vector of electron $`i`$, and $`C_\mu ^{(L)}(\widehat{𝐫}_i)`$ are reduced spherical harmonics . In the following we write $`d_\mu =T_\mu ^{(1)}`$ and $`Q_\mu =T_\mu ^{(2)}`$.
The lowest-order contribution to the term energies arises from the quadrupole-quadrupole interaction $`\widehat{V}_{qq}`$, which varies as $`1/R^5`$. However, the corresponding $`C_5`$ coefficients are only of the order $`10^110^2`$ a.u. , and the dominant contribution appears in the second order in $`\widehat{V}(R)`$, arising from the dipole-dipole interaction $`\widehat{V}_{dd}`$. The second-order dipole interaction is proportional to $`1/R^6`$, and the associated dispersion coefficient $`C_6`$ is of the order of $`10^410^5`$ a.u.. Applying the formalism of degenerate perturbation theory in second order , we obtain an effective Hamiltonian within the two-atom basis Eq. (1)
$$m|H_{\mathrm{eff}}^{\left(2\right)}|n=2^{}\delta _{mn}+m|\widehat{V}_{qq}|n+\underset{\mathrm{\Psi }_i}{}\frac{m|\widehat{V}_{dd}|\mathrm{\Psi }_i\mathrm{\Psi }_i|\widehat{V}_{dd}|n}{2^{}E_i}.$$
(5)
The intermediate molecular state $`|\mathrm{\Psi }_i`$ with unperturbed energy $`E_i`$ runs over a complete set of two-atom states, excluding the model-space states Eq. (1). The formalism of the generalized Bloch equation would allow the inclusion in the model space of the other three atomic states in the $`ns`$ manifold and would account for the mixing of the different fine-structure levels; but such a large model space is not necessary for $`R>10`$ a.u.. The position of the avoided level crossing can be estimated from $`R_{l.c.}2C_6/(_{ns(3/2)_1}_{ns(3/2)_2})10`$ a.u.. The $`\mathrm{\Omega }=4`$ molecular term is unique and being unaffected by avoided crossings, the region of applicability is extended to $`Rn`$ a.u., before the electronic clouds start to overlap. The effect of the quadrupole-quadrupole interaction on the term energy can be disregarded at values of $`RC_6/C_510^3`$ a.u.. The quadrupole-quadrupole correction is discussed by Doery et al. .
Using the Wigner-Eckart theorem, we can represent the matrix element of the dipole-dipole term in the effective second-order Hamiltonian as
$`{\displaystyle \underset{\mathrm{\Psi }_i}{}}{\displaystyle \frac{M_1M_2;\mathrm{\Omega }|\widehat{V}_{dd}|\mathrm{\Psi }_i\mathrm{\Psi }_i|\widehat{V}_{dd}|M_1^{}M_2^{};\mathrm{\Omega }}{2^{}E_i}}=`$ (7)
$`{\displaystyle \frac{1}{R^6}}{\displaystyle \underset{II^{}}{}}C_6^{J_aJ_b}(1)^{J_a+J_b}{\displaystyle \frac{2}{3}}{\displaystyle \underset{\lambda \mu }{}}w_\lambda ^1w_\mu ^1𝒜_{\lambda \mu }^{J_a}(M_1,M_1^{})𝒜_{\lambda \mu }^{J_b}(\mathrm{\Omega }M_1,\mathrm{\Omega }M_1^{}).`$
The dipole weights $`w_\mu ^1`$ are $`w_{+1}^1=w_1^1=1`$, and $`w_0^1=2`$. $`J_a`$ and $`J_b`$ are the corresponding total angular momenta of intermediate atomic states of atoms 1 and 2, and
$$𝒜_{\lambda \mu }^I(M_1,M_1^{})=\left(\begin{array}{ccc}2& 1& I\\ M_1& \mu & m\end{array}\right)\left(\begin{array}{ccc}I& 1& 2\\ m& \lambda & M_1^{}\end{array}\right),$$
(8)
where $`m=\frac{1}{2}(M_1+M_1^{}+\lambda \mu )`$. The intermediate (uncoupled) dispersion coefficients are
$$C_6^{J_aJ_b}=\frac{3}{\pi }_0^{\mathrm{}}𝑑\omega S_{J_a}(i\omega )S_{J_b}(i\omega ).$$
(9)
The reduced dynamic dipole polarizability $`S_I(i\omega )`$ of purely imaginary argument is defined as the sum over atomic states $`|\alpha _IIM_I`$ with total angular momentum $`I`$ and energy $`_{\alpha _I}`$
$$S_I(i\omega )=\underset{\alpha _I}{}\frac{(^{}_{\alpha _I})ns(3/2)_2d\alpha _II\alpha _IIdns(3/2)_2}{(^{}_{\alpha _I})^2+\omega ^2}.$$
(10)
Here $`\alpha _I`$ stands for all quantum numbers of the intermediate state, except for the total angular momentum $`I`$, and $`idj`$ are the reduced electric-dipole matrix elements, defined by the Wigner-Eckart theorem. Three sums with $`I=1,2,3`$ are allowed by electric-dipole selection rules.
We proceed now to construct the dipole polarizability functions $`S_I(i\omega )`$, Eq. (10), evaluate the uncoupled dispersion coefficients (9), and set and diagonalize the second-order effective Hamiltonian (5).
The functions $`S_I(i\omega )`$ satisfy several sum rules. In particular, the static tensor dipole polarizability $`\alpha _{zz}(M)`$ of the $`ns(3/2)_2`$ state may be expressed as
$$\alpha _{zz}(M)=2\underset{I}{}(1)^I\left(\begin{array}{ccc}2& 1& I\\ M& 0& M\end{array}\right)^2\times S_I(0).$$
(11)
The static tensor dipole polarizabilities $`\alpha _{zz}(M=1)`$ and $`\alpha _{zz}(M=2)`$ of metastable noble-gas atoms Ne through Xe have been measured by Molof et al. to within an error of 2%. In the present calculations the values of $`S_I(0)`$ are adjusted to reproduce these experimental values. In addition, as $`\omega \mathrm{}`$, the reduced polarizabilities satisfy the nonrelativistic Thomas-Reiche-Kuhn (TRK) sum rule
$$\underset{n}{}f_{kn}=\frac{2}{15}\underset{I}{}(1)^{I+1}S_I(i\mathrm{})=N,$$
(12)
$`N`$ being the number of electrons in the atom. Our constructed polarizabilities satisfy the sum rule.
It is instructive to consider the action of a one-particle operator on the reference particle-hole Slater determinant $`ns(3/2)_2`$ in the independent electron approximation. Such an operator can (i) annihilate the reference particle-hole pair; (ii) promote a valence electron from $`ns_{1/2}`$ state to another valence state $`mp_{1/2,3/2}`$, the state of the $`(n1)p_{3/2}`$ hole being unchanged; (iii) deexcite the $`(n1)p_{3/2}`$ hole into some other hole state, the valence state remaining the same, and (iv) create another particle-hole pair in addition to the reference pair. According to such a classification it is convenient to break the polarizability function, Eq. (10), into three contributions $`S_I^k`$ corresponding to a number $`k`$ of particle-hole pairs in the intermediate state $`|\alpha _II`$
$`S_I=S_I^0+S_I^1+S_I^2.`$
Since an electric-dipole transition from $`ns(3/2)_2`$ to the closed-core state $`{}_{}{}^{1}S_{0}^{}`$ is prohibited by the angular selection rules, $`S_I^00`$.
The sum $`S_I^1`$ is separated into contributions from the intermediate states in the lowest $`np`$ fine-structure multiplet $`\left(S_I^1\right)_{np}`$, and the rest of the sum $`\left(S_I^1\right)^{}`$
$`S_I^1=\left(S_I^1\right)_{np}+\left(S_I^1\right)^{}.`$
The first term is the dominant contribution. We calculate $`\left(S_I^1\right)_{np}`$ using experimental values of transition energies and decay rates, and adjusted branching ratios. The rest of the sum over valence states (including bound and continuum states) $`(S_I^1)^{}`$ is estimated in the Dirac-Hartree-Fock (DHF) approximation. The metastable state $`|vh;J`$ in lowest order is represented as a combination of the $`|v=ns_{1/2}`$ particle state and the hole state $`|h=(n1)p_{3/2}`$, coupled to the total angular momentum $`J=2`$, $`[((n1)p_{3/2})^1ns_{1/2}]_2`$. The lowest-order energy of such a state is $`^{}=\epsilon _v\epsilon _h`$, $`\epsilon _i`$ being the energy of the DHF orbital $`|i`$. The intermediate state is represented as $`|ma;I`$, a particle state $`m`$ coupled with hole state $`a`$ to the total momentum $`I`$. Explicitly,
$`(S_I^1)^{}(i\omega )=(1)^{J+I}[I][J]{\displaystyle \underset{ma}{}}^{}`$ $`(\delta _{ha}{\displaystyle \frac{(\epsilon _v\epsilon _m)vdm^2}{(\epsilon _v\epsilon _m)^2+\omega ^2}}\left\{\begin{array}{ccc}J& 1& I\\ j_m& j_h& j_v\end{array}\right\}^2`$
$`+\delta _{mv}{\displaystyle \frac{(\epsilon _a\epsilon _h)hda^2}{(\epsilon _a\epsilon _h)^2+\omega ^2}}\left\{\begin{array}{ccc}J& 1& I\\ j_a& j_v& j_h\end{array}\right\}^2),`$
where $`[K]2K+1`$, the summation is performed over the core orbitals $`a`$ and excited states $`m`$, excluding states of the lowest $`np`$ multiplet, and $`J=2`$. The first sum is associated with excitation of the valence electron (case (ii)) while the second sum with deexcitation of the hole state (case (iii)). To arrive at this result we disregarded the coupling between levels within the same fine-structure multiplet. For example, for Ne the $`[(2p_{3/2})^14p_{1/2}]_2`$, $`[(2p_{3/2})^14p_{3/2}]_2`$, and $`[(2p_{1/2})^14p_{3/2}]_2`$ states are summed over independently, even though the correct lowest-order wave-function is a linear combination of them. This approximation corresponds to a disregard of the small difference between the energies of the coupled and uncoupled states. Since the contribution $`(S_I^1)^{}`$ is relatively small, such an estimate suffices at the present level of accuracy. Numerical evaluation of $`(S_I^1)^{}`$ has been performed using a B-spline basis set in the $`V_{N1}`$ DHF potential, with the hole in the $`(n1)p_{3/2}`$ core orbital .
We separate the sum $`S_I^2`$ over the core excited states into two contributions
$`S_I^2=\left(S_I^2\right)_{\mathrm{core}}+\left(S_I^2\right)_{\mathrm{cntr}}.`$
The first term is associated with the dynamic polarizability $`\alpha _g(i\omega )`$ of the closed-shell ground state $`{}_{}{}^{1}S_{0}^{}`$ and the second term is a corrective counter term.
$$\left(S_I^2\right)_{\mathrm{core}}(i\omega )=\left(1\right)^{I+J+1}\frac{[I]}{2}\alpha _g\left(i\omega \right).$$
(15)
We use the semiempirical dynamic polarizabilities for the ground states of noble-gas atoms of Kumar and Meath . The estimated uncertainty of these core polarizabilities is less than 1%. The primary role of the core polarizability is to provide the correct limit Eq. (12) at $`\omega \mathrm{}`$. The relative importance of the core-excitation contribution increases for heavier systems; for example, in a similar calculation for Fr , core excitations contribute 23% of the $`C_6`$ dispersion coefficient. The high-frequency limit, Eq. (12), is accurately reproduced by the present total reduced dynamic polarizabilities $`S_I(i\omega )`$. We obtain for Ne 9.98, for Ar 17.95, for Kr 35.95, and for Xe 53.96 compared to the nonrelativistically exact values 10, 18, 36, and 54 respectively.
In the $`ns(3/2)_2=[((n1)p_{3/2})^1ns_{1/2}]_2`$ state core excitations to the occupied magnetic sub-states of the $`ns_{1/2}`$ particle state are not allowed by the Pauli exclusion principle, and neither are the core excitations from the empty magnetic sub-state of the hole $`(n1)p_{3/2}`$. To remove these transitions from the core polarizability contribution $`\left(S_I^2\right)_{\mathrm{core}}`$, we introduce a counter term $`\left(S_I^2\right)_{\mathrm{cntr}}`$. Explicitly in the independent-electron model
$`\left(S_I^2\right)_{\mathrm{cntr}}(i\omega )`$ $`=`$ $`[I][J](1)^{I+J}({\displaystyle \underset{a}{}}{\displaystyle \frac{(\epsilon _v\epsilon _a)adv^2}{(\epsilon _v\epsilon _a)^2+\omega ^2}}\left\{\begin{array}{ccc}1& J& I\\ j_h& j_a& j_v\end{array}\right\}^2`$
$`+`$ $`{\displaystyle \underset{m}{}}{\displaystyle \frac{(\epsilon _m\epsilon _h)hdm^2}{(\epsilon _m\epsilon _h)^2+\omega ^2}}\left\{\begin{array}{ccc}1& J& I\\ j_v& j_m& j_h\end{array}\right\}^2).`$
We estimate the small counter-term using the Dirac-Hartree-Fock approximation.
The largest contribution to the sums $`S_I`$ arises from the intermediate states in the lowest $`np`$ fine-structure multiplet. The determination of electric-dipole matrix elements involved in the sum $`\left(S_I^1\right)_{np}`$ requires a knowledge of both decay rates and branching ratios in the manifold. The relevant lifetimes have been measured to within an error less than 1% for Ne , Ar and Kr , and less than 3% for Xe . However, the branching ratios $`B`$ are not established to the same precision. The most accurate measurements of $`B`$ in Ne , have an error bar of approximately 4-5%, which would introduce an uncertainty of 4-5% in the static polarizabilities, and 8-10% inaccuracy in the values of $`C_6`$. To reduce the consequent errors, the experimental values of the static polarizability, accurate to 2%, were chosen as the reference data.
The branching ratios of transitions to the $`ns(3/2)_2`$ state have been adjusted as follows. The sum $`S_3(0)`$ includes only one intermediate state in the $`np`$ manifold, $`np(5/2)_3`$, and very small ab initio corrections. The $`np(5/2)_3`$ state has a single decay channel, so that the sum $`S_3(0)`$ is known with the experimental precision of the decay rate. The sums $`S_2(0)`$ and $`S_1(0)`$ can be deduced from the experimental values of the static tensor polarizability as
$`S_1(0)`$ $`=`$ $`{\displaystyle \frac{9}{14}}S_3(0)+5\alpha _{zz}(1){\displaystyle \frac{5}{4}}\alpha _{zz}(2)`$
$`S_2(0)`$ $`=`$ $`{\displaystyle \frac{5}{14}}S_3(0){\displaystyle \frac{15}{4}}\alpha _{zz}(2).`$
Removing small ab initio and semiempirical core-excitation contributions from these sums, the sums $`\left(S_1^1\right)_{np}(0)`$ and $`\left(S_2^1\right)_{np}(0)`$ are obtained. The branching ratios $`B`$ for four states involved in the $`J=1`$ sum and three states in the $`J=2`$ sum were multiplied by a uniform scaling factor. Branching ratios for Ne for the $`J=1`$ levels were multiplied by 1.0035, and for the $`J=2`$ level by 0.905 in order to reproduce the experimental values of the static tensor polarizabilities. We modified the recommended values of $`B`$ for Ar by multiplying the branching ratios of the $`J=1`$ states by 0.973 and of the $`J=2`$ states by 0.965; the values used in the calculations are listed in Table I. For Kr the velocity-gauge branching ratios, tabulated in Ref. from calculations by Aymar and Coulombe , were multiplied by 1.127 for the $`J=1`$ states and by 1.0016 for the $`J=2`$ states. For Xe, velocity-gauge values of $`B`$ calculated in Ref. were multiplied by 0.927 for the $`J=1`$ states and by 0.929 for the $`J=2`$ states. The adjusted data for Ar and Xe are listed in Table I. Doery et al. have compiled the input data for Ne and Kr, which have to be similarly modified.
We employ the constructed reduced polarizabilities $`S_I(i\omega )`$ to calculate the intermediate uncoupled dispersion coefficients $`C_6^{J_aJ_b}`$ by quadrature using Eq. (9). The coefficients are listed in Table II. They are to be used if the entire molecular Hamiltonian, including quadrupole-quadrupole and higher multipoles or perturbation-theory orders is to be diagonalized. Finally, the molecular terms are obtained by the diagonalization of $`H_{\mathrm{eff}}^{\left(2\right)}`$, given by Eq. (5). Neglect of the small corrections due to the quadrupole-quadrupole interaction results in parameterization of term energies in the form
$`U(R)=2^{}C_6/R^6.`$
The calculated dispersion coefficients $`C_6`$ for various molecular symmetries are listed in Table III. Since the region close to $`\omega =0`$ contributes the most to the values of the integral in Eq. (9), the uncertainty in the values of $`C_6`$ is approximately 4%, reflecting the 2% experimental error in the static dipole tensor polarizabilities . The values of the $`C_6`$ coefficients grow monotonically from Ne to Xe, due to the reduction in the energy separations between the metastable states and the $`np`$-manifold. For heavier systems the anisotropy in $`C_6`$, arising from relativistic effects becomes increasingly marked, from 6.5% in Ne to 16% in Xe.
Long-range dispersion coefficients for two interacting metastable Ne atoms were evaluated recently by Doery et al. . The $`C_6`$ coefficients were calculated from the diagonalization of the molecular dipole-dipole Hamiltonian in the model space containing the lowest $`np`$-manifold, so limiting the intermediate states to the lowest $`np`$-manifold states in the present formulation. Experimental values of decay rates and branching ratios were used to deduce the electric-dipole matrix elements. The precision of the calculated values of $`C_6`$ is about 8-10% due to the large uncertainty in the branching ratios. The values of $`C_6`$ from Ref. for different molecular symmetries vary between 1951 and 1956 a.u., exhibiting much less anisotropy than the present results which range between 1877 and 1999 a.u.. The difference can be traced to the anisotropy in the static dipole polarizabilities. Indeed, utilizing input data from Ref. we obtain $`\alpha _{zz}(M=1)=192`$ and $`\alpha _{zz}(M=2)=189`$ a.u. if we include only the $`np`$ manifold as in Ref. . While $`\alpha _{zz}(M=1)`$ agrees with the experimental value 192(4), the $`\alpha _{zz}(M=2)`$ is overestimated by three standard deviations compared to the experimental value 180(3) a.u..
The accuracy of the dispersion coefficients could be improved by applying relativistic all-order many-body methods to calculate transition amplitudes between $`nsnp`$ manifolds. Such ab initio calculations are intrinsically more challenging than for alkali-metal atoms; the accurate experimental lifetimes would provide an excellent gauge of accuracy.
Our values of $`C_6`$ coefficients will be useful in studies of cold collisions of metastable rare-gas atoms . For example, we can estimate the rate coefficient for Penning ionization by ignoring spin-polarization and assuming that every trajectory that surmounts the angular momentum barrier leads to ionization . The corresponding rate coefficient is given by
$`k=6.35\times 10^9{\displaystyle \frac{C_6^{1/3}T^{1/6}}{\mu ^{1/2}}}\mathrm{c}m^3s^1,`$
where $`\mu `$ is the reduced mass measured in units of the electron mass and $`T`$ is the temperature. Combined with short-range potentials a number of other properties could be determined. For example, scattering lengths of elastic collisions could be found, providing input for mean-field equations describing dilute quantum gases.
This work was supported by the U.S. Department of Energy, Division of Chemical Sciences, Office of Energy Research. Thanks are due to M. R. Doery, S. Kotochigova, and J. F. Babb for useful discussions. The authors are grateful to W. R. Johnson for providing the B-spline routine for the $`V_{N1}`$ Dirac-Hartree-Fock potential.
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# 1. Introduction
## 1. Introduction
The observed smallness of the cosmological constant $`\lambda `$ in Einstein’s equations poses a fine-tuning problem already in classical field theory coupled to gravity. E.g., when the Higgs field rolls down its potential, the energy density of the vacuum and thus the effective value of $`\lambda `$ changes by a large amount. Other expected contributions to $`\lambda `$ that are suspiciously absent include those from condensates in QCD.
But perhaps the most mysterious aspect of the problem is that $`\lambda `$ does not seem to receive contributions from the quantum mechanical ground state energies
$$\rho _k=\frac{\mathrm{}\omega }{2},\omega ^2=k^2+m^2$$
of the oscillators with momentum $`k`$ of the massless and light fields of the Standard Model. Summing these contributions up to some large-momentum cutoff $`\mathrm{\Lambda }m`$, one finds in the case of a single–component bosonic field :
$`\lambda =8\pi G{\displaystyle _0^\mathrm{\Lambda }}{\displaystyle \frac{k^2dk}{2\pi ^2}}\rho _k{\displaystyle \frac{\mathrm{\Lambda }^4}{2\pi }}l_P^2.`$ (1.1)
Here, $`G`$ is the Newton constant and $`l_P=\sqrt{\mathrm{}G}1.710^{35}m`$ is the Planck length (setting $`c=1`$). Experimentally, it presently seems that
$`\lambda (210^{33}eV)^2(10^{26}m)^2`$ (1.2)
(setting $`\mathrm{}=G=1`$). $`10^{26}m`$ is the order of magnitude of the curvature radius of the universe, which is roughly the inverse Hubble constant. But even if one considered only the contributions of the two helicity states of the massless photon to the cosmological constant (1.1), then in order to explain such a small value of $`\lambda `$ one would need a momentum–cutoff as small as
$$\mathrm{\Lambda }\frac{1}{100}eV\frac{1}{20\mu m}.$$
So the minimum wavelength would have to be as large as $`20\mu m`$. Already the wavelength of visible light is much smaller than 20 $`\mu m`$. In this sense the observed smallness of the Hubble expansion parameter seems inconsistent even with what we can see with our bare eyes.
Supersymmetry could explain a zero cosmological constant, because the vacuum energies of the superpartners cancel each other. Supersymmetry looks so much like the missing piece in the puzzle that it has been questioned whether supersymmetry is really broken . But if it isn’t broken, it is of course hard to explain why we do not see superpartners of the Standard Model fields .
The recently revived suggestion that we live on a 4–dimensional brane that is embedded in a higher-dimensional bulk opens up a new perspective and a way out (under an assumption stated in section 2). It will be proposed below that supersymmetry is indeed unbroken up to micrometer scales – but only in the bulk supergravity theory. By a simple argument, supersymmetry breaking in the bulk at micrometer scale is derived from supersymmetry breaking at the $`TeV`$ scale on the brane, which carries the Standard Model fields.
Due to a mechanism first proposed in and re–invented in , the vacuum energy of the brane fields is shown not to contribute to the $`4d`$ cosmological constant. Rather, it is absorbed in the curvature transverse to the brane. Only the vacuum energy of the bulk supergravity fields is argued to contribute to the $`4d`$ cosmological constant.
This vacuum energy is estimated. The result is a relation between the Planck mass $`m_{Planck}`$, the scale $`m_{BraneSusy}`$ of supersymmetry breaking on the brane, the scale $`m_{BulkSusy}`$ of supersymmetry breaking in the bulk sector and the Hubble expansion rate $`H_0`$: roughly,
$`\mathrm{log}{\displaystyle \frac{m_{Planck}}{m_{BraneSusy}}}{\displaystyle \frac{1}{2}}\mathrm{log}{\displaystyle \frac{m_{Planck}}{m_{BulkSusy}}}{\displaystyle \frac{1}{4}}\mathrm{log}({\displaystyle \frac{m_{Planck}}{H_0}})`$ (1.3)
(a more detailed relation is given in the text). Based on the known values of $`m_{Planck}`$ and $`H_0`$, this relation predicts gravitinos or other superpartners of the supergravity multiplet with masses of order $`10^3eV`$ (which is inside experimental bounds ) and a supersymmetry breaking scale on the brane of $`26`$ $`TeV`$. Conversely, based on the assumption that supersymmetry is restored in the Standard Model at energies not too much above the weak scale, the relation explains the observed small value of the cosmological constant.
The setup is introduced in section 2. In section 3 it is argued that the $`4d`$ cosmological constant is zero as long as the bulk supergravity theory is treated classically. In section 4 it is shown that the quantum mechanical ground state energy of the supergravity sector produces a cosmological constant that is within a few orders of magnitude of its observed value. Precise matching yields the predictions for supersymmetry breaking in the Standard Model and in the bulk sector, as explained in section 5. Section 6 contains conclusions.
## 2. The setup
We consider a 3–brane soliton that is embedded in a ($`4+n`$)–dimensional bulk spacetime (figure 1). We assume that the $`n`$ extra dimensions are compactified on some manifold $``$. The Standard Model fields are assumed to live only on the brane, while gravity lives in the bulk. Let Vol($``$) be the volume of the compactification manifold, and let Vol($``$) be the volume of the ball $``$ inside $``$ that intersects with the brane. Since we are going to consider non-supersymmetric branes inside supersymmetric bulks (as, e.g., in ), we will identify the size (i.e. the thickness) of the brane with $`l_{BraneSusy}m_{BraneSusy}^1`$, the scale of supersymmetry breaking in the Standard Model. So roughly, Vol($``$) $`(l_{BraneSusy})^n`$.
Similarly as in , because the Einstein action is integrated over Vol($``$) while the Standard Model action is integrated only over Vol($``$), the 4–dimensional Planck length $`l_{Planck}`$ is related to $`l_{BraneSusy}`$ by
$`({\displaystyle \frac{l_{Planck}}{l_{BraneSusy}}})^2=({\displaystyle \frac{m_{BraneSusy}}{m_{Planck}}})^2{\displaystyle \frac{\text{Vol(}\text{)}}{\text{Vol(}\text{)}}}`$ (2.1)
(assuming a $`(4+n)`$–dimensional Newton constant of order one).
Let us first consider the supersymmetric version of the story. So we assume that we have a supersymmetric brane inside a supersymmetric compactification manifold. In string theory, this is achieved by considering a compactification on a Calabi-Yau 3–manifold that involves branes parallel to the 4–dimensional space–time, as in . At distances much larger than the size of the compactification manifold, only a four–dimensional supersymmetric effective theory of Standard Model fields plus their superpartners coupled to $`4d`$ supergravity is seen.
For concreteness, we may assume a metric in the vicinity of the brane of the form
$$ds^2=dr^2+f(r)\widehat{g}_{\mu \nu }dx^\mu dx^\nu +g(r)d\mathrm{\Omega }_5^2$$
where $`r`$ denotes the distance from the brane, $`x^\mu `$ are the space-time coordinates parallel to the brane, $`\widehat{g}_{\mu \nu }`$ is the $`4d`$ metric parallel to the brane, and $`f(r),g(r)`$ are some functions. Supersymmetry of the effective $`4d`$ theory implies that the $`4d`$ metric $`\widehat{g}_{\mu \nu }`$ is Ricci–flat (we are assuming that there are no 4–form gauge field strengths or expectation values of other supergravity fields), i.e. the effective $`4d`$ cosmological constant is zero.
Now suppose that we cut out a region of radius $`l_{BraneSusy}`$ around the brane. The basic assumption under which the arguments in the next section apply is that, at the level of classical supergravity, we can consistently do the following: we can replace the supersymmetric brane soliton solution by a stable non–supersymmetric one (perhaps of the type of ), such that the bulk fields smoothly connect to a solution at $`rl_{BraneSusy}`$ that does not break supersymmetry on the $`4d`$ slices parallel to the brane.
In other words, we assume that there are consistent string compactifications that involve space-time filling stable non-BPS branes, such that $`4`$-dimensional supersymmetry is unbroken away from the brane at least in the classical supergravity approximation. The construction of explicit examples must be left for future work.
In the case of one extra dimension, examples of supergravity solutions that smoothly interpolate between a supersymmetric and a non–supersymmetric region are the kink solutions of $`5d`$ gauged supergravity discussed in .
Supersymmetry is now broken not only in the bulk theory in the vicinity of the brane. It is also broken in the world–brane theory that contains the Standard Model fields and lives in the non–supersymmetric gravitational background. This will result in a brane vacuum energy of the order of $`(m_{BraneSusy})^4`$.
## 3. Classical Supergravity Approximation
Let us first explain why the vacuum energy on the brane does not curve the $`4d`$ metric $`\widehat{g}_{\mu \nu }`$ parallel to the brane (i.e., why it does not create a $`4d`$ cosmological constant) as long as the bulk supergravity theory is treated classically (see ). Although the bulk theory is treated classically, the world–brane theory containing the Standard Model fields is assumed to be treated fully quantum mechanically. Corrections from loops of the bulk fields are very interesting and will be discussed in the next section.
The bulk has been separated into two regions: the non–supersymmetric neighborhood of the brane $`M^4\times `$, where $`M^4`$ is the Minkowski space parallel to the brane; and the supersymmetric region, i.e. the rest of the bulk $`M^4\times ()`$. The classical supergravity equations of motion can be solved separately for each region, and can then be matched at their interface at $`r=l_{BraneSusy}`$.
In the bulk region, the $`4d`$ metric $`\widehat{g}_{\mu \nu }`$ parallel to the brane must still be Ricci–flat because of supersymmetry on the $`4d`$ slices parallel to the brane.
As for the brane region, there may be a singularity or horizon near the center. Let us therefore restrict the discussion to the region $`ϵrl_{BraneSusy}`$, where $`ϵ`$ is a cutoff that hides the singularity or horizon. The issue of boundary conditions at $`r=ϵ`$ will be commented on below.
In this non–supersymmetric brane region, the brane is a source of vacuum energy $`\rho `$ of order $`(m_{BraneSusy})^4`$ that arises from the world–brane fields. Let us assume some distribution $`\rho (r)`$ around $`r=0`$ with width of the order $`l_{BraneSusy}`$. $`\rho (r)`$ enters the Einstein equations like an $`r`$–dependent cosmological constant:
$$R_{mn}\frac{1}{2}g_{mn}R=\widehat{\lambda }(r)g_{mn}$$
with
$$\widehat{\lambda }(r)=8\pi G\rho (r)\lambda _{flux}(r).$$
Here we have included another $`r`$–dependent contribution $`\lambda _{flux}(r)`$ that arises when the brane is a source of electric or magnetic flux.
For simplicity, we focus on the example of a single extra dimension, assume a constant dilaton and neglect the other supergravity fields; the generalization is straightforward. We make the metric ansatz
$$ds^2=dr^2+e^{2\alpha (r)}\widehat{g}_{\mu \nu }dx^\mu dx^\nu .$$
In this ansatz, the $`4d`$ metric $`\widehat{g}`$ is taken to be $`r`$–independent. The 5–dimensional Ricci tensor can be written (cmp. with ):
$$R_{\mu \nu }^{(5)}=\widehat{R}_{\mu \nu }^{(4)}\widehat{g}_{\mu \nu }e^{2\alpha (r)}(\ddot{\alpha }+4\dot{\alpha }^2)$$
(a “dot” means $`\frac{d}{dr}`$). Plugging this into the Einstein equation for the 4-dimensional components ($`\mu ,\nu `$) and using the equation for the $`(r,r)`$ component to eliminate $`\ddot{\alpha }`$,
$$4(\ddot{\alpha }+\dot{\alpha }^2)=\frac{2}{3}\widehat{\lambda },$$
we obtain:
$$\widehat{R}_{\mu \nu }^{(4)}=k^2\widehat{g}_{\mu \nu }\text{where}k^2=e^{2\alpha }(\frac{1}{2}\widehat{\lambda }+3\dot{\alpha }^2)$$
is an integration constant that is by definition the $`4d`$ cosmological constant ($`k^2,\widehat{\lambda }`$ may be negative). So the equations for $`\alpha `$ have a one–parameter family of solutions, labelled by the constant $`4d`$ curvature $`k^2`$. However, matching at $`r=l_{BraneSusy}`$ to the solution in the supersymmetric region requires that we pick the solution that is Ricci–flat in $`4d`$, i.e. $`k=0`$. For this solution, the vacuum energy on the brane is completely absorbed by the warp factor
$$\dot{\alpha }^2=\frac{1}{6}\widehat{\lambda }(r),$$
and therefore does not curve the $`4d`$ metric parallel to the brane. So the vacuum energy does not lead to a $`4d`$ cosmological constant. For $`n`$ extra dimensions, the discussion is similar.
This is the mechanism of Rubakov and Shaposhnikov , recently rediscovered in . We have supplemented it by a matching condition at $`r=l_{BraneSusy}`$ that picks out the solution with vanishing $`4d`$ cosmological constant without fine–tuning. This is a generalization of the suggestion in of “supersymmetry on the Planck brane” in the context of the Randall-Sundrum model. Higher–order corrections will make the differential equations for $`\alpha `$ more complicated, and there may be regions in parameter space where no solutions exist ; let us assume that conditions are favorable and solutions exist.
We have not discussed boundary conditions for $`\alpha `$ at the cutoff $`r=ϵ`$, where the supergravity approximation presumably breaks down. However, whatever boundary conditions must be imposed – the assumption that they can be satisfied is part of the assumption that we have already made in the previous section: that there are consistent string compactifications that involve stable non-BPS branes and leave 4-dimensional supersymmetry unbroken away from the brane at the classical level. Again, it remains to construct explicit examples.
## 4. Supergravity at One Loop
Let us now go beyond the classical supergravity approximation. This is the main new step taken in this paper and it will lead to our numerical results.
As mentioned in the introduction, the ground state energies
$$\frac{\mathrm{}\omega }{2}\text{with}\omega ^2=k^2+m^2$$
of modes of light fields with momentum $`k`$ should give a quantum mechanical contribution to the cosmological constant. We have already demonstrated that the ground state energy of the Standard Model fields does not contribute to the $`4d`$ cosmological constant, so it only remains to compute the vacuum energy produced by the bulk supergravity fields: the gravitino, the dilaton, antisymmetric tensor fields etc.
As long as supersymmetry is unbroken in the bulk, these vacuum energy contributions cancel. Now, breaking supersymmetry in the region of the bulk near the brane also breaks supersymmetry in the effective $`4d`$ theory, obtained by integrating over the compactification manifold $``$. But because of the small overlap of the wave functions of the supergravity fields with the brane, the mass scale $`m_{BulkSusy}`$ of supersymmetry breaking in the bulk sector of the $`4d`$ effective theory will be suppressed with respect to the scale of supersymmetry breaking on the brane by the same volume factor that we already found in (2.1),
$`({\displaystyle \frac{m_{BulkSusy}}{m_{BraneSusy}}})^2{\displaystyle \frac{\text{Vol(}\text{)}}{\text{Vol(}\text{)}}}({\displaystyle \frac{m_{BraneSusy}}{m_{Planck}}})^2.`$ (4.1)
One way of seeing this is to consider a scalar field $`\mathrm{\Phi }`$ in the supergravity multiplet and assume that it has a large mass of order $`m_{BraneSusy}`$ inside the region where supersymmetry is broken: we take its $`(4+n)`$–dimensional Lagrangean to be of the form
$$_m\mathrm{\Phi }^m\mathrm{\Phi }+\theta (r_{BraneSusy}r)m_{BraneSusy}^2\mathrm{\Phi }^2.$$
$`\theta `$ is the step function: $`\theta (x)=0`$ for $`x<0`$ and $`\theta (x)=1`$ for $`x0`$. Integrating this Lagrangean over the compactification manifold, the kinetic term acquires a prefactor Vol($``$) while the mass term only acquires a prefactor Vol($``$). After normalizing $`\mathrm{\Phi }`$ to have a standard kinetic term, its mass is
$$m_\mathrm{\Phi }^2m_{BraneSusy}^2\frac{\text{Vol(}\text{)}}{\text{Vol(}\text{)}}.$$
This implies relation (4.1).<sup>1</sup><sup>1</sup>1The relation $`m_\mathrm{\Phi }\frac{m_{BraneSusy}^2}{m_{Planck}}`$ could also have been derived without reference to branes; in this case the suppression factor is simply due to the smallness of Newton’s constant.
So the hierarchy between the scales of supersymmetry breaking in the bulk supergravity sector and supersymmetry breaking in the Standard Model that lives on the brane is the same as the hierarchy between the scale of supersymmetry breaking on the brane and the Planck scale.
Already in the introduction we have discussed the relation (1.1) between the momentum cutoff $`\mathrm{\Lambda }`$ in the sum over vacuum energies and the value of the cosmological constant $`\lambda `$. In the case of $`N`$ massless bosonic propagating degrees of freedom, the relation changes to
$`\lambda N{\displaystyle \frac{\mathrm{\Lambda }^4}{2\pi }}l_P^2.`$ (4.2)
$`\lambda `$ is related to the Hubble expansion rate $`H_0`$ of the universe by
$$\lambda =3\mathrm{\Omega }_\mathrm{\Lambda }H_0^2\text{with}\mathrm{\Omega }_\mathrm{\Lambda }\frac{2}{3}$$
being the value suggested by observation . In a first estimate we may identify the cutoff $`\mathrm{\Lambda }`$ in (4.2) with the scale $`m_{BulkSusy}`$ of supersymmetry breaking in the bulk.<sup>2</sup><sup>2</sup>2The previous version of this paper used this first estimate (i.e. it set $`|Q|1`$ below) to suggest supersymmetry breaking scales of $`410`$ $`TeV`$ on the brane and $`10^310^2`$ $`eV`$ in the bulk (compare with section 5). Then equation (4.2) implies (converting $`l_{Planck}m_{Planck}^1`$):
$`({\displaystyle \frac{m_{BulkSusy}}{m_{Planck}}})^2{\displaystyle \frac{6\mathrm{\Omega }_\mathrm{\Lambda }\pi }{N}}({\displaystyle \frac{H_0}{m_{BulkSusy}}})^2.`$ (4.3)
A more precise calculation involves the various masses of order $`m_{BulkSusy}`$ of the supergravity fields. Then $`k`$ in (1.1) is integrated not only up to $`m_{BulkSusy}`$, but up to $`k_{max}m_{BraneSusy}`$, which is the fundamental scale in our setup. Formula (1.1) generalizes to (see e.g. for a discussion):
$`\lambda {\displaystyle \frac{l_P^2}{2\pi }}\times {\displaystyle \underset{i}{}}(1)^{F_i}\{k_{max}^4{\displaystyle \frac{1}{2}}m_i^4\mathrm{ln}{\displaystyle \frac{k_{max}}{m_i}}+\mathrm{}\},`$ (4.4)
where $`i`$ counts the propagating degrees of freedom in the supergravity multiplet, $`m_i`$ are their masses of order $`m_{BulkSusy}`$ after supersymmetry breaking, and $`(1)^{F_i}`$ is $`+1`$ for bosonic and $`1`$ for fermionic degrees of freedom. The $`k_{max}^4`$ terms cancel since there is an equal number of bosons and fermions.<sup>3</sup><sup>3</sup>3A possible term of the form $`k_{max}^2_i(1)^{F_i}m_i^2`$ that may appear in a more general calculation should also vanish, since supersymmetry is broken spontaneously by the non-supersymmetric soliton inside the supersymmetric bulk. The conclusion is then that $`N`$ in (4.3) is replaced by
$`N=Q{\displaystyle \underset{i}{}}(1)^{F_i}({\displaystyle \frac{m_i}{m_{BulkSusy}}})^4`$ (4.5)
with
$$Q\frac{1}{2}\mathrm{ln}\frac{m_{BraneSusy}}{m_{BulkSusy}}.$$
Together, with (4.1), (4.3) yields the relation claimed in the introduction (where we have set $`\mathrm{\Omega }_\mathrm{\Lambda }=\frac{2}{3}`$ and roughly approximated $`(\frac{|N|}{4\pi })^{\frac{1}{2}}`$ by 1):
$`\mathrm{log}({\displaystyle \frac{m_{Planck}}{m_{BraneSusy}}})={\displaystyle \frac{1}{2}}\mathrm{log}({\displaystyle \frac{m_{Planck}}{m_{BulkSusy}}})={\displaystyle \frac{1}{4}}\mathrm{log}(\sqrt{{\displaystyle \frac{|N|}{6\mathrm{\Omega }_\mathrm{\Lambda }\pi }}}{\displaystyle \frac{m_{Planck}}{H_0}})L.`$ (4.6)
## 5. The Numbers
Let us now plug in the numbers. We use
$`m_{Planck}`$ $``$ $`10^{19}GeV`$ (5.1)
$`\sqrt{\lambda }=\sqrt{3\mathrm{\Omega }_\mathrm{\Lambda }}H_0`$ $``$ $`210^{33}eV`$ (5.2)
What is $`N`$? Type IIB supergravity multiplets, e.g., have 128 bosonic and 128 fermionic degrees of freedom. Without going into details, it seems safe to assume that $`|\frac{N}{Q}|`$ in (4.5) is somewhere between 1 and 128. With $`|Q|20`$, this gives
$$\sqrt{\frac{|N|}{2\pi }}2\text{to}20$$
Since there may be other hidden factors of $`\pi ,\frac{1}{2},`$ etc. that were missed by our crude analysis, the actual errors in the relation (4.6) may even be somewhat (but not much) larger. Being optimistic about them, we infer that
$$L15.4\pm 0.2$$
This yields the predictions
$`m_{BraneSusy}`$ $``$ $`2TeV6TeV`$ (5.3)
$`m_{BulkSusy}`$ $``$ $`{\displaystyle \frac{1}{2}}10^3eV{\displaystyle \frac{1}{2}}10^2eV.`$ (5.4)
The example with the intermediate value of $`L=15.4`$ is plotted in figure 2.
So we expect a mass of order $`10^3eV`$ for the gravitino or at least for some members of the supergravity multiplet. This corresponds to a Compton wavelegth of the order of a fraction of a millimeter. Similarly as in the case of millimeter–size extra dimensions , the presence in the bulk of gravitinos or dilatons in the micrometer range is not ruled out by experiment: while the brane physics (the Standard Model) has been probed down to the weak scale, the bulk physics (gravity) has only been probed down to centimeter scales. The lower experimental bound on the gravitino mass appears to be only $`10^5eV`$ . The effects of these new fields might show up in short–distance measurements of gravity in the $`\mu m`$ range in the near future .
Remarkably, the predicted scale of supersymmetry breaking in the Standard Model is roughly where it is expected to be, in order to insure that the running coupling constants meet in supersymmetric Grand Unification. This is very nontrivial; a priori it could have come out many orders of magnitude off the mark. Reversing the logic, if we assume a probable scale of supersymmetry breaking between 1 and 100 TeV, then we can predict the value of the cosmological constant within a few orders of magnitude of the value that seems to have been measured!
Let us finally note that our derivation and results apply just as well to the case of a single extra dimension as in the Horava–Witten model or in the Randall–Sundrum model .
## Conclusion
It seems that the proposal that we live on a non–supersymmetric brane that is embedded in a supersymmetric higher–dimensional string compactification can explain the observed small value of the cosmological constant, provided that the scale of supersymmetry breaking in the Standard Model is roughly 2–6 TeV. It remains to construct explicit examples of such compactifications and to show that they are consistent.
Our explanation for the small cosmological constant can be tested by searching for signs of a gravitino, a dilaton or other supergravity fields with masses of order $`10^3`$ $`eV`$. We can thus look forward to a number of surprises in future tests of Einstein gravity in the micrometer range.
Acknowledgements:
I thank W. Lerche and P. Mayr for discussions. I also thank I. Antoniadis, R. Barbieri, S. Dimopoulos, S. Kachru, G. Veneziano, H. Verlinde and E. Witten for comments on the previous version of this paper. This work is supported in part by a Heisenberg fellowship of the DFG.
Note: After this paper appeared, I was notified by the authors of that they mention on page 5 the possibility of using infinite volume extra dimensions to suppress the breaking of $`5d`$ bulk supersymmetry and the $`5d`$ bulk cosmological constant. (Note however that we discuss $`4d`$ supersymmetry and the $`4d`$ cosmological constant on slices parallel to the brane.)
I also became aware of , where non-supersymmetric branes inside a supersymmetric bulk are constructed.
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# Stripes in Doped Antiferromagnets: Single-Particle Spectral Weight
## Abstract
Recent photoemission (ARPES) experiments on cuprate superconductors provide important guidelines for a theory of electronic excitations in the stripe phase. Using a cluster perturbation theory, where short-distance effects are accounted for by exact cluster diagonalization and long-distance effects by perturbation (in the hopping), we calculate the single-particle Green’s function for a striped $`t`$-$`J`$ model. The data obtained quantitatively reproduce salient (ARPES-) features and may serve to rule out ”bond-centered” in favor of ”site-centered” stripes.
There is by now substantial experimental evidence for a tendency of doped holes in high-temperature superconductors (HTSC) to form stripes, leaving behind locally antiferromagnetic (AF) domains. The stripes can be static like in La<sub>1.48</sub>Nd<sub>0.4</sub>Sr<sub>0.12</sub>CuO<sub>4</sub> (Nd-LSCO) or dynamic like in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-x</sub> . At present, there is an intensive discussion and controversy, whether stripes are directly connected with and beneficial for the microscopic mechanism in HTSC . What is clear, however, is that the apparent presence of static or dynamic stripes crucially influences low-energy excitations and thus the foundations for such a microscopic theory. Evidence for this has recently been accumulated by angle-resolved photoemission spectroscopy (ARPES) both on the static stripes in the Nd-LSCO system and on dynamic domain walls in the LSCO compound . The electronic structure revealed by ARPES contains characteristic features consistent with other cuprates, such as the flat band at low energy near the Brillouin zone face ($`𝐤=(\pi ,0)`$). In Nd-LSCO, the frequency-integrated spectral weight is confined inside 1D-segments in $`𝐤`$-space, deviating strongly from the more rounded Fermi surface expected from band calculations.
In this Letter, we present a numerical study of the single-electron excitations in a striped phase via cluster perturbation theory (CPT) and a detailed comparison with recent ARPES results. The basic idea of our application of the CPT is indicated in Fig.(1): it is based on dividing the 2D plane into alternating clusters of metallic stripes and AF domains. The individial clusters are modeled by the microscopic $`t`$-$`J`$ hamiltonian and solved exactly via exact diagonalization (ED). Then the inter-cluster hopping linking the alternating metallic and AF domains is incorporated perturbatively via CPT on the basis of the exact cluster Green’s functions thus yielding the spectral function of the infinite 2D plane in a striped phase. Using this CPT approach, the important short-distance interaction effects within the stripes are taken into account exactly while longer-ranged hopping effects are treated perturbatively. A study of the properties of experimentally observed stripe phases solely by ED is precluded by the prohibitively large unit cells. The manageable clusters for ED are simply too small to accommodate even a single such unit cell.
Our main results are: (i) close to $`𝐤=(\pi ,0)`$ we see, like in experiments, a two-component electronic feature (see Fig.(2)): a sharp low-energy feature close to $`E_F`$ and a more broad feature at higher binding energies. Both features can be explained by the mixing of metallic and antiferromagnetic bands at this $`𝐤`$-point. (ii) the excitation near $`(\pi /2,\pi /2)`$ is at higher binding energies than the low-energy excitation at $`(\pi ,0)`$ and of reduced weight. (iii) the integrated spectral weight of the cluster-stripe calculation resembles the quasi-one-dimensional segments in momentum space (see Fig.(3a)) as seen in the Nd-LSCO experiment. Also in agreement with the Nd-LSCO experiment our calculation finds the low-energy excitations near $`(\pm \pi ,0)`$ and $`(0,\pm \pi )`$ (Fig.(3b)). Interestingly, this agreement with experiment occurs only for so called ”site-centered” metallic stripes (as shown in Fig.(1)) and not for ”bond-centered” metallic stripes. This seems important since it has been argued , that, for bond-centered stripes, superconductivity is expected to survive stripe ordering. In the DMRG calculations by White and Scalapino as well as in dynamical mean-field (DMFT) studies by Fleck et al. , bond-centered and site-centered domains are very close in energy (ground-state). However, our technique allows to distinguish them dynamically.
The computational technique for our calculation of the single-particle spectral weight $`A(𝐤,\omega )`$ and the Green’s function is illustrated in Fig.(1). We solve the AF cluster (here, $`N\times 3`$) and the metallic cluster (here $`N\times 1`$ with a hole filling $`n_h=0.5)`$ by ED and combine the individual clusters to an infinite lattice via CPT as in the homogeneous case. Where it was technically feasible, we extended the unit cell by a factor of two and diagonalized two AF 3-leg clusters with a staggered magnetic field pointing in opposite directions, resulting in a $`\pi `$-phase shifted Néel order in the final configuration. This site-centered ”$`3+1`$” configuration with $`\pi `$-phase shifted Néel order of stripes was first suggested by Tranquada et al. . Bond-centered stripes, on the other hand, are modeled by 2-leg ladders with alternating filling (half-filled, $`n_h=0`$ and doped, $`n_h=1/4`$). In the following we will refer to this bond-centered configuration as ”$`2+2`$”.
Holes can propagate out of the metallic stripes into the AF insulating domains via the inter-cluster hoppings. In the Hubbard model, these hoppings correspond to one-body operators and can be treated within a systematic strong-coupling perturbation expansion. For homogeneous systems such an expansion was constructed in ref.. We take the lowest order contribution:
$$G_{\mathrm{}}(𝐏,z)=\frac{G_{cluster}(z)}{1\epsilon (𝐏)G_{cluster}(z)},$$
(1)
which is of RPA form. Here, $`𝐏`$ is a superlattice wave vector and $`G_{\mathrm{}}`$ is the Green’s function of the ”$`\mathrm{}`$-size” 2D system, however, still in a hybrid representation: real space within a cluster and Fourier-space between the clusters. This is related to the fact that $`G_{\mathrm{}}(𝐏,z)`$ is now an $`M\times M`$ matrix in the space of site indices (in the inhomogeneous stripe configuration of Fig.(1) $`M=N\times 3+N\times 1=4N`$). Likewise, $`\epsilon (𝐏)`$ and $`G_{cluster}`$ are $`M\times M`$ matrices in real space with $`\epsilon (𝐏)`$ standing for the perturbation, which includes hoppings out of the clusters. A true Fourier representation of $`G_{\mathrm{}}`$ in terms of the original reciprocal lattice then yields the lowest order CPT approximation. As discussed in ref. for the homogeneous case, the approximation in eq.(1) is exact for vanishing interaction. When the interactions are turned on, eq.(1) is no longer exact (apart from the local limit, i.e. $`t=t^{}=0`$), but strong interactions are known to be important mainly for short-range correlations. These correlations are incorporated with good accuracy in modest-size clusters and are treated here by ED within the cluster. It has been shown by Sénéchal et al. that the CPT reproduces the spectral weight of the 1D and 2D Hubbard models in quantitative agreement with exact results.
To allow for larger cluster sizes $`N`$, we diagonalize the $`t`$-$`J`$ model and take its spectral (Green’s) function as our local $`G_{cluster}`$ in eq.(1) as an approximation to the Hubbard model’s (cluster-) Green’s function. A comparison of the Hubbard and $`t`$-$`J`$ model’s spectral function on small clusters shows that the strong low energy peaks have similar dispersion and weight in both models, the main difference being a transfer of incoherent high energy spectral weight from momenta near $`(\pi ,\pi )`$ to $`(0,0)`$. The $`t`$-$`J`$ hamiltonian is defined as
$$H=\underset{ij,\sigma }{}t_{i,j}\widehat{c}_{i,\sigma }^{}\widehat{c}_{j,\sigma }+J\underset{<i,j>}{}(𝐒_i𝐒_j\frac{n_in_j}{4}).$$
(2)
The hopping matrix element $`t_{i,j}`$ is nonzero only for nearest ($`t`$) and next-nearest neighbors ($`t^{}`$). The second sum counting the Heisenberg interaction $`J`$ runs over all nearest neighbor pairs. No double occupancy is allowed. We have chosen commonly accepted values for the ratio $`t^{}/t=0.2`$, $`J/t=0.4`$, where $`t0.5eV`$. In this Letter we present calculations for systems with $`N=8`$ based on diagonalizations of a $`N\times 3=24`$-site half-filled 3-leg ladder and a quarter-filled $`8`$-site chain. Results for smaller $`N=6`$ do not differ much from $`N=8`$ results. $`N=6`$, however, is somewhat pathological, since it has an odd number of electrons in the quarter-filled chain.
We proceed to the discussion of the spectra: Fig.(2a) shows the experimental ARPES results for LSCO at the superconductor-insulator transition (doping $`x=0.05`$) . To enhance the structure in the obtained spectra, the authors of ref. plotted the second derivative of the ARPES spectrum, so areas with high second derivative are marked white and areas with low curvature (i.e. flat intensity) are black. This result is compared with the theoretical CPT calculation for different stripe configurations with overall doping of $`x=1/8`$. Figs.(2b,c) are for the ”$`3+1`$” site-centered configuration. Fig.(2b) shows the result for a ”$`3+1`$” configuration with alternating (i.e. $`\pi `$-phase shifted) Néel order between the 3-leg ladders (induced by a staggered magnetic field $`B=0.1t`$) without next-nearest neighbor hopping, Fig.(2c) shows the result for the ”$`3+1`$” stripe configuration with next nearest-neighbor hopping $`t^{}=0.2t`$, however without Néel order (due to the reduced symmetry, a $`t^{}`$ diagonalization with staggered field is not technically feasible). Fig.(2d) shows the result for a bond-centered ”$`2+2`$” stripe configuration. We observe that the spectra for the site-centered ”$`3+1`$” configuration (Fig.(2b,c)) are in surprisingly good agreement with experiment, and, that the $`t^{}=0`$ calculation (Fig.(2b)) results in much more coherent bands due to the enforced Néel order in the 3-leg ladders. Similar to recent DMFT calculations , the sharp excitation near the Fermi surface around $`(\pi ,0)`$, that has been interpreted by Ino et al. as the quasiparticle peak in the SC state, is visible as well as a dispersive band at higher binding energies which (at least away from $`(\pi ,0)`$; see discussion below) can be interpreted as remnants of the insulating valence band resulting from the AF domains. Especially in the Néel ordered configuration, we observe a very coherent and pronounced band. This clear dispersion is also visible in the $`(\pi ,\pi )`$ direction near $`(\pi /2,\pi /2)`$. We note that, in agreement with the experimental result (Fig.(2a)), the excitation at $`(\pi ,0)`$ is at significantly lower binding energy than the excitation at $`(\pi /2,\pi /2)`$. Neither the 2D $`t`$-$`J`$ model nor a 2D $`t`$-$`t^{}`$-$`t^{\prime \prime }`$-$`J`$ model, with its parameters fitted to the insulating state, can reproduce this result . This is a crucial effect of the stripe assumption: With the stripes oriented along the $`y`$-direction, the metallic band is dispersionless in $`x`$-direction. Therefore, at $`𝐤=(\pi ,0)`$, the minimum of the metallic spinon band (located at $`(k_x,0)`$ for any $`k_x`$) hybridizes with the top of the insulating valence band resulting in a two-peak structure with one peak pushed to higher and the other pushed to lower binding energies (see below). At $`𝐤=(\pi /2,\pi /2)`$, on the other hand, the metallic band has crossed the Fermi surface (its $`k_F`$ being $`\pi /4`$) and no mixing takes place. Finally, the ”$`2+2`$” bond-centered stripe configuration (Fig.(2d)) does not show much resemblance to the experimental result. Its main band is much more two-dimensional, normal metal-like, comparable to the dispersion of a 2D tight-binding band.
Fig.(3a) plots the integrated spectral weight $`n(𝐤)`$ for the ”$`3+1`$” site-centered stripe configuration with $`t^{}=0`$. Although not as clear as in the Nd-LSCO ARPES experiment (from ref.), the ”Fermi surface” is rather one-dimensional in structure. Like in Nd-LSCO the low energy excitations (shown in Fig.(3b), calculated by integrating over a $`\mathrm{\Delta }\omega =0.2t`$ window below the Fermi-energy for each $`𝐤`$-point in the Brillouin zone) are located near the $`(\pm \pi ,0)`$, $`(0,\pm \pi )`$ points in momentum space. In Fig.(3b) we notice the $`8\times 8`$ square lattice of bright points. This is the repeated Brillouin zone of the supercell consisting of $`(3+1+3+1)\times N=8\times 8`$ lattice sites (due to the Néel order in $`x`$-direction). Clearly, the low energy excitations in the momentum space of the supercell are near the $`(\pm \pi ,0)`$ and $`(0,\pm \pi )`$ points as well. Fig.(3c) shows the integrated spectral weight $`n(𝐤)`$ of the ”$`2+2`$$`t`$-$`J`$ stripe calculation. Here, the CPT ”Fermi surface” is much more rounded, similar to the quasi-2D Fermi surface known from band calculations. The low energy excitations are located isotropically around the ”Fermi surface” as well (Fig.(3d)). The loss of one-dimensionality observed for this bond-centered ”$`2+2`$” stripe configuration is accompanied by a substantial enhancement of spectral weight near $`(\pi /2,\pi /2)`$. From these $`A(𝐤,\omega )`$ results we conclude that, at least in the Nd-LSCO system, the stripes are site-centered and of ”$`3+1`$” type.
With our technique we are able to resolve for each excitation, whether its main origin is from the insulating or the metallic part of the stripe configuration: In Fig.(4) we compare the spectra for the unperturbed stripe configuration with inter-cluster hopping set equal to zero (in Fig.(4a): solid curve for the AF domains, shaded curve for the 1D metal) with the result of our CPT calculation (solid line in Fig.(4b)) with inter-cluster hopping $`t`$ (for the Néel ordered ”$`3+1`$” configuration). All of the spectral weight in inverse photoemission ($`\omega >0`$) naturally stems from the chain since the half-filled 3-leg $`t`$-$`J`$ ladder does not have target states for an inverse photoemission process. The stripes are oriented along the $`y`$-direction. Therefore, we can conclude that peaks (in the calculation with $`t=0`$), that show a dispersion along $`(0,0)`$ to $`(\pi ,0)`$ direction stem from the AF domains (solid line in (4a)) whereas the metallic excitations prior to the mixing (shaded curve in (4a)) are dispersionless. By comparing the three curves, we therefore conclude that the sharp quasiparticle peak near $`(\pi ,0)`$ results from the mixing of the (dispersionless) $`(k_x,0)`$ minimum of the metallic spinon band and the top of the insulating valence band situated at $`(\pi ,0)`$. Going from $`(\pi ,0)`$ to $`(\pi ,\pi )`$, the metallic band becomes dispersive and crosses, in agreement with experiment the Fermi surface at $`k_y=\pi /4`$ (since it is quarter-filled). The dispersion of the insulating band, however, is in the opposite direction. For this reason, in the final spectrum, we observe that the sharp quasiparticle peak at $`(\pi ,0)`$ becomes dispersive going into $`(0,\pi )`$ direction and eventually crosses the Fermi surface, however, with diminishing weight due to the absence of mixing with the insulating band. This effect is best visible in the grayscale plot of Fig.(2b). This finding may serve to clarify questions raised in the experimental ref. concerning the origin of the quasiparticle peak at $`(\pi ,0)`$.
To summarize, the single-particle spectral weight $`A(𝐤,\omega )`$ was calculated for different stripe configurations employing an application of the cluster perturbation technique for inhomogeneous systems. This technique allows to obtain all $`𝐤`$-points and, therefore, allows for a detailed comparison with ARPES data. The $`A(𝐤,\omega )`$ results for the ”$`3+1`$” site-centered configuration display salient features observed in experiments such as a two-peak structure around $`(\pi ,0)`$ with a sharp excitation close to the Fermi-energy and a broader feature at higher binding energies, a quasi-1D distribution of spectral weight, and low energy excitations located around the $`(\pi ,0)`$-points in the Brillouin zone. The theoretical results suggest to rule out the alternative bond-centered ”$`2+2`$” configuration.
The authors acknowledge financial support from BMBF (05SB8WWA1) and DFG (HA1537/17-1). The calculations were carried out at the high-performance computing centers HLRS (Stuttgart) and LRZ (München).
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# External symmetry in general relativity
## 1 Introduction
In general relativity the development of the quantum field theory in curved spacetimes give rise to many difficult problems related to the physical interpretation of the one-particle quantum modes that may indicate how to quantize the field. This is because the form and the properties of the particular solutions of the free field equations, in the cases when these can be analytically solved , are strongly dependent on the procedure of separation of variables and, implicitly, on the choice of the local chart. Moreover, when the fields have spin the situation is more complicated since then the field equations and, therefore, the form of their particular solutions depend, in addition, on the tetrad gauge in which one works . In these conditions it would be helpful to use the traditional method of the quantum theory in flat spacetime based on the complete sets of commuting operators that determine the quantum modes as common eigenstates and give physical meaning to the constants of the separation of variables which are just the eigenvalues of these operators. A good step in this direction could be to proceed like in special relativity looking for the generators of the geometric symmetries similar to the familiar momentum, angular momentum and spin operators of the Poincaré covariant field theories .
However, the relativistic covariance in the sense of general relativity is too general to play the same role as the Lorentz or Poincaré covariance in special relativity. In other respects, the tetrad gauge invariance of the theories with spin represents another kind of general symmetry that is not able to produce itself conserved observables . For this reason we have to concentrate only upon some special transformations which should form a well-defined Lie group with significant parameterization from the geometric point of view. Obviously, these must be just the isometry transformations that point out the symmetry of the background giving us the specific Killing vectors . The physical fields take over this symmetry transforming according to appropriate representations of the isometry group. In the case of the scalar vector or tensor fields these representations are completely defined by the well-known transformations rules under coordinate transformations since the isometries are in fact particular automorphysms. It remains open the problem of the behavior under isometries of the fields with half integer spin which explicitly depend on the tetrad gauge fixing. Another important problem is how to define the generators of these representations for any spin. It is known that there is a standard operator-valued representation of the isometry group in the space of scalar functions whose generators can be written with the help of the Killing vectors in a similar manner as the orbital angular momentum operators of special relativity. Then it is natural to ask how could be defined the corresponding spin parts of the generators of the representations according which the fields with spin may transform.
Our aim here is to propose a way to solve these problems. We start with the idea that if we intend to study the symmetry of a physical theory we must take into consideration the whole geometric context, including the positions of the local frames given by the tetrad fields, since the spin is measured just with respect to the axes of these frames. Therefore, the symmetry transformations must preserve not only the form of the metric tensor but the tetrad gauge too. These have to be isometries combined with suitable tetrad gauge transformations in such a manner to leave invariant the tetrad field components. Thus we define the external symmetry group and we show that this is locally isomorphic with the isometry group, having the same structure constants. Moreover, there are arguments that in fact this is isomorphic with the universal covering group of the isometry one.
The next step is to define the operator-valued representations of the external symmetry group carried by spaces of fields with spin. We point out that these are induced by the linear finite-dimensional representations of the $`SL(2,C)`$ group. This is the motive why the symmetry transformations which leave invariant the field equations have generators with a composite structure. These have the usual orbital terms of the scalar representation and, in addition, specific spin terms which depend on the choice of the tetrad gauge even in the case of the fields with integer spin. In general, the spin and orbital terms do not commute to each other apart of some special gauge fixings where the fields transform manifestly covariant under external symmetry transformations.
These general results allow us to study two important examples, namely the central symmetry and the maximal symmetry of the de Sitter (dS) and anti-de Sitter (AdS) spacetimes. In the case of the central geometries we use central charts with Cartesian coordinates and the Cartesian tetrad gauge which allowed us recently to find new analytical solutions of the Dirac equation . We show that in this gauge fixing the central symmetry becomes global and, consequently, the spin terms of its generators are the same as those of special relativity . This is important from the technical point of view since in the largely used diagonal tetrad gauge in spherical coordinates we obtain that the spin terms are partially hidden. For the dS and AdS spacetimes we calculate the generators of the representations of the external symmetry group in central charts with our Cartesian gauge and in Minkowskian charts with another gauge where the fields behave manifestly covariant under the Lorentz symmetry .
We start in the second section with a brief review of the relativistic covariance and gauge symmetry which will be treated together introducing the group of the combined transformations defined as gauge transformations followed by authomorphysms. The next section is devoted to our approach. Therein we define the external symmetry transformations, we show that these form a group and we study the operator-valued representations of this group and its Lie algebra. In Sec.4 and 5 we discuss the mentioned examples.
We present our proposal at the level of the relativistic quantum mechanics in the sense of general relativity avoiding to consider the specific problems of the quantum field theory or to use too complicated mathematical methods. We work in natural units with $`\mathrm{}=c=1`$.
## 2 Relativistic covariance
In the Lagrangian field theory in curved spacetimes the relativistic covariant equations of scalar, vector or tensor fields arise from actions that are invariant under general coordinate transformations. Moreover, when the fields have spin in the sense of the $`SL(2,C)`$ symmetry then the action must be invariant under tetrad gauge transformations . The first step to our approach we propose here is to embed both these kind of transformations into new ones, called combined transformations, that will help us to understand the relativistic covariance in its most general terms.
### 2.1 Gauge transformations
Let us consider the curved spacetime $`M`$ and a local chart (natural frame) of coordinates $`x^\mu ,\mu =0,1,2,3`$. Given a gauge, we denote by $`e_{\widehat{\mu }}(x)`$ the tetrad fields that define the local frames, in each point $`x`$, and by $`\widehat{e}^{\widehat{\mu }}(x)`$ those of the corresponding coframes. These have the usual orthonormalization properties
$$\widehat{e}_\alpha ^{\widehat{\mu }}e_{\widehat{\nu }}^\alpha =\delta _{\widehat{\nu }}^{\widehat{\mu }},\widehat{e}_\alpha ^{\widehat{\mu }}e_{\widehat{\mu }}^\beta =\delta _\alpha ^\beta ,e_{\widehat{\mu }}e_{\widehat{\nu }}=\eta _{\widehat{\mu }\widehat{\nu }},\widehat{e}^{\widehat{\mu }}\widehat{e}^{\widehat{\nu }}=\eta ^{\widehat{\mu }\widehat{\nu }},$$
(1)
where $`\eta =`$diag$`(1,1,1,1)`$ is the Minkowki metric. From the line element
$$ds^2=\eta _{\widehat{\mu }\widehat{\nu }}d\widehat{x}^{\widehat{\mu }}d\widehat{x}^{\widehat{\nu }}=g_{\mu \nu }(x)dx^\mu dx^\nu ,$$
(2)
expressed in terms of 1-forms, $`d\widehat{x}^{\widehat{\mu }}=\widehat{e}_\nu ^{\widehat{\mu }}dx^\nu `$, we get the components of the metric tensor of the natural frame,
$$g_{\mu \nu }=\eta _{\widehat{\alpha }\widehat{\beta }}\widehat{e}_\mu ^{\widehat{\alpha }}\widehat{e}_\nu ^{\widehat{\beta }},g^{\mu \nu }=\eta ^{\widehat{\alpha }\widehat{\beta }}e_{\widehat{\alpha }}^\mu e_{\widehat{\beta }}^\nu .$$
(3)
These raise or lower the natural vector indices, i.e., the Greek ones ranging from 0 to 3, while for the local vector indices, denoted by hat Greeks and having the same range, we must use the Minkowski metric. The derivatives $`\widehat{}_{\widehat{\nu }}=e_{\widehat{\nu }}^\mu _\mu `$ satisfy the commutation rules
$$[\widehat{}_{\widehat{\mu }},\widehat{}_{\widehat{\nu }}]=e_{\widehat{\mu }}^\alpha e_{\widehat{\nu }}^\beta (\widehat{e}_{\alpha ,\beta }^{\widehat{\sigma }}\widehat{e}_{\beta ,\alpha }^{\widehat{\sigma }})\widehat{}_{\widehat{\sigma }}=C_{\widehat{\mu }\widehat{\nu }}^{\widehat{\sigma }}\widehat{}_{\widehat{\sigma }}$$
(4)
defining the Cartan coefficients which halp us to write the conecttion components in the local frames as
$$\widehat{\mathrm{\Gamma }}_{\widehat{\mu }\widehat{\nu }}^{\widehat{\sigma }}=e_{\widehat{\mu }}^\alpha e_{\widehat{\nu }}^\beta (\widehat{e}_\gamma ^{\widehat{\sigma }}\mathrm{\Gamma }_{\alpha \beta }^\gamma \widehat{e}_{\beta ,\alpha }^{\widehat{\sigma }})=\frac{1}{2}\eta ^{\widehat{\sigma }\widehat{\lambda }}(C_{\widehat{\mu }\widehat{\nu }\widehat{\lambda }}+C_{\widehat{\lambda }\widehat{\mu }\widehat{\nu }}+C_{\widehat{\lambda }\widehat{\nu }\widehat{\mu }}).$$
(5)
The notation $`\mathrm{\Gamma }_{\alpha \beta }^\gamma `$ stands for the usual Christoffel symbols involved in the formulas of the covariant derivatives $`_\mu =_{;\mu }`$ .
The Minkowski metric $`\eta _{\widehat{\mu }\widehat{\nu }}`$ remains invariant under the transformations of the gauge group of this metric, $`G(\eta )=O(3,1)`$. This has as subgroup the Lorentz group, $`L_+^{}`$, of the transformations $`\mathrm{\Lambda }[A(\omega )]`$ corresponding to the transformations $`A(\omega )SL(2,C)`$ through the canonical homomorphism . In the standard covariant parameterization, with the real parameters $`\omega ^{\widehat{\alpha }\widehat{\beta }}=\omega ^{\widehat{\beta }\widehat{\alpha }}`$, we have
$$A(\omega )=e^{\frac{i}{2}\omega ^{\widehat{\alpha }\widehat{\beta }}S_{\widehat{\alpha }\widehat{\beta }}},$$
(6)
where $`S_{\widehat{\alpha }\widehat{\beta }}`$ are the covariant basis-generators of the $`sl(2,C)`$ Lie algebra which satisfy
$$[S_{\widehat{\mu }\widehat{\nu }},S_{\widehat{\sigma }\widehat{\tau }}]=i(\eta _{\widehat{\mu }\widehat{\tau }}S_{\widehat{\nu }\widehat{\sigma }}\eta _{\widehat{\mu }\widehat{\sigma }}S_{\widehat{\nu }\widehat{\tau }}+\eta _{\widehat{\nu }\widehat{\sigma }}S_{\widehat{\mu }\widehat{\tau }}\eta _{\widehat{\nu }\widehat{\tau }}S_{\widehat{\mu }\widehat{\sigma }}).$$
(7)
For small values of $`\omega ^{\widehat{\alpha }\widehat{\beta }}`$ the matrix elements of the transformations $`\mathrm{\Lambda }`$ can be written as
$$\mathrm{\Lambda }[A(\omega )]_{\widehat{\nu }}^{\widehat{\mu }}=\delta _{\widehat{\nu }}^{\widehat{\mu }}+\omega _{\widehat{\nu }}^{\widehat{\mu }}+\mathrm{}.$$
(8)
Now we assume that $`M`$ is orientable and time-orientable such that $`L_+^{}`$ can be considered as the gauge group of the Minkowski metric . In these conditions the fields with spin can be defined as in the case of the flat spacetime, with the help of the finite-dimensional linear representations of the $`SL(2,C)`$ group . Then any field $`\psi _\rho :MV_\rho `$ is defined over $`M`$ with values in the vector space $`V_\rho `$ of the representation $`\rho `$ which determines the spin content of $`\psi _\rho `$. In the following we systematically use the bases of $`V_\rho `$ labeled only by spinor or vector local indices defined with respect to the axes of the local frames given by the tetrad fields. Generally, these will not be written explicitly apart the cases when this is requested by the concrete calculation needs.
The relativistic covariant field equations are derived from actions ,
$$𝒮[\psi _\rho ,e]=d^4x\sqrt{g}(\psi _\rho ,D_{\widehat{\mu }}\psi _\rho ),g=|det(g_{\mu \nu })|,$$
(9)
with covariant derivatives,
$$D_{\widehat{\alpha }}=\widehat{}_{\widehat{\alpha }}+\frac{i}{2}\rho (S_{\widehat{\gamma }}^{\widehat{\beta }})\widehat{\mathrm{\Gamma }}_{\widehat{\alpha }\widehat{\beta }}^{\widehat{\gamma }},$$
(10)
that assure the invariance of the whole theory under the gauge transformations,
$`\widehat{e}_\mu ^{\widehat{\alpha }}(x)`$ $``$ $`\widehat{e}_\mu ^{\widehat{\alpha }}(x)=\mathrm{\Lambda }[A(x)]_{\widehat{\beta }}^{\widehat{\alpha }}\widehat{e}_\mu ^{\widehat{\beta }}(x),`$
$`e_{\widehat{\alpha }}^\mu (x)`$ $``$ $`e_{}^{}{}_{\widehat{\alpha }}{}^{\mu }(x)=\mathrm{\Lambda }[A(x)]_{\widehat{\alpha }}^{\widehat{\beta }}e_{\widehat{\beta }}^\mu (x),`$ (11)
$`\psi _\rho (x)`$ $``$ $`\psi _\rho ^{}(x)=\rho [A(x)]\psi _\rho (x),`$
determined by the mappings $`A:MSL(2,C)`$ the values of which are the local $`SL(2,C)`$ transformations $`A(x)A[\omega (x)]`$. These mappings can be organized as a group, $`𝒢`$, with respect to the multiplication $`\times `$ defined as $`(A^{}\times A)(x)=A^{}(x)A(x)`$. The notation $`Id`$ stands for the mapping identity, $`Id(x)=1SL(2,C)`$, while $`A^1`$ is the inverse of $`A`$, $`(A^1)(x)=[A(x)]^1`$.
### 2.2 Combined transformations
The general coordinate transformations are automorphysms of $`M`$ which, in the passive mode, can be seen as changes of the local charts corresponding to the same domain of $`M`$ . If $`x`$ and $`x^{}`$ are the coordinates of a point in two different charts then there is a mapping $`\varphi `$ between these charts giving the coordinate transformation, $`xx^{}=\varphi (x)`$. These transformations form a group with respect to the composition of mappings, $``$, defined as usual, i.e. $`(\varphi ^{}\varphi )(x)=\varphi ^{}[\varphi (x)]`$. We denote this group by $`𝒜`$, its identity map by $`id`$ and the inverse mapping of $`\varphi `$ by $`\varphi ^1`$.
The automothysms change all the components carrying natural indices including those of the tetrad fields . This means that, from the physical point of view, these transformations may change the positions of the local frames with respect to the natural ones. If we assume that the physical experiment makes reference to the axes of the local frame then it could appear situations when several correction of the positions of the local frames should be needed before (or after) each general coordinate transformation. Obviously, these have to be done with the help of suitable gauge transformation associated to the authomorphysms. Thus we arrive to the necessity of introducing the combined transformations denoted by $`(A,\varphi )`$ and defined as gauge transformations, given by $`A𝒢`$, followed by automorphysms, $`\varphi 𝒜`$. In this new notation the pure gauge transformations will appear as $`(A,id)`$ while the automorphysms will be denoted from now by $`(Id,\varphi )`$.
The effect of a combined transformation $`(A,\varphi )`$ upon our basic fields, $`\psi _\rho ,e`$ and $`\widehat{e}`$ is $`xx^{}=\varphi (x),e(x)e^{}(x^{}),\widehat{e}(x)\widehat{e}^{}(x^{})`$ and $`\psi _\rho (x)\psi _\rho ^{}(x^{})=\rho [A(x)]\psi _\rho (x)`$ where $`e^{}`$ are the transformed tetrads of the components
$$e_{\widehat{\alpha }}^\mu [\varphi (x)]=\mathrm{\Lambda }[A(x)]_{\widehat{\alpha }}^{\widehat{\beta }}e_{\widehat{\beta }}^\nu (x)\frac{\varphi ^\mu (x)}{x^\nu },$$
(12)
while the components of $`\widehat{e}^{}`$ have to be calculated according to Eqs.(1). Thus we have written down the most general transformation laws that leave invariant the action in the sense that $`𝒮[\psi _\rho ^{},e^{}]=𝒮[\psi _\rho ,e]`$. The field equations derived from $`𝒮`$, written in operator form as $`(E_\rho \psi _\rho )(x)=0`$, covariantly transform under these transformations,
$$(E_\rho \psi _\rho )(x)(E_\rho ^{}\psi _\rho ^{})(x^{})=\rho [A(x)](E_\rho \psi _\rho )(x),$$
(13)
since the operators $`E_\rho `$ involve covariant derivatives .
The association among the transformations of the groups $`𝒢`$ and $`𝒜`$ must lead to a new group with a specific multiplication. In order to find how looks this new operation it is convenient to introduce the compositions among the mappings $`A`$ and $`\varphi `$ (taken only in this order) giving new mappings, $`A\varphi `$, defined as $`(A\varphi )(x)=A[\varphi (x)]`$. The calculation rules $`Id\varphi =Id`$, $`Aid=A`$ and $`(A^{}\times A)\varphi =(A^{}\varphi )\times (A\varphi )`$ are obvious. With these ingredients we define the new multiplication
$$(A^{},\varphi ^{})(A,\varphi )=((A^{}\varphi )\times A,\varphi ^{}\varphi ).$$
(14)
It is clear that now the identity is $`(Id,id)`$ while the inverse of a pair $`(A,\varphi )`$ reads
$$(A,\varphi )^1=(A^1\varphi ^1,\varphi ^1).$$
(15)
First of all we observe that the operation $``$ is well-defined and represents the composition among the combined transformations since these can be expressed, according to their definition, as $`(A,\varphi )=(Id,\varphi )(A,id)`$. Furthermore, we can convince ourselves that if we perform successively two arbitrary combined transformations, $`(A,\varphi )`$ and $`(A^{},\varphi ^{})`$, then the resulting transformation is just $`(A^{},\varphi ^{})(A,\varphi )`$ as given by Eq.(14). This means that the combined transformations form a group with respect to the multiplication $``$. It is not difficult to verify that this group, denoted by $`\stackrel{~}{𝒢}`$, is the semidirect product $`\stackrel{~}{𝒢}=𝒢\mathrm{}𝒜`$ where $`𝒢`$ is the invariant subgroup while $`𝒜`$ is an usual one.
In the theories involving only vector and tensor fields we do not need to use the combined transformations defined above since the theory is independent on the positions of the local frames. This can be easily shown even in our approach where we use field components with local indices. Indeed, if we perform a combined transformation $`(A,\varphi )`$ then any tensor field of rank $`(p,q)`$,
$$\psi _{\widehat{\beta }_1,\widehat{\beta }_2,\mathrm{},\widehat{\beta }_q}^{\widehat{\alpha }_1,\widehat{\alpha }_2,\mathrm{},\widehat{\alpha }_p}=\widehat{e}_{\mu _1}^{\widehat{\alpha }_1}\mathrm{}\widehat{e}_{\mu _p}^{\widehat{\alpha }_p}e_{\widehat{\beta }_1}^{\nu _1}\mathrm{}e_{\widehat{\beta }_q}^{\nu _q}\psi _{\nu _1,\nu _2,\mathrm{},\nu _q}^{\mu _1,\mu _2,\mathrm{},\mu _p},$$
(16)
transforms according to the representation
$$\rho _{\widehat{\alpha }_1,\widehat{\alpha }_2,\mathrm{},\widehat{\alpha }_p;\widehat{\beta }_1^{},\widehat{\beta }_2^{},\mathrm{},\widehat{\beta }_q^{}}^{\widehat{\beta }_1,\widehat{\beta }_2,\mathrm{},\widehat{\beta }_q;\widehat{\alpha }_1^{},\widehat{\alpha }_2^{},\mathrm{},\widehat{\alpha }_p^{}}(A)=\mathrm{\Lambda }_{\widehat{\beta }_1^{}}^{\widehat{\beta }_1}(A)\mathrm{}\mathrm{\Lambda }_{\widehat{\alpha }_1}^{\widehat{\alpha }_1^{}}(A)\mathrm{},$$
(17)
such that the resulting transformation law of the components carrying natural indices,
$$\psi _{\nu _1,\mathrm{}}^{\mu _1,\mathrm{}}(x^{})=\frac{x^{\mu _1}}{x^{\sigma _1}}\mathrm{}\frac{x^{\tau _1}}{x^{\nu _1}}\mathrm{}\psi _{\tau _1,\mathrm{}.}^{\sigma _1,\mathrm{}.}(x),$$
(18)
is just the familiar one . In other words, in this case the effect of the combined transformations reduces to that of their authomorohysms. However, when the half integer spin fields are involved this is no more true and we must use the combined transformations of $`\stackrel{~}{𝒢}`$ if we want to keep under control the positions of the local frames.
## 3 External symmetry
In general, the symmetry of any manifold $`M`$ is given by its isometry group whose transformations leave invariant the metric tensor in any chart. The scalar field transforms under isometries according to the standard scalar representation generated by the orbital generators related to the Killing vectors of $`M`$ . In the following we propose a possible generalization of this theory of symmetry to the fields with spin, defining the external symmetry group and its representations.
### 3.1 Isometries
There are conjectures when several coordinate transformations, $`xx^{}=\varphi _\xi (x)`$, depend on $`N`$ independent real parameters, $`\xi ^a`$ ($`a,b,c\mathrm{}=1,2,\mathrm{},N`$), such that $`\xi =0`$ corresponds to the identity map, $`\varphi _{\xi =0}=id`$. The set of these mappings can be organized as a Lie group , $`G𝒢`$, if they accomplish the composition rule
$$\varphi _\xi ^{}\varphi _\xi =\varphi _{f(\xi ^{},\xi )},$$
(19)
where the functions $`f:G\times GG`$ define the group multiplication. These must satisfy $`f^a(0,\xi )=f^a(\xi ,0)=\xi ^a`$ and $`f^a(\xi ^1,\xi )=f^a(\xi ,\xi ^1)=0`$ where $`\xi ^1`$ are the parameters of the inverse mapping of $`\varphi _\xi `$, $`\varphi _{\xi ^1}=\varphi _\xi ^1`$. Moreover, the structure constants of $`G`$ can be calculated as
$$c_{abc}=\left(\frac{f^c(\xi ,\xi ^{})}{\xi ^a\xi ^b}\frac{f^c(\xi ,\xi ^{})}{\xi ^b\xi ^a}\right)_{|\xi =\xi ^{}=0}.$$
(20)
For small values of the group parameters the infinitesimal transformations, $`x^\mu x^\mu =x^\mu +\xi ^ak_a^\mu (x)+\mathrm{}`$, are given by the vectors $`k_a`$ whose components,
$$k_a^\mu =\frac{\varphi _\xi ^\mu }{\xi ^a}_{|\xi =0},$$
(21)
satisfy the identities
$$k_a^\mu k_{b,\mu }^\nu k_b^\mu k_{a,\mu }^\nu +c_{abc}k_c^\nu =0,$$
(22)
resulting from Eqs.(19) and (20).
In the following we restrict ourselves to consider only the isometry transformations, $`x^{}=\varphi _\xi (x)`$, which leave invariant the components of the metric tensor , i.e.
$$g_{\alpha \beta }(x^{})\frac{x^\alpha }{x^\mu }\frac{x^\beta }{x^\nu }=g_{\mu \nu }(x).$$
(23)
These form the isometry group $`GI(M)`$ which is the Lie group giving the symmetry of the spacetime $`M`$. We consider that this has $`N`$ independent parameters and, therefore, $`k_a,a=1,2,\mathrm{}N`$, are independent Killing vectors (which satisfy $`k_{a\mu ;\nu }+k_{a\nu ;\mu }=0`$). Then their corresponding Lie derivatives form a basis of the Lie algebra $`i(M)`$ of the group $`I(M)`$ .
However, in practice we are interested to find the operators of the relativistic quantum theory related to these geometric objects which describe the symmetry of the background. For this reason we focus upon the operator-valued representations of the group $`I(M)`$ and its algebra. The scalar field $`\psi :MC`$ transforms under isometries as $`\psi (x)\psi ^{}[\varphi _\xi (x)]=\psi (x)`$. This rule defines the representation $`\varphi _\xi T_\xi `$ of the group $`I(M)`$ whose operators have the action $`\psi ^{}=T_\xi \psi =\psi \varphi _\xi ^1`$. Hereby it results that the operators of infinitesimal transformations, $`T_\xi =1i\xi ^aL_a+\mathrm{}`$, depend on the basis-generators,
$$L_a=ik_a^\mu _\mu ,a=1,2,\mathrm{},N,$$
(24)
which are completely determined by the Killing vectors. From Eq.(22) we see that they obey the commutation rules
$$[L_a,L_b]=ic_{abc}L_c,$$
(25)
given by the structure constants of $`I(M)`$. In other words they form a basis of the operator-valued representation of the Lie algebra $`i(M)`$ in a carrier space of scalar fields. Notice that in the usual quantum mechanics the operators similar to the generators $`L_a`$ are called often orbital generators.
### 3.2 The group of external symmetry
Now the problem is how may transform under isometries the whole geometric framework of the theories with spin where we explicitly use the local frames. Since the isometry is a general coordinate transformation it changes the relative positions of the local and natural frames. This fact may be an impediment when one intends to study the symmetries of the theories with spin induced by those of the background. For this reason it is natural to suppose that the good symmetry transformations we need are combined transformations in which the isometries are preceded by appropriate gauge transformations such that not only the form of the metric tensor should be conserved but the form of the tetrad field components too.
Thus we arrive at the main point of our proposal. We introduce the external symmetry transformations, $`(A_\xi ,\varphi _\xi )`$, as combined transformations involving isometries and corresponding gauge transformations necessary to preserve the gauge. We assume that in a fixed gauge, given by the tetrad fields $`e`$ and $`\widehat{e}`$, $`A_\xi `$ is defined by
$$\mathrm{\Lambda }[A_\xi (x)]_{\widehat{\beta }}^{\widehat{\alpha }}=\widehat{e}_\mu ^{\widehat{\alpha }}[\varphi _\xi (x)]\frac{\varphi _\xi ^\mu (x)}{x^\nu }e_{\widehat{\beta }}^\nu (x),$$
(26)
with the supplementary condition $`A_{\xi =0}(x)=1SL(2,C)`$. Since $`\varphi _\xi `$ is an isometry Eq.(23) guarantees that $`\mathrm{\Lambda }[A_\xi (x)]L_+^{}`$ and, implicitly, $`A_\xi (x)SL(2,C)`$. Then the transformation laws of our fields are
$$(A_\xi ,\varphi _\xi ):\begin{array}{ccccc}\hfill x& \hfill & \hfill x^{}& =& \varphi _\xi (x),\hfill \\ \hfill e(x)& \hfill & \hfill e^{}(x^{})& =& e[\varphi _\xi (x)],\hfill \\ \hfill \widehat{e}(x)& \hfill & \hfill \widehat{e}^{}(x^{})& =& \widehat{e}[\varphi _\xi (x)],\hfill \\ \hfill \psi _\rho (x)& \hfill & \hfill \psi _\rho ^{}(x^{})& =& \rho [A_\xi (x)]\psi _\rho (x).\hfill \end{array}$$
(27)
The mean virtue of these transformations are that they leave invariant the operators of the field equations, $`E_\rho `$, since the components of the tetrad fields and, consequently, the covariant derivatives do not change their form.
For small $`\xi ^a`$ the covariant $`SL(2,C)`$ parameters of $`A_\xi (x)A[\omega _\xi (x)]`$ can be written as $`\omega _\xi ^{\widehat{\alpha }\widehat{\beta }}(x)=\xi ^a\mathrm{\Omega }_a^{\widehat{\alpha }\widehat{\beta }}(x)+\mathrm{}`$ where, according to Eqs.(6), (8) and (26), we have
$$\mathrm{\Omega }_a^{\widehat{\alpha }\widehat{\beta }}\frac{\omega _\xi ^{\widehat{\alpha }\widehat{\beta }}}{\xi ^a}_{|\xi =0}=\left(\widehat{e}_\mu ^{\widehat{\alpha }}k_{a,\nu }^\mu +\widehat{e}_{\nu ,\mu }^{\widehat{\alpha }}k_a^\mu \right)e_{\widehat{\lambda }}^\nu \eta ^{\widehat{\lambda }\widehat{\beta }}.$$
(28)
We must specify that these functions are antisymmetric if and only if $`k_a`$ are Killing vectors. This indicates that the association among isometries and the gauge transformations defined by Eq.(26) is correct.
It remains to show that the transformations $`(A_\xi ,\varphi _\xi )`$ form a Lie group related to $`I(M)`$. Starting with Eq.(26) we find that
$$(A_\xi ^{}\varphi _\xi )\times A_\xi =A_{f(\xi ^{},\xi )},$$
(29)
and, according to Eqs.(19) and (29), we obtain
$$(A_\xi ^{},\varphi _\xi ^{})(A_\xi ,\varphi _\xi )=(A_{f(\xi ^{},\xi )},\varphi _{f(\xi ^{},\xi )}),$$
(30)
and $`(A_{\xi =0},\varphi _{\xi =0})=(Id,id)`$. Thus we have shown that the pairs $`(A_\xi ,\varphi _\xi )`$ form a Lie group with respect to the operation $``$. We say that this is the external symmetry group of $`M`$ and we denote it by $`S(M)\stackrel{~}{𝒢}`$. From Eq.(30) we understand that $`S(M)`$ is locally isomorphic with $`I(M)`$ and, therefore, the Lie algebra of $`S(M)`$, denoted by $`s(M)`$, is isomorphic with $`i(M)`$ having the same structure constants. In our opinion, $`S(M)`$ must be isomorphic with the universal covering group of $`I(M)`$ since it has anyway the topology induced by $`SL(2,C)`$ which is simply connected. In general, the number of group parameters of $`I(M)`$ or $`S(M)`$ (which is equal to the number of the independent Killing vectors of $`M`$) can be $`0N10`$.
The form of the external symmetry transformations is strongly dependent on the choice of the local chart as well as that of the tetrad gauge. If we change simultaneously the gauge and the coordinates with the help of a combined transformation $`(A,\varphi )`$ then each $`(A_\xi ,\varphi _\xi )S(M)`$ transforms as
$$(A_\xi ,\varphi _\xi )(A_\xi ^{},\varphi _\xi ^{})=(A,\varphi )(A_\xi ,\varphi _\xi )(A,\varphi )^1$$
(31)
which means that
$`A_\xi ^{}`$ $`=`$ $`\left\{\left[\left(A\varphi _\xi \right)\times A_\xi \right]\times A^1\right\}\varphi ^1,`$ (32)
$`\varphi _\xi ^{}`$ $`=`$ $`\left(\varphi \varphi _\xi \right)\varphi ^1.`$ (33)
### 3.3 Representations
The last of Eqs.(27) which gives the transformation law of the field $`\psi _\rho `$ defines the operator-valued representation $`(A_\xi ,\varphi _\xi )T_\xi ^\rho `$ of the group $`S(M)`$,
$$(T_\xi ^\rho \psi _\rho )[\varphi _\xi (x)]=\rho [A_\xi (x)]\psi _\rho (x).$$
(34)
The invariance of the field equations under these transformations requires to have
$$T_\xi ^\rho E_\rho (T_\xi ^\rho )^1=E_\rho .$$
(35)
Since $`A_\xi (x)SL(2,C)`$ we say that this representation is induced by the representation $`\rho `$ of $`SL(2,C)`$ . As we have shown in Sec.2.2, if $`\rho `$ is a vector or tensor representation (having only integer spin components) then the effect of the transformation (34) upon the components carrying natural indices is due only to $`\varphi _\xi `$. However, for the representations with half integer spin the presence of $`A_\xi `$ is crucial since there are no natural indices. Moreover, this allows us to define the generators of the representations (34) for any spin.
The basis-generators of the representations of the Lie algebra $`s(M)`$ are the operators
$$X_a^\rho =i\frac{T_\xi ^\rho }{\xi ^a}_{|\xi =0}=S_a^\rho +L_a,$$
(36)
which appear as sums among the orbital generators defined by Eq.(24) and the spin terms which have the action
$$(S_a^\rho \psi _\rho )(x)=\rho [S_a(x)]\psi _\rho (x).$$
(37)
This is determined by the form of the local $`sl(2,C)`$ generators,
$$S_a(x)=i\frac{A_\xi (x)}{\xi ^a}_{|\xi =0}=\frac{1}{2}\mathrm{\Omega }_a^{\widehat{\alpha }\widehat{\beta }}(x)S_{\widehat{\alpha }\widehat{\beta }},$$
(38)
which depend on the functions (28). Furthermore, if we derive Eq.(29) with respect to $`\xi `$ and $`\xi ^{}`$ then from Eqs.(8), (20) and (28), after a few manipulations, we find the identities
$$\eta _{\widehat{\alpha }\widehat{\beta }}\left(\mathrm{\Omega }_a^{\widehat{\alpha }\widehat{\mu }}\mathrm{\Omega }_b^{\widehat{\beta }\widehat{\nu }}\mathrm{\Omega }_b^{\widehat{\alpha }\widehat{\mu }}\mathrm{\Omega }_a^{\widehat{\beta }\widehat{\nu }}\right)+k_a^\mu \mathrm{\Omega }_{b,\mu }^{\widehat{\mu }\widehat{\nu }}k_b^\mu \mathrm{\Omega }_{a,\mu }^{\widehat{\mu }\widehat{\nu }}+c_{abc}\mathrm{\Omega }_c^{\widehat{\mu }\widehat{\nu }}=0.$$
(39)
Hereby it results that
$$[S_a^\rho ,S_b^\rho ]+[L_a,S_b^\rho ][L_b,S_a^\rho ]=ic_{abc}S_c^\rho ,$$
(40)
and, according to Eq.(25), we find the expected commutation rules
$$[X_a^\rho ,X_b^\rho ]=ic_{abc}X_c^\rho .$$
(41)
Thus we have obtained the basis-generators of a operator-valued representation of $`s(M)`$ induced by the linear representation $`\rho `$ of $`sl(2,C)`$. All the operators of this representation commute with the operator $`E_\rho `$ since, according to Eqs.(35) and (36), we have
$$[E_\rho ,X_a^\rho ]=0,a=1,2,\mathrm{},N.$$
(42)
Therefore, for defining quantum modes we can use the set of commuting operators containing the Casimir operators those of the Cartan subalgebra and $`E_\rho `$.
The action of the operators $`X_a^\rho `$ depends on the choice of many elements: the natural coordinates, the tetrad gauge, the group parameterization and the representation $`\rho `$. What is important here is that they are strongly dependent on the tetrad gauge fixing even in the case of the representations with integer spin. This is because the covariant parameterization of the $`sl(2,C)`$ algebra is defined with respect to the axes of the local frames. In general, if we consider the representation $`(A_\xi ,\varphi _\xi )T_\xi ^\rho `$ and we perform a combined transformation (31) then it results the equivalent representation, $`(A_\xi ^{},\varphi _\xi ^{})T_\xi ^\rho `$. Its generators calculated from Eqs.(32) indicate that in this case the equivalence relations are much more complicated than those of the usual theory of linear representations. Without to enter in other technical details we specify that if we change only the gauge with the help of the transformation $`(A,id)`$ then the local $`sl(2,C)`$ generators (38) transform as
$`S_a(x)S_a^{}(x)`$ $`=`$ $`A(x)S_a(x)A(x)^1`$ (43)
$`+k_a^\sigma (x)\mathrm{\Lambda }[A(x)]_{\widehat{\alpha }\widehat{\mu },\sigma }\mathrm{\Lambda }[A(x)]_{\widehat{\beta }}^{\widehat{\mu }}S^{\widehat{\alpha }\widehat{\beta }},`$
while the orbital parts do not change their form. This means that the gauge transformations change, in addition, the commutation relations among the spin and orbital parts of the generators $`X_a^\rho `$. Hence we can draw the conclusion that the choice of the tetrad gauge which defines the local frames may have important consequences upon the measurement of the local spin effects.
There are gauge fixings where the local $`sl(2,C)`$ generators $`S_a(x)`$, $`a=1,2,\mathrm{},n`$ ($`nN`$), corresponding to a subgroup $`HS(M)`$, are independent on $`x`$ and, therefore, $`[S_a^\rho ,L_b]=0`$ for all $`a=1,2,\mathrm{},n`$ and $`b=1,2,\mathrm{},N`$. Then the operators $`S_a^\rho `$ are just the basis-generators of an usual linear representation of $`H`$ and the field $`\psi _\rho `$ behaves manifestly covariant under the external symmetry transformations of this subgroup. Of course, when $`H=S(M)`$ we say simply that the field $`\psi _\rho `$ is manifest covariant.
The simplest examples are the manifest covariant fields of special relativity. Here since the spacetime $`M`$ is flat the metric in Cartesian coordinates is $`g_{\mu \nu }=\eta _{\mu \nu }`$ and one can use the inertial (local) frames with $`e_\nu ^\mu =\widehat{e}_\nu ^\mu =\delta _\nu ^\mu `$. Then the isometries are just the transformations $`x^{}=\mathrm{\Lambda }[A(\omega )]xa`$ of the Poincaré group, $`𝒫_+^{}=T(4)\mathrm{}L_+^{}`$ . If we denote by $`\xi ^{(\mu \nu )}=\omega ^{\mu \nu }`$ the $`SL(2,C)`$ parameters and by $`\xi ^{(\mu )}=a^\mu `$ those of the translation group $`T(4)`$, then we find that $`S(M)\stackrel{~}{𝒫}_+^{}=T(4)\mathrm{}SL(2,C)`$ is just the universal covering group of $`𝒫_+^{}`$. Furthermore, it is a simple exercise to calculate the basis-generators
$`X_{(\mu )}^\rho `$ $`=`$ $`i_\mu ,`$ (44)
$`X_{(\mu \nu )}^\rho `$ $`=`$ $`\rho (S_{\mu \nu })+i(\eta _{\mu \alpha }x^\alpha _\nu \eta _{\nu \alpha }x^\alpha _\mu ),`$ (45)
which show us that $`\psi _\rho `$ transforms manifestly covariant.
In general, there are many cases of curved spacetimes for which one can choose suitable local frames allowing one to introduce manifest covariant fields with respect to a subgroup $`HS(M)`$ or even the whole group $`S(M)`$. In our opinion, this is possible only when $`H`$ or $`S(M)`$ are at most subgroups of $`\stackrel{~}{𝒫}_+^{}`$.
## 4 The central symmetry
Let us take as first example the spacetimes $`M`$ which have spherically symmetric static chart that will be referred here as central charts. These manifolds have the isometry group $`I(M)=T(1)SO(3)`$ of time translations and space rotations.
### 4.1 Central charts
In a central chart with Cartesian coordinates $`x^0=t`$ and $`x^i`$ ($`i,j,k\mathrm{}=1,2,3`$), the metric tensor is time-independent and transforms manifestly covariant under the rotations $`RSO(3)`$ of the space coordinates,
$$t^{}=t,x^i=R_j^i(\omega )x^j=x^i+\omega _j^ix^j+\mathrm{},$$
(46)
denoted simply by $`xx^{}=Rx`$. Here the most general form of the line element,
$$ds^2=g_{\mu \nu }(x)dx^\mu dx^\nu =A(r)dt^2[B(r)\delta _{ij}+C(r)x^ix^j]dx^idx^j,$$
(47)
may involve three functions, $`A`$, $`B`$ and $`C`$, depending only on the Euclidean norm of $`\stackrel{}{x}`$, $`r=|\stackrel{}{x}|`$. In applications it is convenient to replace these functions by new ones, $`u`$, $`v`$ and $`w`$, defined as
$$A=w^2,B=\frac{w^2}{v^2},C=\frac{w^2}{r^2}\left(\frac{1}{u^2}\frac{1}{v^2}\right).$$
(48)
Other useful central charts are those with spherical coordinates, $`r`$, $`\theta `$, $`\varphi `$, commonly associated with the Cartesian space ones. Here the line elements are
$$ds^2=w^2dt^2\frac{w^2}{u^2}dr^2\frac{w^2}{v^2}r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2).$$
(49)
In these charts we see that the advantage of the new functions we have introduced is of simple transformation laws under the isotropic dilatations which change only the radial coordinate, $`rr^{}(r)`$, without to affect the central symmetry of the line element. These transformations,
$$u^{}(r^{})=u(r)\left|\frac{dr^{}(r)}{dr}\right|,v^{}(r^{})=v(r)\frac{r^{}(r)}{r},w^{}(r^{})=w(r),$$
(50)
allow one to choose desired forms for the functions $`u,v`$ and $`w`$.
### 4.2 The Cartesian gauge
The Cartesian gauge in central charts was mentioned long time ago but it is less used in concrete problems since it leads to complicated calculations in spherical coordinates. However, in Cartesian coordinates this gauge has the advantage of explicitly pointing out the global central symmetry of the manifold. In Refs. we have proposed a version of Cartesian gauge in central charts with Cartesian coordinates that preserve the manifest covariance under rotations (46) in the sense that the corresponding 1-forms transform as
$$d\widehat{x}^{\widehat{\mu }}d\widehat{x}^{\widehat{\mu }}=\widehat{e}_\alpha ^{\widehat{\mu }}(x^{})dx^\alpha =(Rd\widehat{x})^{\widehat{\mu }}.$$
(51)
If the line element has the form (47) then the most general choice of the tetrad fields with the above property is
$`\widehat{e}_0^0=\widehat{a}(r),\widehat{e}_i^0=\widehat{e}_0^i=0,\widehat{e}_j^i=\widehat{b}(r)\delta _{ij}+\widehat{c}(r)x^ix^j+\widehat{d}(r)ϵ_{ijk}x^k,`$ (52)
$`e_0^0=a(r),e_i^0=e_0^i=0,e_j^i=b(r)\delta _{ij}+c(r)x^ix^j+d(r)ϵ_{ijk}x^k,`$ (53)
where, according to (3), (47) and (48), we must have
$`\widehat{a}=w,\widehat{b}={\displaystyle \frac{w}{v}}\mathrm{cos}\alpha ,\widehat{c}={\displaystyle \frac{1}{r^2}}\left({\displaystyle \frac{w}{u}}{\displaystyle \frac{w}{v}}\mathrm{cos}\alpha \right),\widehat{d}={\displaystyle \frac{1}{r}}{\displaystyle \frac{w}{v}}\mathrm{sin}\alpha ,`$ (54)
$`a={\displaystyle \frac{1}{w}},b={\displaystyle \frac{v}{w}}\mathrm{cos}\alpha ,c={\displaystyle \frac{1}{r^2}}\left({\displaystyle \frac{u}{w}}{\displaystyle \frac{v}{w}}\mathrm{cos}\alpha \right),d={\displaystyle \frac{1}{r}}{\displaystyle \frac{v}{w}}\mathrm{sin}\alpha .`$ (55)
The angle $`\alpha `$ is an arbitrary function of $`r`$ which is not explicitly involved in the expression of the metric tensor since it represents the angle of an arbitrary rotation of the local frame around the direction of $`\stackrel{}{x}`$, that does not change the relative position of $`\stackrel{}{x}`$ with respect to this frame.
When one defines the metric tensor such that $`g_{\mu \nu }^{}{}_{|r=0}{}^{}=\eta _{\mu \nu }`$ then $`u(0)^2=v(0)^2=w(0)^2=1`$. Moreover, it is natural to take $`\alpha (0)=0`$. In other respects, from Eqs.(54) and (55) we see that the function $`w`$ must be positively defined in order to keep the same sense for the time axes of the natural and local frames. In addition, it is convenient to consider that the function $`u`$ is positively defined too. However, the function $`v=\eta _P|v|`$ has the sign given by the relative parity $`\eta _P`$ which takes the value $`\eta _P=1`$ when the space axes of the local frame at $`x=0`$ are parallel with those of the natural frame, and $`\eta _P=1`$ if these are antiparalel.
Now we have all the elements we need to calculate the generators of the representations $`T^\rho `$ of the group $`S(M)`$. If we denote by $`\xi ^{(0)}`$ the parameter of the time translations and by $`\xi ^{(i)}=ϵ_{ijk}\omega ^{jk}/2`$ the parameters of the rotations (46), we find that the local $`sl(2,C)`$ generators of Eq.(38) are just the $`su(2)`$ ones, i.e. $`S_{(ij)}(x)=S_{ij}`$, such that the basis-generators read
$$X_{(0)}^\rho =i_t,X_{(i)}^\rho =\frac{1}{2}ϵ_{ijk}\rho (S_{jk})+L_{(i)}$$
(56)
where $`L_{(i)}=iϵ_{ijk}x^j_k`$ are the usual components of the orbital angular momentum. Thus we obtain that the group $`S(M)=T(1)SU(2)`$ is the universal covering group of $`I(M)`$. The physical significance of the basis-generators is the usual one, namely $`X_{(0)}^\rho `$ is the Hamiltonian operator while $`X_{(i)}^\rho J_{(i)}^\rho `$ are the components of the whole angular momentum operator of the field $`\psi _\rho `$ which transforms manifestly covariant.
We can conclude that, in our Cartesian gauge, the local frames play the same role as the usual Cartesian rest frames of the central sources in flat spacetime since their axes are just those of projections of the angular momenta.
### 4.3 The diagonal gauge
In other gauge fixings the basis-generators are quite different. A tetrad gauge largely used in central charts with spherical coordinates is the diagonal gauge defined by the 1-forms
$$d\widehat{x}_s^0=wdt,d\widehat{x}_s^1=\frac{w}{u}dr,d\widehat{x}_s^2=r\frac{w}{v}d\theta ,d\widehat{x}_s^3=r\frac{w}{v}\mathrm{sin}\theta d\varphi .$$
(57)
In this gauge the angular momentum operators of the canonical basis (where $`J_{(\pm )}=J_{(1)}\pm iJ_{(2)}`$) are
$$J_{(\pm )}^\rho =\frac{e^{\pm i\varphi }}{\mathrm{sin}\theta }\rho (S_{23})+L_{(\pm )},J_{(3)}^\rho =L_{(3)}.$$
(58)
Thus we obtain a representation of $`SU(2)`$ where the spin terms do not commute with the orbital ones and, therefore, the field $`\psi _\rho `$ does not transform manifestly covariant under rotations. In this case we can say that the spin part of the central symmetry remains partially hidden because of the diagonal gauge which determines special positions of the local frames with respect to the natural one. However, when this is an impediment one can change anytime this gauge into the Cartesian one by using a simple local rotation. For the flat spacetimes these transformations and their effects upon the Dirac equation are studied in Ref.. We note that the form of the spin generators as well as that of the mentioned rotation depend on the enumeration of the 1-forms (57).
## 5 The dS and AdS symmetries
The backgrounds with highest external symmetry are the dS and the AdS spacetimes. We shall briefly discuss simultaneously both these manifolds which will be denoted by $`M_ϵ`$ where $`ϵ=1`$ for dS case and $`ϵ=1`$ for AdS one. Our goal here is to calculate the generators of the representations of the group $`S(M_ϵ)`$ induced by those of $`SL(2,C)`$.
The dS and AdS spacetimes are hyperboloids in the $`(4+1)`$ or $`(3+2)`$-dimensional flat spacetimes, $`M_ϵ^5`$, of coordinates $`Z^A,A,B,\mathrm{}=0,1,2,3,5`$, and the metric $`\eta (ϵ)=\mathrm{diag}(1,1,1,1,ϵ)`$. The equation of the hyperboloid of radius $`r_0=1/\widehat{\omega }`$ reads
$$\eta _{AB}(ϵ)Z^AZ^B=ϵr_{0}^{}{}_{}{}^{2}.$$
(59)
From their definitions it results that the dS or AdS spacetimes are homogeneous spaces of the pseudo-orthogonal groups $`SO(4,1)`$ or $`SO(3,2)`$ which play the role of gauge groups of the metric $`\eta (ϵ)`$ (for $`ϵ=1`$ and $`ϵ=1`$ respectively) and represent just the isometry groups of these manifolds, $`G[\eta (ϵ)]=I(M_ϵ)`$. Then it is natural to use the covariant real parameters $`\omega ^{AB}=\omega ^{BA}`$ since in this parameterization the orbital basis-generators of the representations of $`G[\eta (ϵ)]`$ carried by the spaces of the functions over $`M_ϵ^5`$ have the usual form
$$L_{AB}^5=i\left[\eta _{AC}(ϵ)Z^C_B\eta _{BC}(ϵ)Z^C_A\right].$$
(60)
They will give us directly the orbital basis-generators of the representations of $`S(M_ϵ)`$ in the carrier spaces of the functions defined over dS or AdS spacetimes.
### 5.1 Central charts
The hyperboloid equation can be solved in Cartesian dS/AdS coordinates, $`x^0=t`$ and $`x^i`$ ($`i=1,2,3`$), which satisfy
$`Z^5`$ $`=`$ $`\widehat{\omega }^1\chi _ϵ(r)\{\begin{array}{ccc}\mathrm{cosh}\widehat{\omega }t\hfill & \mathrm{if}\hfill & ϵ=1\hfill \\ \mathrm{cos}\widehat{\omega }t\hfill & \mathrm{if}\hfill & ϵ=1\hfill \end{array}`$ (63)
$`Z^0`$ $`=`$ $`\widehat{\omega }^1\chi _ϵ(r)\{\begin{array}{ccc}\mathrm{sinh}\widehat{\omega }t\hfill & \mathrm{if}\hfill & ϵ=1\hfill \\ \mathrm{sin}\widehat{\omega }t\hfill & \mathrm{if}\hfill & ϵ=1\hfill \end{array}`$ (66)
$`Z^i`$ $`=`$ $`x^i,`$
where we have denoted $`\chi _ϵ(r)=\sqrt{1ϵ\omega ^2r^2}`$. The line elements
$`ds^2`$ $`=`$ $`\eta _{AB}(ϵ)dZ^AdZ^B`$
$`=`$ $`\chi _ϵ(r)^2dt^2{\displaystyle \frac{dr^2}{\chi _ϵ(r)^2}}r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),`$
are defined on the radial domains $`D_r=[0,1/\sqrt{\widehat{\omega }})`$ or $`D_r=[0,\mathrm{})`$ for dS or AdS respectively.
We calculate the Killing vectors and the orbital generators of the external symmetry in the Cartesian coordinates defined by Eq.(63) and the mentioned parameterization of $`I(M_ϵ)`$ starting with the identification $`\xi ^{(AB)}=\omega ^{AB}`$. Then, from Eqs.(24) and (60), after a few manipulations, we obtain the orbital basis-generators
$`L_{(05)}`$ $`=`$ $`{\displaystyle \frac{iϵ}{\widehat{\omega }}}_t,`$ (68)
$`L_{(j5)}`$ $`=`$ $`{\displaystyle \frac{iϵ}{\widehat{\omega }}}\chi _ϵ(r)\left(\begin{array}{c}\mathrm{cosh}\widehat{\omega }t\hfill \\ \mathrm{cos}\widehat{\omega }t\hfill \end{array}\right)_j+{\displaystyle \frac{ix^j}{\chi _ϵ(r)}}\left(\begin{array}{c}\mathrm{sinh}\widehat{\omega }t\hfill \\ \mathrm{sin}\widehat{\omega }t\hfill \end{array}\right)_t,`$ (73)
$`L_{(0j)}`$ $`=`$ $`{\displaystyle \frac{i}{\widehat{\omega }}}\chi _ϵ(r)\left(\begin{array}{c}\mathrm{sinh}\widehat{\omega }t\hfill \\ \mathrm{sin}\widehat{\omega }t\hfill \end{array}\right)_j+{\displaystyle \frac{ix^j}{\chi _ϵ(r)}}\left(\begin{array}{c}\mathrm{cosh}\widehat{\omega }t\hfill \\ \mathrm{cos}\widehat{\omega }t\hfill \end{array}\right)_t,`$ (78)
$`L_{(ij)}`$ $`=`$ $`i(x^i_jx^j_i).`$ (79)
Furthermore, we consider the Cartesian tetrad gauge defined by Eqs.(52) - (55) where, according to Eq.(5.1), we have
$$u(r)=\chi _ϵ(r)^2,v(r)=w(r)=\chi _ϵ(r).$$
(80)
In addition we take $`\alpha =0`$. In this gauge we obtain the following local $`sl(2,C)`$ generators
$`S_{(05)}(x)`$ $`=`$ $`0,`$ (81)
$`S_{(j5)}(x)`$ $`=`$ $`S_{0j}\left(\begin{array}{c}\mathrm{sinh}\widehat{\omega }t\hfill \\ \mathrm{sin}\widehat{\omega }t\hfill \end{array}\right)+{\displaystyle \frac{1}{r^2}}[\chi _ϵ(r)1][ϵ{\displaystyle \frac{S_{jk}x^k}{\widehat{\omega }}}\left(\begin{array}{c}\mathrm{cosh}\widehat{\omega }t\hfill \\ \mathrm{cos}\widehat{\omega }t\hfill \end{array}\right)`$ (89)
$`{\displaystyle \frac{S_{0k}x^kx^j}{\chi _ϵ(r)}}\left(\begin{array}{c}\mathrm{sinh}\widehat{\omega }t\hfill \\ \mathrm{sin}\widehat{\omega }t\hfill \end{array}\right)],`$
$`S_{(0j)}(x)`$ $`=`$ $`S_{0j}\left(\begin{array}{c}\mathrm{cosh}\widehat{\omega }t\hfill \\ \mathrm{cos}\widehat{\omega }t\hfill \end{array}\right)+{\displaystyle \frac{1}{r^2}}[\chi _ϵ(r)1][{\displaystyle \frac{S_{jk}x^k}{\widehat{\omega }}}\left(\begin{array}{c}\mathrm{sinh}\widehat{\omega }t\hfill \\ \mathrm{sin}\widehat{\omega }t\hfill \end{array}\right)`$ (97)
$`{\displaystyle \frac{S_{0k}x^kx^j}{\chi _ϵ(r)}}\left(\begin{array}{c}\mathrm{cosh}\widehat{\omega }t\hfill \\ \mathrm{cos}\widehat{\omega }t\hfill \end{array}\right)],`$
$`S_{(ij)}(x)`$ $`=`$ $`S_{ij}.`$ (98)
With their help we can write the action of the spin terms (37) and, implicitly, that of the basis-generators $`X_{(AB)}^\rho =S_{(AB)}^\rho +L_{(AB)}`$ of the representations of $`S(M_ϵ)`$ induced by the representations $`\rho `$ of $`SL(2,C)`$. Hereby it is not difficult to show that $`S(M_ϵ)`$ is isomorphic with the universal covering group of $`I(M_ϵ)`$ which in both cases ($`ϵ=\pm 1`$) is a subgroup of the $`SU(2,2)`$ group. As was expected, in the central charts and Cartesian gauge the fields transform manifestly covariant only under the transformations of the subgroup $`SU(2)S(M_ϵ)`$.
### 5.2 Minkowskian charts
Another possibility is to solve the hyperboloid equation (59) in Minkowskian charts where the coordinates, $`x^\mu `$, are defined by
$$Z^5=\widehat{\omega }^1\stackrel{~}{\chi }_ϵ(s),Z^\mu =x^\mu ,$$
(99)
with $`\stackrel{~}{\chi }_ϵ(s)=\sqrt{1+ϵ\widehat{\omega }^2s^2}`$ and $`s^2=\eta _{\mu \nu }x^\mu x^\nu `$. In these coordinates it is convenient to identify the hat indices with the usual ones and to do not rise or lower these indices. Then we find that the metric tensor,
$$g_{\mu \nu }(x)=\eta _{\mu \nu }\frac{ϵ\widehat{\omega }^2}{\stackrel{~}{\chi }_ϵ(s)^2}\eta _{\mu \alpha }x^\alpha \eta _{\nu \beta }x^\beta ,$$
(100)
transforms manifestly covariant under the global $`L_+^{}`$ transformations, $`x^\mu x^\mu =\mathrm{\Lambda }_\nu ^\mu x^\nu `$. Moreover, the whole theory remains manifest covariant if we use the tetrad fields in the Lorentz gauge defined as
$$e_\nu ^\mu (x)=\delta _\nu ^\mu +h_ϵ(s)\eta _{\nu \alpha }x^\alpha x^\mu ,\widehat{e}_\nu ^\mu (x)=\delta _\nu ^\mu +\widehat{h}_ϵ(s)\eta _{\nu \alpha }x^\alpha x^\mu ,$$
(101)
where
$$h_ϵ(s)=\frac{1}{s^2}\left[\stackrel{~}{\chi }_ϵ(s)1\right],\widehat{h}_ϵ(s)=\frac{1}{s^2}\left[\frac{1}{\stackrel{~}{\chi }_ϵ(s)}1\right].$$
(102)
First we calculate the $`SO(4,1)`$ or $`SO(3,2)`$ orbital generators,
$`L_{(\mu 5)}`$ $`=`$ $`{\displaystyle \frac{iϵ}{\widehat{\omega }}}\stackrel{~}{\chi }_ϵ(s)_\mu ,`$ (103)
$`L_{(\mu \nu )}`$ $`=`$ $`i(\eta _{\mu \alpha }x^\alpha _\nu \eta _{\nu \alpha }x^\alpha _\mu ),`$ (104)
which are independent on the gauge fixing. We observe that in the Minkowskian charts $`_t`$ is no more a Killing vector as in the case of the central ones. However, here we have another advantage namely that of the Lorentz gauge in which the local $`sl(2,C)`$ generators of Eq.(37) have the form
$`S_{(\mu 5)}(x)`$ $`=`$ $`{\displaystyle \frac{ϵ}{\widehat{\omega }s^2}}\left[\stackrel{~}{\chi }_ϵ(s)1\right]S_{\mu \alpha }x^\alpha ,`$ (105)
$`S_{(\mu \nu )}(x)`$ $`=`$ $`S_{\mu \nu },`$ (106)
showing that the field $`\psi _\rho `$ transforms manifestly covariant under the whole $`SL(2,C)`$ subgroup of $`S(M_ϵ)`$. Since these representations are induced just by those of $`SL(2,C)`$ we can say that in this gauge the manifest covariance is maximal.
## 6 Concluding remarks
We have presented here the theory of external symmetry in general relativity. Starting with the group $`I(M)`$ which gives the symmetry of the background, we have defined the group $`S(M)`$ showing that its Lie algebra, $`s(M)`$, is isomorphic to $`i(M)`$, having the same structure constants. We have pointed out that the fields with spin transform according to the representations of the $`S(M)`$ group induced by linear representations of $`SL(2,C)`$. This allowed us to calculate the generators of these representations which have specific spin terms. In this way we have obtained the operators associated to the external symmetries which commute with the operator of the relativistic covariant field equation. We have thus the opportunity to choose suitable sets of commuting operators which should determine the quantum modes.
In other respects, since the concrete form of these generators depends on the choice of both the natural and local frames, the commutation rules among their spin and orbital parts are determined by the tetrad gauge. Consequently, the results of the local measurement of the spin observables may depend on the positions of the local frame. This suggests that it should be interesting to investigate new inertial spin effects in other possible tetrad gauge fixings.
However, in our opinion, the most important domain of the further developmens is that of the external symmetry of the quantum field theory in curved spacetimes where the generators of the symmetry transformations must be the one-particle operators corresponding to the external symmetries through the Noether theorem.
Acknowledgments
I would like to thank Mircea Bundaru and Mihai Visinescu for useful comments and enlightening discussions about some sensitive problems appearing here.
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# Interface dynamics from experimental data
## I Introduction
Inverse techniques have a wide range of applicability ranging from geophysics to nonlinear time analysis and statistics . The common philosophy behind these methods is the extraction of equations of motion starting from successive experimental time series of some dynamical variable in addition to basic assumptions such as determinism. If a reasonably general form of the equations is guessed either by symmetry arguments or by general considerations, the “true” parameters are then determined by minimizing a cost function quantifying the distance between experimental observations and corresponding reconstructed quantities, the latter being implicitly dependent upon the parameters. Among such approaches, the Least Square method is the most popular one.
A typical system which can be treated using reconstruction techniques is the case where an observational noise is superimposed to a standard deterministic evolution. In this case the system is expected to evolve under the action of deterministic system and stochasticity come only from our measurement apparatus. The particular case where the dynamics underlying the system is chaotic has also received considerable attention due to its widespread occurence in natural systems , and the importance of treating the presence of the noise with due care has already been emphasized .
The alternative possibility of dynamical noise occurs whenever the noise is a built-in component of the equations of motion. This is a far more difficult problem since one has to deal with stochastic rather than deterministic equations. An important such case, which is rather ubiquitous in nature, is the Langevin dynamics where variables evolve subject both to dissipative generalized forces and to a fluctuating part . In this latter instance, the presence of dynamical noise can drastically modify the dynamics and hence hampers the efficiency of the usual reconstruction techniques based on deterministic ideas.
In our work, we focus on a particular class of Langevin dynamics which has its origin in a seminal paper on interface dynamics but has ever since displayed relations to a variety of physical systems such as for instance bacterial colonial growth, immiscible fluids, directed polymer and superconductors .
The Kardar-Parisi-Zhang (KPZ) equation was introduced as a coarse-grained mesoscopic description for the growth of a rough surface under the deposition of particles driven by gravity. The crucial ingredient introduced in the KPZ and not present in the corresponding linear counterpart, namely the Edward-Wilkinson (EW) equation , is a non-linear term which takes into account the fact that the growth is normal to the surface. The KPZ equation can be mapped into various other models. A Cole-Hopf change of variables maps it into a directed polymer diffusion equation subject to a random potential , while the identification of the local gradient with a velocity leads to the Burgers’s equation for a vorticity-free velocity field . Furthermore it is believed that the KPZ equation has the same large-scale behavior of the Kuramoto-Sivashinsky equation in $`1+1`$ dimensions , while in higher dimensionality the situation is much less clear . Nonetheless, in spite of the gigantic effort devoted to the KPZ equation in the past decade, a complete understanding of its properties is still lacking.
The aim of the present paper is to introduce a new inverse approach to the KPZ equation. A previous attempt due to Lam and Sander was based on the standard Least Square (LS) reconstruction method. These authors used this approach directly on numerically simulated experimental surfaces without a preventive test of the performance of the method itself. Lam and Shin subsequently showed that the standard discretizations used in was not adequate. We shall argue below that, even with the improvements given in , the classical identification procedure devised in Ref. is not properly suited for Langevin dynamics since it is based on deterministic equation ideas. By an explicit computation using the LS technique applied to a $`1+1`$ KPZ equation, we shall review their method and point out what we consider its main deficiencies.
We then go on to introduce a different approach based on the Fokker-Plank equation (FPE) associated to each Langevin equation ,. The advantage of this viewpoint is that one can construct deterministic relations among correlation functions which however still carry informations regarding the fluctuating nature of the original quantities. Those equations can then be easily analyzed within a least squares framework as in the LS method.
The paper is organized as follows. In Sec. II, the KPZ equation is rapidly recalled along with its numerical real space approximations in $`1+1`$ dimensions while the LS approach is reviewed in section III. Sec. IV contains the basic equations of our modified method which is then applied in Sec. V. Numerical results are then given in Sec. VI and some concluding remarks are provided in Sec. VII. More technical points are finally confined in the Appendices. Appendix A shows why the least square method fails for sufficiently large noise amplitudes and Appendix B presents some results concerning renormalized interfaces and their corresponding renormalized equations.
## II Interface dynamics
We consider a one-dimensional line of total length $`L`$ and a surface of height $`h(x,t)`$ at position $`x`$ and time $`t`$. The continuum $`1+1`$ KPZ equation then reads:
$`_th(x,t)`$ $`=`$ $`c+\nu _x^2h(x,t)+{\displaystyle \frac{\lambda }{2}}\left[_xh(x,t)\right]^2+\eta (x,t),`$ (1)
where $`\eta (x,t)`$ is an uncorrelated white noise
$`\eta (x,t)\eta (x^{},t^{})`$ $`=`$ $`2D\delta (xx^{})\delta (tt^{}).`$ (2)
The average $``$ is taken on different realizations of the noise. In Eq.(1) and (2), $`c`$, $`\nu `$, $`\lambda `$ and $`D`$ are coupling parameters ($`c`$ is often set to zero because of the invariance of Eq.(1) under rescaling $`hh+ct`$). For $`\lambda =0`$, Eq.(1) reduces to the Edward-Wilinson (EW) equation which can be solved exactly.
In writing Eq.(1) either a regularization in the correlation given in Eq. (2) (such as for instance a spatially correlated noise ) or the introduction of a minimal length scale $`a`$ is always tacitly assumed. In the latter case, one is then naturally led to consider a discretization of the continuum equation at a given cutoff length scale $`a`$. In that case, (a) the noise term $`\eta (x,t)`$ is discretized
$`\eta _i(t)`$ $`=`$ $`\sqrt{{\displaystyle \frac{D}{a}}}\theta _i(t),`$ (3)
where $`\theta _i(t)`$ is a random noise
$`\theta _i(t)\theta _j(t^{})`$ $`=`$ $`2\delta _{i,j}\delta (tt^{}).`$ (4)
with $`\delta _{i,j}`$ the Kronecker symbol; and (b) Eq. (1) is written for a discrete variable $`h_i(t)`$ ($`i=1,\mathrm{},N=L/a`$) with periodic boundary conditions
$`{\displaystyle \frac{dh_i}{dt}}`$ $`=`$ $`c+\nu _{\mathrm{eff}}F_i^\nu [h]+{\displaystyle \frac{\lambda _{\mathrm{eff}}}{2}}F_i^\lambda [h]+\sqrt{D_{\mathrm{eff}}}\theta _i(t).`$ (5)
Here $`\nu _{\mathrm{eff}}=\nu /a^2`$, $`\lambda _{\mathrm{eff}}=\lambda /a^2`$, $`D_{\mathrm{eff}}=D/a`$. $`F_i^\nu [h]`$ and $`F_i^\lambda [h]`$ are proper discretizations of the linear $`_x^2h`$ and non-linear $`(_xh)^2`$ terms respectively. We note that the exact meaning of “proper discretization” has been the object of some investigations .
In all practical applications, a further temporal discretization , is also performed on Eq. (5)
$`h_i(t+\delta t)`$ $`=`$ $`h_i(t)+\delta t\left(c+\nu _{\mathrm{eff}}F_i^\nu [h(t)]+{\displaystyle \frac{\lambda _{\mathrm{eff}}}{2}}F_i^\lambda [h(t)]\right)+\sqrt{2D_{\mathrm{eff}}\delta t}r_i,`$ (6)
where $`r_i`$ is a Gaussian random generator of unit variance and $`\delta t`$ the discretization time step.
In $`d=1+1`$ it is known that the steady state solution $`P[h]`$ for the probability distribution of the heights in the KPZ equation is identical to the EW stationary distribution due to a fluctuation-dissipation theorem . It was shown , that the correct stationary discrete probability namely
$`P[h]`$ $`=`$ $`𝒩^1\mathrm{exp}[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\nu }{Da}}{\displaystyle \underset{i=1}{\overset{N}{}}}(h_ih_{i+1})^2],`$ (7)
where $`𝒩^1`$ is a normalization factor, can be obtained by taking
$`F_i^\nu [h]`$ $`=`$ $`h_{i+1}+h_{i1}2h_i,`$ (8)
and
$`F_i^\lambda [h]`$ $`=`$ $`{\displaystyle \frac{1}{3}}[(h_{i+1}h_i)^2+(h_{i+1}h_i)(h_ih_{i1})+(h_ih_{i1})^2].`$ (9)
The standard choice $`F_i^\lambda [h]=1/4(h_{i+1}h_{i1})^2`$ on the other hand fails to reproduce Eq. (7) and suffers of other problems as well . A necessary (albeit not sufficient) condition for the identification with the continuum counterpart Eq. (1), is clearly that the correct steady state (i.e. independent of $`\lambda `$) is recovered. For this reason, we shall exploit for the new identification procedure as well as for the LS scheme, equations (8) and (9) hereafter instead of the standard choice which was used in .
## III The Least Squares error model method
Before introducing our method we first review the LS error method used in Ref.. We consider experimental surfaces coarse-grained at length scale $`a`$ described by the interface heights $`h_i^{\mathrm{obs}}(t)`$, ($`i=1,\mathrm{},N`$) which are sampled $`M`$ times i.e. at discrete times $`t=t_k=k\mathrm{\Delta }t`$ ($`k=1,\mathrm{},M+1`$). Note that the sampling time $`\mathrm{\Delta }t`$ is the time interval between two experimental observations and it is clearly different from the discretization time $`\delta t`$ of Eq.(6). For surfaces obtained by numerical simulations $`\mathrm{\Delta }t`$ is typically a multiple of $`\delta t`$. We note that in Ref., the authors used $`\mathrm{\Delta }t`$ equal to $`\delta t`$ which is a rather particular case.
For the sake of simplicity, we assume here that measurements are free from observational noise. It must be emphasized that, in the presence of measurement noise, our method performs a priori better then the LS scheme since it is based on spatial-averaged values which are less affected by errors on local height measurements.
Our purpose is to determine the coefficients $`c`$, $`\nu `$, $`\lambda `$, $`D`$ at the given length scale $`a`$ in Eq.(3).
Let us first neglect the dynamical noise in Eq.(5). We then obtain a standard identification problem of the coupling parameters governing a deterministic non-linear equation which can be cast in the compact form:
$`{\displaystyle \frac{dh_i}{dt}}`$ $`=`$ $`{\displaystyle \underset{\alpha =1}{\overset{p}{}}}\mu _\alpha F_i^\alpha [h],`$ (10)
where in the present case $`p=3`$ and $`\mu _1=c`$, $`\mu _2=\nu _{\mathrm{eff}}`$, $`\mu _3=\lambda _{\mathrm{eff}}`$, (8) and (9) are used for $`F_i^2[h]`$ and $`F_i^3[h]`$ whereas $`F_i^1[h]=1`$.
Optimal parameters are then determined by minimizing a cost function $`𝒥`$ such as the sum-square difference
$`𝒥`$ $`=`$ $`{\displaystyle \frac{1}{NM}}{\displaystyle \underset{k=1}{\overset{M}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left[h_i^{\mathrm{obs}}(t_{k+1})h_i^{\mathrm{pred}}(t_{k+1})\right]^2,`$ (11)
which quantifies the distance between experimental observations $`h_i^{\mathrm{obs}}(t_k)`$ and equivalent reconstructed quantities $`h_i^{\mathrm{pred}}(t_k)`$. The latter quantities are computed from Eq. (10) for given parameters and are thus generally implicit functions of the parameters. However, if the sampling time $`\mathrm{\Delta }t`$ is small enough, then $`F_i^\alpha [h]`$ are nearly constant between two measurements and the amplitudes $`h_i^{\mathrm{pred}}(t_{k+1})`$ can be related to the parameters $`\mu _\alpha `$ by
$`h_i^{pred}(t_{k+1})`$ $`=`$ $`h_i^{obs}(t_k)+\mathrm{\Delta }t{\displaystyle \underset{\alpha =1}{\overset{p}{}}}\mu _\alpha F_i^\alpha [h^{\mathrm{obs}}(t_k)]`$ (12)
In this case, the cost function $`𝒥(\{\mu \})`$ itself becomes explicit and quadratic in the parameters. Optimal parameters can thus be evaluated through a simple matrix inversion. Indeed the extremal value of $`𝒥`$ is inferred by
$`{\displaystyle \frac{𝒥}{\mu _\alpha }}|_\mu ^{}`$ $`=`$ $`0.`$ (13)
The solution for the optimal parameters $`\{\mu ^{}\}`$ is then given by a matrix equation:
$`\mu _\alpha ^{}`$ $`=`$ $`{\displaystyle \underset{\beta =1}{\overset{p}{}}}A_{\alpha \beta }^1B_\beta ,`$ (14)
where we have defined
$`A_{\alpha \beta }`$ $`=`$ $`{\displaystyle \frac{1}{NM}}{\displaystyle \underset{k=1}{\overset{M}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}F_i^\alpha F_i^\beta ,`$ (15)
$`B_\alpha `$ $`=`$ $`{\displaystyle \frac{1}{NM}}{\displaystyle \underset{k=1}{\overset{M}{}}}{\displaystyle \underset{i=1}{\overset{N}{}}}\left[{\displaystyle \frac{h_i^{obs}(t_{k+1})h_i^{obs}(t_k)}{\mathrm{\Delta }t}}\right]F_i^\alpha ,`$ (16)
and where functions $`F_i^\alpha `$ are clearly expressed at $`h_1^{\mathrm{obs}}(t_k),\mathrm{},h_N^{\mathrm{obs}}(t_k)`$.
This classical least squares method is an easy and natural approach and it works fairly well in the absence of any noise. In the presence of noise however, it has its main drawback in the fact that it approximates time derivatives by finite differences. If the dynamics is governed by a deterministic equation and measurements are performed with a negligible observational noise, this simply imposes the choice of a sampling time much smaller than the characteristic or relaxation time of the process.
Lam and Sander assumed that if the sampling time $`\mathrm{\Delta }t`$ is small enough, the above method could be extended to a Langevin equation (i.e. with dynamical noise). The amplitude of the noise can then be inferred from Eq.(11) when $`J`$ is taken at the minimum values of the parameters, that is
$`D`$ $`=`$ $`{\displaystyle \frac{1}{2\mathrm{\Delta }t}}aJ(\{\mu ^{}\}))`$ (17)
However it was already observed in dynamical systems that even with pure measurement noise the above method can cause large errors. This is expected to be the case for dynamical noise as well. Two main reasons for this could be advocated. First if $`\mathrm{\Delta }t`$ is too large, the linear approximation (12) which explicitly relates the observed quantities breaks down. Because of the dynamical noise term, this happens a priori for shorter times intervals in a Langevin equation compared with its deterministic counterpart. Second, even in the favorable case in which $`\mathrm{\Delta }t`$ is small , such a method is efficient only if large sizes and small noise amplitudes are used. This is explained in Appendix A where a simple zero-dimensional case is explicitly worked out with the method of Lam and Sander.
## IV Stochastic approach for Model identification.
We now turn to our method which is based on the simple observation that all the information present in the Langevin equation (5) are also contained in the corresponding Fokker-Planck equation (FPE) :
$`_tP[h,t]`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}D_{\mathrm{eff}}{\displaystyle \frac{^2}{h_i^2}}P[h,t]{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \frac{}{h_i}}(F_i[h]P[h,t]),`$ (18)
where
$`F_i[h]`$ $`=`$ $`c+\nu _{\mathrm{eff}}F_i^\nu [h]+{\displaystyle \frac{1}{2}}\lambda _{\mathrm{eff}}F_i^\lambda [h],`$ (19)
and
$`P[h,t]`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}\delta (h_ih_i(t)),`$ (20)
where the solution $`h_i(t)`$ is associated to a particular noise configuration $`\theta _i(t)`$.
In equation (18) the second term on the r.h.s. characterizes the deterministic behavior of the system whereas the first term contains stochastics effects. We derive a first general equation involving the parameters $`c`$ and $`\lambda `$. Using Eq. (18), the time derivative of the ensemble average of $`h_i(t)`$ can be easily shown to be:
$`{\displaystyle \frac{dh_i(t)}{dt}}`$ $`=`$ $`{\displaystyle 𝒟hF_i[h]P[h,t]},`$ (21)
where $`𝒟h_{i=1}^Ndh_i`$. If we denote by $`g^{(1)}(t)=\frac{1}{N}_ih_i(t)`$ the mean height at time $`t`$ averaged over the noise, its time derivative can be written after some simple algebra as
$`{\displaystyle \frac{dg^{(1)}(t)}{dt}}`$ $`=`$ $`c+{\displaystyle \frac{\lambda _{\mathrm{eff}}}{6}}\left[2g_0^{(2)}(t)+g_1^{(2)}(t)\right],`$ (22)
where we have defined
$`g_l^2(t)`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta h_i(t)\delta h_{i+l}(t),`$ (23)
in the variables $`\delta h_i=h_ih_{i+1}`$. Note that there are only $`N1`$ independent $`\delta h_i`$ variables due to periodic boundary conditions and to the fact that $`_{i=1}^N\delta h_i=0`$.
The above result prompts a convenient change of variables from $`h_1,..,h_N`$ to $`\delta h_1,..,\delta h_{N1},\overline{h}1/N_{i=1}^Nh_i`$ followed by an integration over $`\overline{h}`$. Physically this is related to the fact that our system is infinitely degenerate with respect to the average height. Note that the stationary probability Eq. (7) is now Gaussian and well defined in the new variables $`\delta h_1,\mathrm{},\delta h_{N1}`$. The corresponding probability $`\stackrel{~}{P}[\delta h]`$ is the solution of a modified FPE:
$`_t\stackrel{~}{P}[\delta h,t]`$ $`=`$ $`2D_{\mathrm{eff}}{\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \frac{^2}{\delta h_i^2}}\stackrel{~}{P}[\delta h,t]`$ (24)
$``$ $`{\displaystyle \underset{i=1}{\overset{N1}{}}}{\displaystyle \frac{}{\delta h_i}}\left(G_i[\delta h]\stackrel{~}{P}[\delta h,t]\right)2D_{\mathrm{eff}}{\displaystyle \underset{i=2}{\overset{N1}{}}}{\displaystyle \frac{^2}{\delta h_i\delta h_{i1}}}\stackrel{~}{P}[\delta h,t],`$ (25)
where we have defined
$`G_i`$ $`=`$ $`F_iF_{i+1}=\nu _{\mathrm{eff}}G_i^\nu +{\displaystyle \frac{1}{2}}\lambda _{\mathrm{eff}}G_i^\lambda ,`$ (26)
with
$`G_i^\nu `$ $`=`$ $`\delta h_{i+1}+\delta h_{i1}2\delta h_i,`$ (27)
and
$`G_i^\lambda `$ $`=`$ $`{\displaystyle \frac{1}{3}}\left[\delta h_{i1}^2\delta h_{i+1}^2\delta h_i(\delta h_{i+1}\delta h_{i1})\right].`$ (28)
We are now in the position to derive our second basic result.
Integrating Eq. (24) over all variables but (say) $`\delta h_j`$, one gets, for the single variable probability
$`p(\delta h_j)`$ $`=`$ $`{\displaystyle 𝒟\delta h_j\stackrel{~}{P}[\delta h,t]},`$ (29)
where the shorthand notation $`𝒟\delta h_j=_{ij=1}^{N1}d\delta h_i`$ was again exploited, the following equation:
$`_tp(\delta h_j,t)`$ $`=`$ $`2D_{\mathrm{eff}}{\displaystyle \frac{^2}{\delta h_j^2}}p(\delta h_j,t){\displaystyle \frac{}{\delta h_j}}\pi (\delta h_j,t),`$ (30)
where the non-local term $`\pi (\delta h_j,t)`$ is defined as
$`\pi (\delta h_j,t)`$ $`=`$ $`{\displaystyle 𝒟\delta h_jG_j[\delta h]\stackrel{~}{P}[\delta h,t]}.`$ (31)
The last step is to introduce the Fourier transform of $`p(\delta h_j,t)`$ which can be reckoned as a generating function for all moments of the distribution. Specifically on defining
$`\widehat{p_j}(q,t)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑\delta h_je^{\mathrm{i}q\delta h_j}p(\delta h_j,t),`$ (32)
we find a simple equation for the average $`\widehat{p}(q,t)`$ over all sites of $`\widehat{p_j}(q,t)`$ :
$`_t\widehat{p}(q,t)`$ $`=`$ $`2D_{\mathrm{eff}}q^2\widehat{p}(q,t)\mathrm{i}q\widehat{\pi }(q,t),`$ (33)
in which $`\widehat{\pi }(q,t)`$ is the Fourier transform of $`\pi (\delta h_j,t)`$ averaged over all sites. One can then expand Eq. (33) in powers of $`q`$ and obtain an infinite hierarchy (closure problem) in the correlation functions. The first two non-trivial orders ($`O(q^2)`$ and $`O(q^3)`$) are
$`{\displaystyle \frac{dg_0^{(2)}(t)}{dt}}`$ $`=`$ $`4\nu _{\mathrm{eff}}[g_1^{(2)}(t)g_0^{(2)}(t)]+4D_{\mathrm{eff}},`$ (34)
and
$`{\displaystyle \frac{dg_{00}^{(3)}(t)}{dt}}`$ $`=`$ $`3\nu _{\mathrm{eff}}[g_{11}^{(3)}(t)+g_{01}^{(3)}(t)2g_{00}^{(3)}(t)]+{\displaystyle \frac{1}{2}}\lambda _{\mathrm{eff}}[g_{001}^{(4)}(t)g_{111}^{(4)}(t)],`$ (35)
where we have defined the following higher order correlation functions
$`g_{lm}^{(3)}(t)`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta h_i(t)\delta h_{i+l}(t)\delta h_{i+m}(t),`$ (36)
$`g_{lmn}^{(4)}(t)`$ $`=`$ $`{\displaystyle \frac{1}{N}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta h_i(t)\delta h_{i+l}(t)\delta h_{i+m}(t)\delta h_{i+n}(t).`$ (37)
It is worth mentioning that $`\lambda `$ does not explicitly appear in equation (34). As one can explicitly check, this is a feature associated to the particular discretization Eq. (9) and it would have not been the case had we used the standard discretization for $`F_i^\lambda [h]`$. This is clearly in turn related to the fact that the steady state probability distribution Eq. (7) is independent of $`\lambda `$. We also note that in the $`1+1`$ case we are considering, the explicit steady state solution of Eq.(33) is known, and depends only on a single parameter $`\frac{D}{\nu }`$. As a consequence, the steady version of Eq.(33) cannot be used here to identify $`\nu `$ and $`D`$. In the $`2+1`$ case where such a peculiar feature is not present, the stationary solution depends on $`\lambda `$ as well and parameters identification can exploit the steady analogue equation.
## V Parameters identification
Our aim is to implement an identification procedure which could be exploited on real experiments. For this reason, we assume that the experimental surface is constituted by a finite number of sites $`N`$ with lattice spacing $`a`$ (corresponding to a size $`L=Na`$), and it is measured during a finite time $`T_{obs}`$ every sampling time $`\mathrm{\Delta }t`$. Again $`\mathrm{\Delta }t`$ is a priori different from the discretization time $`\delta t`$ when the data is produced numerically (note that in real experiments $`\delta t`$ is not even defined!). We shall test the robustness and efficiency of the scheme with respect to the size $`L`$ and sampling time $`\mathrm{\Delta }t`$.
Identification methods are often based on minimizing a cost function defined through dynamical constraints. This is clearly the case of the least square method as explained in section III. Here we derive dynamical constraints using Eqns. (22) and (34) which contain informations of the original Langevin equation including mean values and fluctuations around mean values. The present identification is thus based on deterministic equations. This constitutes a crucial difference with respect to the previous reconstruction method directly based on stochastic equations. Another important feature is that the observed quantities we use in our reconstruction scheme are dealing with averaged site values. Hence the fluctuations of all these terms, which derive from stochastic quantities, are reduced typically by a factor $`1/\sqrt{N}`$, and self-averaging is expected to be more effective.
Let us derive the constraints we use. First, the total observation time $`T_{obs}`$ is divided into $`q`$ equal slices $`[T_1,T_2]`$,$`\mathrm{}`$,$`[T_q,T_{q+1}]`$ with $`\mathrm{\Delta }T=T_{j+1}T_j`$. Let us integrate (22) and (34) on each slice $`[T_j,T_{j+1}]`$:
$`{\displaystyle \frac{\mathrm{\Delta }g^{(1)}}{T_{j+1}T_j}}`$ $`=`$ $`c+{\displaystyle \frac{\lambda _{\mathrm{eff}}}{6}}{\displaystyle \frac{1}{T_{j+1}T_j}}{\displaystyle _{T_j}^{T_{j+1}}}𝑑t\left[2g_0^{(2)}(t)+g_1^{(2)}(t)\right],`$ (38)
$`{\displaystyle \frac{\mathrm{\Delta }g_0^{(2)}}{T_{j+1}T_j}}`$ $`=`$ $`4\nu _{\mathrm{eff}}{\displaystyle \frac{1}{T_{j+1}T_j}}{\displaystyle _{T_j}^{T_{j+1}}}𝑑t\left[g_1^{(2)}(t)g_0^{(2)}(t)\right]+4D_{\mathrm{eff}}.`$ (39)
If functions and integrals in Eqns. (38) and (39) are computed using experimental data, these discrete equations provide $`2q`$ relations between the parameters to identify. From these constraints, two cost functions are built in a way already described in Sec.III with $`p=2`$. The corresponding $`2\times 2`$ equations then yield $`c`$ and $`\lambda `$ from one cost function and $`\nu `$ and $`D`$ from the other.
We now explain how functions and integrals in Eqns. (38) and (39) are obtained experimentally. Starting with a same initial surface e.g. a flat surface, we will grow the surface $``$ times. Because of the stochastic nature of the phenomenon, this produces $``$ different observations or realizations of the same process . Such a procedure, which can be performed very easily in real experiments, allows the computation, at sampling times $`t=t_k=k\mathrm{\Delta }t`$ ($`k=1,\mathrm{},M+1`$), of $`[g^{(1)}]_{\mathrm{exp}}`$, $`[g_0^{(2)}]_{\mathrm{exp}}`$ and $`[g_1^{(2)}]_{\mathrm{exp}}`$. Indeed, these quantities are the averages over $``$ different realizations of the spatial average height and correlations of the first neighbors. The number $``$ of realizations need not be large: if the total number $`N`$ of sites is sufficiently large, the experimental values are rather close to the corresponding theoretical predictions $`g^{(1)}(t)`$, $`g_0^{(2)}(t)`$ and $`g_1^{(2)}(t)`$. From these functions sampled every $`t=t_k=k\mathrm{\Delta }t`$, integrals in Eqns. (38) and (39) can be efficiently evaluated for small sampling time $`\mathrm{\Delta }t`$. In this case, the smooth functions $`[g^{(1)}]_{\mathrm{exp}}`$, $`[g_0^{(2)}]_{\mathrm{exp}}`$ $`[g_1^{(2)}]_{\mathrm{exp}}`$ can be approximated on the whole time interval $`[T_1,T_{q+1}]`$ by a standard curve fitting algorithm which gives as a by product the time integrals. This method does impose a constraint on the sampling time $`\mathrm{\Delta }t`$. However, this constraint is substantially weaker with respect to that imposed by the LS method as we will show below. This is a considerable advantage of our new procedure.
Two remarks are here in order. Firstly one expects the result to be independent on the number of slices $`q`$ provided that $`q`$ satisfies the following two constraints. On the one hand $`q`$ should be greater than $`2`$ (since two parameters are identified per cost functions) and on the other hand it should be less than $`M=\frac{T_{obs}}{\mathrm{\Delta }t}`$ so that $`\mathrm{\Delta }T`$ cannot be less than the sampling time $`\mathrm{\Delta }t`$. Secondly the identification of $`\nu `$ and $`D`$ could be achieved by using Eq.(33) rather than Eq. (34). We shall see that in our approach the two equations yield virtually identical results.
## VI Results
In order to test the potentiality of the different identification methods, we produce experimental data by simulating Eq.(5) with a standard Euler time integration algorithm with time step $`\delta t=0.01`$, lattice spacing $`a=1`$, and parameters $`\nu =D=1`$ and $`\lambda =3`$. These are the same values used in Ref.. The time step is expected to be sufficiently small for not causing instability problems and the non-linear term $`\lambda `$ is big enough to be well inside the KPZ phase. We find interesting to repeat each calculation few times (typically 5) to give an estimate of the error bars to be associated to each parameter value (this was missing in previous works).
### A LS Method
Let us compute the parameters using the original LS method with the spatial and temporal discretization Eqns. (5) and (6). We exploit the same trick used in Ref. in which a KPZ surface of size $`2L`$ is obtained by a magnification of a fully relaxed surface of size $`L`$ where the height are rescaled by a factor $`2^\alpha `$, ($`\alpha =0.5`$) and linearly interpolated. The obtained surface is then relaxed to stationarity before a successive magnification is attempted. However, unlike Ref. where a single surface of size $`L=32768`$ was computed, we consider $`L=512,1024,2048,4096`$ and linearly extrapolate the results to the limit $`L\mathrm{}`$. The calculation is repeated for increasing values of $`s=\mathrm{\Delta }t/\delta t`$ in order to display the crucial weakness of the method as explained before. Fig. 1 depicts the results for the parameter $`\nu `$ at finite $`L`$. Similar trends are present for $`\lambda `$ and $`D`$. The extrapolated values at $`L\mathrm{}`$ are reported in Table I. The gradual decrease in the precision of the reconstructed parameters is apparent and it shows the loss of accuracy of the LS method as $`\mathrm{\Delta }t`$ increases as previously advertised.
We also considered the LS when the reconstructed quantities are computed at time intervals $`\mathrm{\Delta }T`$ which are multiple of the sampling time $`\mathrm{\Delta }t`$. In fact this test was also carried out by the authors of Ref. (in their notations $`\tau =\mathrm{\Delta }T`$ and $`\mathrm{\Delta }t=\delta t`$) and it will constitute a further source of comparison with our alternative stochastic method (see below). Even in this case there is a decrease in the performance of the procedure as the ratio $`r=\mathrm{\Delta }T/\mathrm{\Delta }t`$ increases, consistently with the results of Ref. . The corresponding extrapolated values are reported in Table II.
### B Stochastic Approach
For a more convenient comparison with the LS method, we use the same sizes and statistics (five different configurations for each size). Our calculations are carried out in the transient rather than in the steady state and are therefore much less time consuming. Again the results are obtained for $`L=512,1024,2048,4096`$ and linearly extrapolated to $`L\mathrm{}`$. For a comparison with the previous calculation, the outcomes of the parameter $`\nu `$ at different size $`L`$ are plotted in Fig.2 for increasing values of the ratio $`s=\mathrm{\Delta }t/\delta t`$, and the corresponding extrapolated values are reported in Table III. One can see that the parameter values are rather insensitive to the changing the ratio $`s=\mathrm{\Delta }t/\delta t`$ as expected. Next we check the performance of our method with respect to the increasing of the ratio $`r=\mathrm{\Delta }T/\mathrm{\Delta }t`$. This is reported in Table IV. As expected, our method outperforms the LS one in all situations.
Since the LS method could in principle be carried out in transient rather than in steady state conditions, one might wonder how it would perform in this case.
To this aim we recompute the parameters using the LS method under these conditions and find that the predicted values are far off with respect to the exact ones. For instance for $`L=4096`$ a typical run yields $`\nu 0.36`$, $`\lambda 0.68`$ and $`D0.005`$ to be compared with a typical result obtained with our method $`\nu 0.99`$, $`\lambda 2.98`$ and $`D1.01`$.
As a final cross-check of our method, we recompute the parameters in the same situation as before but using Eq. (33) rather than Eq. (34) to extract $`\nu `$ and $`D`$ and find nearly identical values.
### C Coarse-graining and KPZ real discretization.
The application of this machinery to experimental surfaces assumes that the system is described by a KPZ-like dynamics. In this case, besides being able to address the issue of whether or not they belong to the KPZ universality class, one would be able to provide a numerical estimates of the coupling parameters which are usually overlooked in studies focussing only on the universality class.
Following Lam and Sander , we produce an interface based on the KPZ discretized model Eq. (5) which is then smoothed by introducing the (discrete) Fourier transform of the heights
$`\widehat{h}_{q_n}(t)`$ $`=`$ $`a{\displaystyle \underset{i=1}{\overset{N}{}}}e^{\mathrm{i}q_nx_i}h_i(t)`$ (40)
A coarse-graining surface at level $`a_s=ba`$ can then be achieved by simplying setting to zero all wavelength components $`\widehat{h}_{q_n}(t)`$ with $`qN_s=L/a_s`$ and transforming back to real space. The obtained smoothed surface is then assumed to be governed by a KPZ equation (renormalizability property). This new KPZ dynamics can then reconstructed along lines similar to those described above before a further time step is carried out on the original surface.
With $`b=2`$, $`\mathrm{\Delta }T/\mathrm{\Delta }t=1`$ and sizes up to $`L=8192`$ averaged over 5 configurations as before, we find $`\nu =1.09\pm 0.04`$, $`\lambda =3.27\pm 0.05`$, and $`D=0.88\pm 0.03`$. Higher values of $`b`$ result in poorer and poorer agreement with the expected values even with higher lattice sizes. The same feature is also present in the original LS procedure as we explicitly checked. In fact this is a general deficiency of the real space discretization as explained in Appendix B: the finite size difference has lost some renormalizability property of the original KPZ continuum equation.
## VII Conclusions
In this paper, we discuss a method for extracting the coupling parameters from a non-linear Langevin equation starting from experimental surfaces representing successive snapshots of the system. We apply this scheme to the KPZ equation in $`1+1`$ dimensions (although it could be extended to any dimensions) and compare it with the previous approach of Ref. , finding the following differences. First of all it does not require large sizes and it is well suited for a transient state. This is expected to be a considerable advantage, notably in numerical work, since the typical time required to reach a steady state increases as $`L^z`$ where $`z`$ is the dynamical exponent ($`3/2`$ in the $`1+1`$ KPZ case). We have explicitly shown how the LS method which works rather well in the aforementioned conditions, fails to provide sensible answers otherwise. Most importantly however is the fact that our approach is stable under the changing of the sampling time, unlike the LS method which is not. We stress the importance of this feature since in typical experimental situations, the sampling time is an externally tuned parameters which has nothing to do with the evolution time of the system. We have discussed the reasons why this is so and provide an intuitive heuristic argument showing why the LS scheme is not expected to work under these more realistic conditions. Finally we implemented a coarse-graining procedure in order to be able to apply our method to experimentally generated profiles. We showed that the agreement with the expected values is much poorer in the present case and we further argued that any real space based approach is doomed to run into this problem, the reason being that they have not a correct renormalization behavior under coarse-graining as explained in Appendix B. We have recently devised an alternative approach based on a Fourier-based scheme which avoids discretization problems. The results of this will be the subject of a forthcoming publication.
###### Acknowledgements.
This work was supported by a joint CNR-CNRS exchange program number 5274. One of us (AG) acknowledges financial support by MURST and INFM.
## A A simple solvable example
This appendix shows, on a simple and solvable example which is a zero-dimensional analogue of Eq.(1), that the method of least squares can be hampered by the presence of dynamical noise. Let us assume that the scalar variable $`X(t)`$ is governed by the following Langevin equation:
$`{\displaystyle \frac{dX}{dt}}`$ $`=`$ $`B(X)+\mu G(X)+\eta (t),`$ (A1)
where $`B`$ and $`G`$ are prescribed functions, $`\eta (t)`$ is an uncorrelated white noise
$`\eta (t)\eta (t^{})`$ $`=`$ $`2D\delta (tt^{}).`$ (A2)
For simplicity, we assume that: (a) measurement noise is negligible, (b) the observed time serie $`X^{\mathrm{obs}}(t_k)`$ with $`t_k=k\mathrm{\Delta }t`$ ($`k=1,\mathrm{},M+1`$) has been produced by the dynamical system (A1) with the value $`\mu =0`$. We ask whether the least square method is capable of identifying the correct coupling parameter $`\mu =0`$.
From the one hand, the least square method first assumes that the data is produced by the deterministic counterpart of Eq.(A1) with an unkown parameter $`\mu `$. In discrete times, this yields a ”predicted” value given by
$`X^{\mathrm{pred}}(t_{k+1})`$ $`=`$ $`X^{\mathrm{obs}}(t_k)+\mathrm{\Delta }t\left[B(X^{\mathrm{obs}}(t_k))+\mu G(X^{obs}(t_k))\right].`$ (A3)
On the other hand, the ”observed” value is given, if the sampling time is small enough, by the discrete time counterpart of Eq.(A1) with $`\mu =0`$ :
$`X^{\mathrm{obs}}(t_{k+1})`$ $`=`$ $`X^{\mathrm{obs}}(t_k)+\mathrm{\Delta }tB(X^{obs}(t_k))+r(t_k)\sqrt{2D\mathrm{\Delta }t},`$ (A4)
where $`r(t_k)`$ is a gaussian random generator of unit variance.
Using both experimental observations $`X^{\mathrm{obs}}(t_k)`$ and reconstructed quantities $`X^{\mathrm{pred}}(t_k)`$, a cost function can be constructed:
$`𝒥`$ $`=`$ $`{\displaystyle \frac{1}{M}}{\displaystyle \underset{k=1}{\overset{M}{}}}\left[X^{\mathrm{obs}}(t_{k+1})X^{\mathrm{pred}}(t_{k+1})\right]^2.`$ (A5)
Using Eqns. (A3) and (A4), the cost function is readly rewritten as
$`𝒥`$ $`=`$ $`{\displaystyle \frac{1}{M}}{\displaystyle \underset{k=1}{\overset{M}{}}}\left[\mathrm{\Delta }t\mu G(X^{\mathrm{obs}}(t_k))r(t_k)\sqrt{2D\mathrm{\Delta }t}\right]^2.`$ (A6)
The minimum value of $`𝒥`$ then satisfies the extremality condition
$`{\displaystyle \frac{𝒥}{\mu }}|_\mu ^{}`$ $`=`$ $`0,`$ (A7)
which provides the value
$`\mu ^{}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2D}{\mathrm{\Delta }t}}}{\displaystyle \frac{\frac{1}{M}_{k=1}^Mr(t_k)G(X^{\mathrm{obs}}(t_k))}{\frac{1}{M}_{k=1}^MG^2(X^{\mathrm{obs}}(t_k))}}.`$ (A8)
Let us estimate the two sums appearing in Eq.(A8). Because of self-averaging, the denominator can be clearly rewritten as
$`{\displaystyle \frac{1}{M}}{\displaystyle \underset{k=1}{\overset{M}{}}}G^2(X^{\mathrm{obs}}(t_k))G^2,`$ (A9)
the average being over the noise $`\eta `$. The second sum in can be separated into two contributions
$`{\displaystyle \frac{1}{M}}{\displaystyle \underset{k=1}{\overset{M}{}}}[r(t_k)G(X^{\mathrm{obs}}(t_k)]`$ $`=`$ $`{\displaystyle \frac{1}{M}}{\displaystyle \underset{k=1}{\overset{M}{}}}r(t_k)[G(X^{obs}(t_k)G]+{\displaystyle \frac{G}{M}}{\displaystyle \underset{k=1}{\overset{M}{}}},r(t_k)`$ (A10)
where using the Central Limit Theorem (CLT) we have that
$`{\displaystyle \frac{G}{M}}{\displaystyle \underset{k=1}{\overset{M}{}}}r(t_k)`$ $``$ $`{\displaystyle \frac{G}{\sqrt{M}}}r_1,`$ (A11)
where $`r_1`$ is a random variable of unit variance. In Eq.(A10), each term of the first sum of the r.h.s. is a random variable which is a product of two independent random variables of zero mean and variance respectively equal to $`1`$ and $`\sigma =G^2G^2`$. A similar argument based on the CLT applies to the evaluation of the first term of the r.h.s. since it is again a sum of $`M`$ random variables with zero average which implies that:
$`{\displaystyle \frac{1}{M}}{\displaystyle \underset{k=1}{\overset{M}{}}}r(t_k)[G(X^{\mathrm{obs}}(t_k)G]`$ $``$ $`{\displaystyle \frac{\sqrt{\sigma }}{\sqrt{M}}}r_2,`$ (A12)
where $`r_2`$ is a random variable of unit variance. Instead of the true value one then finds
$`\mu ^{}`$ $``$ $`{\displaystyle \frac{\sqrt{2D}}{\sqrt{M\mathrm{\Delta }t}}}{\displaystyle \frac{Max(\sqrt{\sigma },G)}{G^2}}.`$ (A13)
With non-vanishing amplitudes $`D`$, the noise hence causes errors for small $`\mathrm{\Delta }t`$ unless large statistics are considered.
## B Non-renormalizability of real space discretizations
Let us consider a one-dimensional surface defined on a lattice of length $`L=Na`$ ($`a`$ being the lattice spacing and $`N`$ being the total number of sites) which is identified by $`h_1,\mathrm{},h_N`$ at positions $`x_1=a,\mathrm{},x_N=Na`$ and having periodic boundary conditions. If we assume a KPZ dynamics, then the equation of motion is given by Eq.(5) and its stationary state by Eq. (7). Let us introduce the (discrete) Fourier transform $`\widehat{h}_q`$ so that
$`h_i(t)`$ $`=`$ $`{\displaystyle \frac{1}{L}}{\displaystyle \underset{n=N/2}{\overset{N/2}{}}}e^{\mathrm{i}q_nx_i}\widehat{h}_{q_n}(t),`$ (B1)
and conversely
$`\widehat{h}_{q_n}(t)`$ $`=`$ $`a{\displaystyle \underset{i=1}{\overset{N}{}}}e^{\mathrm{i}q_nx_i}.h_i(t)`$ (B2)
By using the relation
$`{\displaystyle \underset{i=1}{\overset{N}{}}}e^{\mathrm{i}\left(q_nq_m\right)x_i}`$ $`=`$ $`N\delta _{n,m},`$ (B3)
it is easy to show that the correct corresponding of Eq.(7) for the stationary distribution is
$`\widehat{P_a}[\widehat{h}]`$ $`=`$ $`𝒩^1\mathrm{exp}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\nu }{DL}}{\displaystyle \underset{n=N/2}{\overset{N/2}{}}}q_n^2|\widehat{h}_{q_n}|^2\right].`$ (B4)
In Eq.(B4) the continuum limit $`a0`$ is simply achieved by letting $`N\mathrm{}`$.
We now recall that the proper variables to be used in this context are the $`\delta h_i=h_ih_{i+1}`$ whose Fourier transform $`\delta \widehat{h}_{q_n}`$ are related to the Fourier transform $`\widehat{h}_{q_n}`$ of the heights by
$`\delta \widehat{h}_{q_n}`$ $`=`$ $`\widehat{h}_{q_n}\left(1e^{\mathrm{i}q_na}\right).`$ (B5)
We can exploit the periodic boundary conditions to extend the $`N1`$ sites to $`N`$ (bearing however in mind that only $`N1`$ are independent) and use the the fact that
$`{\displaystyle \frac{1}{a}}{\displaystyle \underset{i=1}{\overset{N}{}}}\delta h_i^2`$ $`=`$ $`{\displaystyle \frac{1}{La^2}}{\displaystyle \underset{n=N/2}{\overset{N/2}{}}}|\delta \widehat{h}_{q_n}|^2,`$ (B6)
to rewrite the probability (B4) in the alternative form
$`\widehat{P_a}[\widehat{h}]`$ $`=`$ $`𝒩^1\mathrm{exp}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\nu }{DLa^2}}{\displaystyle \underset{n=N/2}{\overset{N/2}{}}}|\widehat{h}_{q_n}\left(1e^{\mathrm{i}q_na}\right)|^2\right].`$ (B7)
In the continuum limit $`a0`$ one can explicitly check that Eq.(B7) is equivalent to Eq.(B4) by expanding the term $`e^{\mathrm{i}q_na}`$ in powers of $`a`$ and keeping the lowest non-vanishing order $`O(a)`$. Nevertheless Eq.(B7) leads to problems when one tries to coarse grain the surface. Indeed suppose we now perform a smoothing of the lattice of a factor $`b`$ to obtain a new lattice constant $`a_s=ba`$ and a decreased number of sites $`N_s=N/b`$. This amounts to set to zero all modes from $`N_s`$ to $`N`$ and thus the correct corresponding stationary distribution is according to (B7)
$`\widehat{P_{a_s}}[\widehat{h}]`$ $`=`$ $`𝒩_\mathrm{s}^1\mathrm{exp}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\nu }{DLa^2}}{\displaystyle \underset{n=N_s/2}{\overset{N_s/2}{}}}|\widehat{h}_{q_n}\left(1e^{\mathrm{i}q_na}\right)|^2\right].`$ (B8)
On the other hand, had we started from the “smoothed” lattice and construct the differences $`\delta h_i^\mathrm{s}=h_i^\mathrm{s}h_{i+b}^\mathrm{s}`$ of the coarse-grained heights, we would have found
$`\widehat{P_{a_s}}[\widehat{h}]`$ $`=`$ $`𝒩_\mathrm{s}^1\mathrm{exp}\left[{\displaystyle \frac{1}{2}}{\displaystyle \frac{\nu }{DLa_s^2}}{\displaystyle \underset{n=N_s/2}{\overset{N_s/2}{}}}|\widehat{h}_{q_n}\left(1e^{\mathrm{i}q_nba}\right)|^2\right],`$ (B9)
which differs from the previous one.
It is clear that this is a general problem of all discretizations in real space. In the limit $`a0`$ in fact, one recovers from (B9) the correct expression (B4).
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# 1 Introduction
## 1 Introduction
SUSY/string models contain soft parameters which are in general complex and introduce new sources of CP violation regarding the electric dipole moment (EDM) of the electron and of the neutron. The typical size of these phases in O(1) and they pose a serious EDM problem. Thus the current limits on the electron and the neutron EDM are given by $`|d_e|<4.3\times 10^{27}`$ ecm, $`|d_n|<6.3\times 10^{26}`$ ecm and an order of magnitude analysis shows that the theoretical predictions with phases O(1) are already in excess of the experimental limits. For the minimal supergravity unified model (mSUGRA) the soft SUSY breaking sector is characterized by the parameters $`m_0`$, $`m_{1/2}`$, $`A_0`$ and $`\mathrm{tan}\beta `$, where $`m_0`$ is the universal scalar mass, $`m_{1/2}`$ is the universal gaugino mass, $`A_0`$ is the universal trilinear coupling, and $`\mathrm{tan}\beta `$ is defined by $`\mathrm{tan}\beta =<H_2>/<H_1>`$ where $`H_2`$ is the Higgs that gives mass to the up quark and $`H_1`$ is the Higgs that gives mass to the down quark. In addition one has the Higgs mixing parameter $`\mu `$ which is viewed as the same size as the soft SUSY parameters, and is determined by the constraints of radiative breaking of the electro-weak symmetry. In mSUGRA a set of field redefinitions shows that there are only two independent phases in the theory, and they can be chosen to be $`\alpha _{A_0}`$ and $`\theta _\mu `$ where $`\alpha _{A_0}`$ is the phase of $`A_0`$ and $`\theta _\mu `$ is the phase of $`\mu `$.
The operators that contribute to the electric dipole moments consist of
$`_{}^{}{}_{I}{}^{}={\displaystyle \frac{i}{2}}d_f\overline{\psi }\sigma _{\mu \nu }\gamma _5\psi F^{\mu \nu },_{}^{𝒞}{}_{I}{}^{}={\displaystyle \frac{i}{2}}\stackrel{~}{d}^C\overline{q}\sigma _{\mu \nu }\gamma _5T^aqG^{\mu \nu a}`$
$`_{}^{𝒢}{}_{I}{}^{}={\displaystyle \frac{1}{6}}\stackrel{~}{d}^Gf_{\alpha \beta \gamma }G_{\alpha \mu \rho }G_{\beta \nu }^\rho G_{\gamma \lambda \sigma }ϵ^{\mu \nu \lambda \sigma }`$ (1)
Regarding the color dipole and the purely gluonic dimension six operator one uses the so called naive dimensional analysis $`d_q^C=\frac{e}{4\pi }\stackrel{~}{d}_q^C\eta ^C`$$`d^G=\frac{eM}{4\pi }\stackrel{~}{d}^G\eta ^G`$, where $`\eta ^C`$ $``$ $`\eta ^G`$ $`3.4`$ and $`M`$ =1.19 GeV is the chiral symmetry breaking scale. There are several solutions suggested to control the EDM problem. One possibility is that the phases could be small, or there could be a mass suppression because of the largeness of the sparticle masses. Recently, a new possibility was suggested, i.e., the cancellation mechanism which can control the SUSY EDM problem and there have further developments and applications. The cancellation mechanism works in two stages. First one typically has a cancellation among the $`\stackrel{~}{g},\stackrel{~}{\chi }_i^\pm ,\stackrel{~}{\chi }_k^0`$ exchange contributions to the EDMs. Second there are further cancellation among the electric dipole, the chromoelectric dipole and the purely gluonic contributions. Such cancellations are quite generic in a broad class of SUSY/SUGRA, and in string and D brane models. In addition there are two loop contributions involving axionic Higgs exchange. However, over most of the parameter space such contributions are relatively small.
While most of the analyses to explore the region of cancellations have been numerical in nature, recently there has been an attempt to explore the regions of cancellations also analytically. Such a situation exists in the so called scaling region where $`\mu ^2/M_Z^2>>1`$ and one has $`m_{\chi _1}\stackrel{~}{m}_1,m_{\chi _2}\stackrel{~}{m}_2`$, $`m_{\chi _{3,4}}\mu `$. It was shown in Ref. that in the scaling region one cancellation point in the $`m_0m_{\frac{1}{2}}`$ plane can be promoted to a full trajectory where cancellations occur with only a minor adjustments of parameters. This promotion comes about via the following scaling on $`m_0,m_{\frac{1}{2}}`$
$$m_0\lambda m_0,m_{\frac{1}{2}}\lambda m_{\frac{1}{2}}$$
(2)
With the above scaling and under the constraint of the electro-weak symmetry breaking $`\mu `$ undergoes the following transformation $`\mu \lambda \mu `$ and the total electric dipole $`d_f`$ transforms as $`d_f\lambda ^2d_f`$. Thus the point $`d_f=0`$ is invariant under $`\lambda `$ scaling. Thus if cancellation holds at one point, it holds at other points under scaling by only a small adjustment of parameters and often with no adjustment of parameters at all. As discussed above in mSUGRA one has only two phases after field redefinitions. In the MSSM there are many more phases available. A very general analysis shows that the electric dipole moment of the electron $`d_e`$ depends on 3 phases, while the electric dipole moment of the neutron depends on 9 phases. Together $`d_e`$ and $`d_n`$ depend on 10 phases. The presence of many phases allows for cancellations in larger regions of the parameter space. A similar situation occurs in string and brane models. Of course it may happen that certain models turn out to be free of the EDM problem as is the case in the work of Ref. which also solves the strong CP problem. However, in general large CP phases could exist with a simultaneous resolution to the strong CP problem. For a recent discussion of the possibilities for the resolution of the strong CP problem see Ref..
## 2 SUSY Dark Matter
There are 32 new particles in MSSM and any one of these particles could be an LSP. In SUGRA models, however, one finds that starting with prescribed boundary conditions at the GUT scale with gravity mediated breaking of supersymmetry one finds that the model predicts the lightest neutralino to be the LSP over most of the parameter space of the model. Further, with R parity invariance the LSP will be stable and thus the lightest neutralino is predicted to be a candidate for cold dark matter over most of the parameter space in SUGRA models. Many analyses of supersymmetric dark matter already exist in the literature. These include effects of FCNC constraints from $`bs+\gamma `$, the effects of non-universalities of scalar masses, effects of non-universalities of gaugino masses and effects of co-annihilation. Recently, effects of uncertainties in the WIMP velocity in the direct and in the indirect detection of dark matter have been analyzed and analyses have also been given of the effects of uncertainties of the quark mass densities on the direct detection rates. In this paper we discuss the effects of CP violation on direct detection.
## 3 CP Effects on Dark Matter
The effects of CP violation on the relic density have been discussed in Refs. Here we discuss the effects of CP violation on event rates. The effective Lagrangian with CP violation is gotten from the micropscopic SUGRA lagrangian by integration on the Z, Higgs, and sfermion poles and one finds
$`_{eff}=\overline{\chi }\gamma _\mu \gamma _5\chi \overline{q}\gamma ^\mu (AP_L+BP_R)q+C\overline{\chi }\chi m_q\overline{q}q+D\overline{\chi }\gamma _5\chi m_q\overline{q}\gamma _5q`$
$`+E\overline{\chi }i\gamma _5\chi m_q\overline{q}q+F\overline{\chi }\chi m_q\overline{q}i\gamma _5q`$ (3)
Here A and B are spin dependent terms arising from the Z boson exchange and squark exhange and is given by
$$A=\frac{g^2}{4M_W^2}[|X_{30}|^2|X_{40}|^2][T_{3q}e_qsin^2\theta _W]\frac{|C_{qR}|^2}{4(M_{\stackrel{~}{q1}}^2M_\chi ^2)}\frac{|C_{qR}^{^{}}|^2}{4(M_{\stackrel{~}{q2}}^2M_\chi ^2)}$$
(4)
$$B=\frac{g^2}{4M_W^2}[|X_{30}|^2|X_{40}|^2]e_qsin^2\theta _W+\frac{|C_{qL}|^2}{4(M_{\stackrel{~}{q1}}^2M_\chi ^2)}+\frac{|C_{qL}^{^{}}|^2}{4(M_{\stackrel{~}{q2}}^2M_\chi ^2)}$$
(5)
where $`C_{qR}`$ etc are defined in Ref. and $`X_{n0}`$ give the gaugino-Higgsino content of the LSP and is defined by
$$\chi ^0=X_{10}^{}\stackrel{~}{B}+X_{20}^{}\stackrel{~}{W}+X_{30}^{}\stackrel{~}{H}_1+X_{40}^{}\stackrel{~}{H}_2$$
(6)
where $`\stackrel{~}{B}`$ is the Bino, $`\stackrel{~}{W}`$ is the neutral Wino, and $`\stackrel{~}{H}_1`$ and $`\stackrel{~}{H}_2`$ are the Higgsinos corresponding to the Higgs $`H_1`$ and $`H_2`$. In Eq.(3) C governs the scalar interaction which arises from the CP even Higgs exchange and from the sfermion exhange and gives rise to coherent scattering. It is given by
$$C=C_{\stackrel{~}{f}}+C_{h^0}+C_{H^0}$$
(7)
where
$$C_{\stackrel{~}{f}}(u,d)=\frac{1}{4m_q}\frac{1}{M_{\stackrel{~}{q1}}^2M_\chi ^2}Re[C_{qL}C_{qR}^{}]\frac{1}{4m_q}\frac{1}{M_{\stackrel{~}{q2}}^2M_\chi ^2}Re[C_{qL}^{^{}}C_{qR}^{{}_{}{}^{}}]$$
(8)
$$C_{h^0}(u,d)=(+)\frac{g^2}{4M_WM_{h^0}^2}\frac{\mathrm{cos}\alpha (sin\alpha )}{\mathrm{sin}\beta (cos\beta )}Re\sigma $$
(9)
$$C_{H^0}(u,d)=\frac{g^2}{4M_WM_{H^0}^2}\frac{\mathrm{sin}\alpha (cos\alpha )}{\mathrm{sin}\beta (cos\beta )}Re\rho $$
(10)
In the above (u,d) exhibit the quark flavor in the scattering and $`\alpha `$ stands for the Higgs mixing angle while $`\sigma `$ and $`\rho `$ are given by
$`\sigma =X_{40}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{cos}\alpha +X_{30}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{sin}\alpha `$
$`\rho =X_{40}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{sin}\alpha +X_{30}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{cos}\alpha `$ (11)
The D term in Eq.(3) arises from the exchange of the CP odd Higgs $`A^0`$
$$D(u,d)=C_{\stackrel{~}{f}}(u,d)+\frac{g^2}{4M_W}\frac{cot\beta (tan\beta )}{m_{A^0}^2}Re\omega $$
(12)
while the terms E and F arise only in the presence of CP violation and are given by
$$E(u,d)=T_{\stackrel{~}{f}}(u,d)+\frac{g^2}{4M_W}[(+)\frac{cos\alpha (sin\alpha )}{\mathrm{sin}\beta (cos\beta )}\frac{Im\sigma }{m_{h^0}^2}+\frac{sin\alpha (cos\alpha )}{\mathrm{sin}\beta (cos\beta )}\frac{Im\rho }{m_{H^0}^2}]$$
(13)
$$F(u,d)=T_{\stackrel{~}{f}}(u,d)+\frac{g^2}{4M_W}\frac{cot\beta (tan\beta )}{m_{A^0}^2}Im\omega $$
(14)
where $`\omega `$ is given by
$$\omega =X_{40}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{cos}\beta +X_{30}^{}(X_{20}^{}\mathrm{tan}\theta _WX_{10}^{})\mathrm{sin}\beta $$
(15)
and
$$T_{\stackrel{~}{f}}(q)=\frac{1}{4m_q}\frac{1}{M_{\stackrel{~}{q1}}^2M_\chi ^2}Im[C_{qL}C_{qR}^{}]+\frac{1}{4m_q}\frac{1}{M_{\stackrel{~}{q2}}^2M_\chi ^2}Im[C_{qL}^{^{}}C_{qR}^{{}_{}{}^{}}]$$
(16)
In the limit when CP phases vanish, the above formulae limit correctly to previous analyses in the absence of CP phases. Numerical analysis shows that the coefficients A-F exhibit a strong dependence on CP phases. Typically, however, the terms D, E and F make only small contributions and the terms A, B and C generally dominate the scattering. The analysis of event rates follows the method of Ref.. The analysis including the CP violating phases but without the imposition of the EDM constraints is displayed in Fig.1 where the ratio R/R(0) is plotted as a function of $`\theta _\mu `$, where R/R(0) is the ratio of the event rates with and without CP violation effects. The analysis shows that the CP violating phases can generate variations in the event rates up to 2-3 orders of magnitude. A similar analysis but with inclusion of the EDM constraints in given in Fig.2. Here one finds that the effects are much reduced, i.e., around a factor of 2 variation over the allowed range of phases. In Ref. the analysis included only the two phases $`\alpha _{A0}`$ and $`\theta _\mu `$. However, for nonminimal models we have three $`\xi `$ phases in the the gaugino mass sector. Only one of these three phases, i.e. $`\xi _1`$, enters the expressions of direct detection through the neutralino mass matrix. Among the remaining two phases, $`\xi _2`$ affects the EDM of the electron and of the neutron while $`\xi _3`$ affects only the EDM of the neutron. Using these differential effects generated by $`\xi _1`$, $`\xi _2`$ and $`\xi _3`$ we can arrange cancellations for the EDMs to satisfy the EDM constraints and at the same time generate a large effect on the direct detection of neutralinos.
## 4 Conclusions
In a large class of SUSY, string and brane models there are new sources of CP violation arising from the soft breaking sector of the theory. Since the natural size of these CP phases is O(1) there exists a priori a serious EDM problem. The cancellation mechanism is a possible solution to the EDM problem with large CP phases. Detailed analyses show that there exists a significant part of the parameter space where large CP phases are compatible with the current experiment on the EDMs. The existence of large CP phases can have significant effects on low energy SUSY phenomenology, and in this paper we have discussed the effects of large CP phases on event rates in the direct detection of dark matter. We emphasize that the inclusion of CP phases in the dark matter analysis without the inclusion of EDM constraints can lead to erroneously large effects since the CP effects can change the event rates by several orders of magnitude. With the inclusion of the EDM constraints the CP effects are much smaller although still significant enough to be included in any precision analysis of dark matter. These results are of import in view of the ongoing and future dark matter experiments. In addition to their effects on dark matter, large CP phases will also affect searches for SUSY at the Tevatron, at the LHC and in B physics and it is imperative that one include CP effecs in future SUSY searches to cover the allowed parameter space of models. The cancellation mechanism is a testable idea. Thus if the cancellation idea is right, the EDMs of the electron and of the neutron should become visible with an order of magnitude improvement in experiment. Such a possibility exists with experiments underway to improve the sensitivity of the measurements on the electron and on the neutron electric dipole moments.
Acknowledgements
This research is supported in part by the NSF grant PHY-9901057.
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# Pulsars in the Westerbork Northern Sky Survey
## 1 Introduction
Pulsars can have a steep radio spectrum at frequencies around 1 GHz and can be highly polarised (e.g. Manchester & Taylor 1977) . Synthesis maps can therefore be used to select pulsar candidates on the basis of their steep spectrum and/or high degree of polarisation. These sources are later observed with a high time resolution instrument to search for pulsations. At least two pulsars, PSR B1937+21 Backer et al. (1982) and PSR J0218+4232 Navarro et al. (1995), have been found in this way, while they were missed in regular pulsar surveys because their pulses were smeared out in the detection process due to their small period and their high dispersion measure.
The pulsar population found in this way, may supplement the presently known population, since this method has totally different selection effects. These effects can be investigated by studying the spectral indices and polarisation degrees of known pulsars in these continuum observations. The spectral indices are strongly influenced by scintillation. This may cause the flux density of a pulsar to vary by more than 100 percent on time scales of minutes (diffractive scintillation) to days (refractive scintillation).
In this paper I search for detections of pulsars in the Westerbork Northern Sky Survey (WENSS), a survey performed at 325 MHz. I compare my results with those from Kaplan et al. Kaplan et al. (1998) and Han & Tian Han & Tian (1999), who did similar analyses with data from the 1400 MHz NRAO VLA Sky Survey (NVSS). Sect. 2 describes the WENSS and Sect. 3 describes how the pulsar catalog was correlated with the WENSS source catalog. In Sect. 4 the positions of the non-detected pulsars are searched for flux densities above three times the local noise level. Sect. 5 discusses the remaining non-correlations. Sect. 6 combines the WENSS results with the NVSS correlation studies to determine the spectral indices and compares these with values reported in other literature. Finally, in Sect. 7 the role of scintillation is discussed.
## 2 WENSS
The Westerbork Synthesis Radio Telescope (WSRT) is an east-west array, consisting of fourteen 25 meter dishes Baars & Hooghoudt (1974). A twelve hour observation is needed to get spatial resolution in all directions. Ten of the dishes have fixed positions and are spaced 144 meters apart. The remaining four are moveable, although their mutual distances are usually kept constant. By observing an object several times with different distances between the fixed and moveable subarray, different baselines are obtained and the $`uv`$-coverage improves, yielding an improved synthesised antenna pattern.
The Westerbork Northern Sky Survey (WENSS) is a 325 MHz continuum survey of the sky above declination +30°(Rengelink et al. 1997, de Bruyn et al. in preparation). This area was surveyed using a mosaicing technique. Each field was observed 18 times for 20 seconds spread over 12 hours. Each mosaic was observed on six days with a different spacing between the fixed and moveable subarray to get a uniform spatial distribution. These observations were spread over periods of weeks to years. The resulting flux densities are averaged over all these observations.
The WENSS beam size is 54″$`\times `$ 54″ $`\mathrm{cosec}\delta `$ (FWHM). The final maps have pixel sizes of 21.09″. When a pixel was found with a flux density above five times the local noise level, a two-dimensional Gaussian was fitted to its surroundings. The coordinates of the centroid of the fit, the maximum (peak) flux density and the flux density integrated over the fit were added to the source list. Extra flags in the catalog mark multiple and extended sources. For a point source the peak flux density equals the flux density integrated over the beam. Only the peak flux densities are used in this analysis, since pulsars are intrinsically point sources. Bright sources have positional errors of 1.5″, weaker sources at lower declinations have errors up to about 10″. On average the uncertainty is 5″.
The total bandwidth was 5 MHz. For most regions of the sky the detection limit was between 15 and 25 mJy (five times the local noise level). The polar cap (declination $`>`$ +75°) was surveyed with a larger bandwidth and the detection limit for this area was about 10 – 12 mJy. The WENSS source catalog contains 229420 sources de Bruyn et al. (1998). A total of 18186 of these are located in the polar cap area. Fig. 1 shows the density of sources in the WENSS area.
## 3 Detected pulsars
### 3.1 Positional correlation
The radio pulsar catalog (Taylor et al., 1993, Taylor et al., 1995) contains 84 entries in the part of the sky that was covered by the WENSS area. The typical positional uncertainty for a pulsar with a flux density greater than 10 mJy is 0.1″ or less. In most cases the uncertainty in the pulsar position is negligible compared to the positional uncertainty of WENSS sources.
The pulsar proper motions are neglected, since for each pulsar the change in position between the epoch of discovery and the epoch of the WENSS is less than 0.5″ in right ascension and less than 0.8″ in declination for all these pulsars. This is much smaller than the uncertainties in the WENSS positions.
Seven pulsars in the WENSS area have large positional errors of about 4′. The probability that a WENSS source is located by chance within a circle of three times this positional error is more than 90 percent, if it assumed that the WENSS sources are uniformly distributed over the sky. Therefore, I have excluded these 7 pulsars (PSRs J0417+35, B1639+36B, J1758+30, J1900+30, J1931+30, J2002+30 and J2304+60) from further analysis. PSR B2000+40 was also excluded, since it is located in the Cygnus A region, where no WENSS map could be made.
I have taken the J2000 positions of the remaining 76 pulsars and compared them with the positions of the sources in the WENSS catalog. Twenty-five pulsars have a WENSS source located within three times their combined positional uncertainty $`\sigma `$, with
$$\sigma =\sqrt{\left(\frac{\mathrm{\Delta }\alpha }{\sigma _\alpha }\right)^2+\left(\frac{\mathrm{\Delta }\delta }{\sigma _\delta }\right)^2},$$
(1)
where $`\mathrm{\Delta }\alpha `$ and $`\mathrm{\Delta }\delta `$ are the positional differences in right ascension and declination, respectively, and
$$\sigma _\alpha =\sqrt{\sigma _{\alpha ,\mathrm{WENSS}}^2+\sigma _{\alpha ,\mathrm{PSR}}^2}$$
(2)
and similarly for $`\sigma _\delta `$. Table 1 lists the J2000 positions of the correlated pulsars and WENSS sources and their offset, both in arcseconds and in $`\sigma `$.
In Fig. 1 it can be seen that the pulsars are located in areas with different values for the WENSS source density. The probability of a change coincidence can be approximated by a probability calculation that assumes a uniform distribution of the WENSS sources. In that case, the probability that an individual correlation is just by chance is 0.0012. The binomial probability, that one out of 76 trials gives a chance correlation is 0.083. The probability that two correlations occur by chance is 0.004.
Fig. 2 displays the distribution of positional differences between a pulsar and its nearest WENSS source in units of their combined positional uncertainty $`\sigma `$. There is a clear gap between the correlated pairs (difference less than 3$`\sigma `$) and the non-correlated ones. The distribution of the positional difference ($`\mathrm{\Delta }`$) for the related pairs is
$$P(\mathrm{\Delta })=\mathrm{\Delta }e^{\mathrm{\Delta }^2/2},$$
(3)
which follows after a conversion to polar coordinates. This distribution is also plotted in Fig. 2. A Kolmogorov-Smirnov test assigns a 15 procent probability that our sample is drawn from the distribution (3).
There are two objects with positional differences between 3 and 10 $`\sigma `$. These are in confused regions and will be discussed in Sect. 4.
### 3.2 Flux densities
Table 2 lists the observed flux densities of the correlated WENSS sources and their uncertainties. The flux densities can be compared with known pulsar flux densities, but these are usually measured at frequencies of 400 MHz and higher. By assuming a power law with constant spectral index, these flux density data can be extrapolated. I have used flux density data from Lorimer et al. Lorimer et al. (1995), hereafter LYLG. These flux densities are averaged over many observations spread over years and their uncertainties include variations due to scintillation. LYLG provide data for 24 of the 25 pulsars in Table 1. Flux densities for PSR J0218+4232 are taken from its discovery paper Navarro et al. (1995).
Pulsar flux densities usually obey a power law with a negative exponent ($`S\nu ^\alpha `$ with $`\alpha <0`$) in the frequency range from 325 to 1400 MHz. I have fitted the logarithms of the flux densities with a straight line. These lines are plotted in Fig. 3. From this fit a flux density at 325 MHz is estimated. This estimate is plotted against the flux density of the WENSS counterpart in Fig. 4. It is known that some pulsars have a low frequency turnover, usually located around 100 MHz Malofeev et al. (1994). However, some pulsars exhibit a turnover at a higher frequency: PSR B0329+54 around 300 MHz Lyne & Rickett (1968) and PSR B2021+51 around 400 MHz. The spectrum of PSR B2319+60 is flat between 200 and 600 MHz Malofeev et al. (1994). PSRs B1946+35 and B2154+40 also have a turnover, but it is not clear whether this located between 325 and 400 MHz Malofeev et al. (1994). In these cases the assumption of a constant power law is not correct and this results in an overestimated flux density at 325 MHz.
Navarro et al. Navarro et al. (1995) discovered that the flux density of PSR J0218+4232 has a non-pulsed component. They find that the continuum flux density at 325 MHz varies between 100 and 200 mJy. I used this flux density estimate for a comparison with the WENSS flux density. An estimate from the pulsed flux density was not derived.
### 3.3 Source maps
Fig. 5 shows maps of the WENSS sources which are correlated with a pulsar. Most sources are consistent with a point source convolved with the beam shape. The shapes of PSR B0329+54 and PSR B0809+74 are affected by scintillation (see Sect. 7). The integrated and peak flux densities of PSR J0218+4232 differ by 20 percent. This suggests that the source is extended. However, a two-dimensional fit shows that its shape is not significantly different from the beam shape.
## 4 Marginally detected pulsars
The WENSS source finding routines have a flux density threshold at five times the local noise level ($`5S_{\mathrm{noise}}`$). Since in this paper single trials are done to find a correlation at a very small number of positions, a $`3S_{\mathrm{noise}}`$ source can be marked as a marginal detection. However, source fitting is less acurate. Therefore, a 5 by 5 pixel box (1.75′) centered around the pulsar position is searched for the pixel with the maximum flux density. This pixel was marked as a marginal detection, if its flux density was greater than $`3S_{\mathrm{noise}}`$. The positional uncertainty for these marginal sources is estimated using the same equation as for the fitted sources Rengelink et al. (1997). The positional error for these marginal sources is approximately 42″. Since a 1.75′ box is searched, no marginal detection is missed.
The probability of finding a $`3S_{\mathrm{noise}}`$ (or higher) pixel in a map with Gaussian noise is $`2.710^3`$. About 5 WENSS FWHM beams fit into a 5 $`\times `$ 5 pixel box. Fifty-one boxes were searched. Effectively, 255 trials have been performed. The binomial probability that one $`3S_{\mathrm{noise}}`$ pixel is found is 0.35, the probability that two are found is 0.12 and that three are found is 0.03. Therefore, there is a 50 percent probability that (at least) one of our 14 marginal detection is just a noise fluctuation.
The maps of the marginal detections are shown in Fig. 6. The location of the pixel with maximum flux density is also indicated. The shapes of these sources is not as point-like as the strong detections with $`S>5S_{\mathrm{noise}}`$. The flux densities of the marginal detections are listed in Table 3, together with an estimate of the pulsar flux density based on a similar extrapolation of the pulsar spectrum as done in Sect. 3.2. The estimate for PSR J1518+4904 is based on its measured flux density at 370 MHz Sayer et al. (1997), since its spectrum as plotted by Kramer et al. Kramer et al. (1999) shows evidence for a low frequency turnover.
The ratios of the extrapolated pulsar flux density and the WENSS source flux density are displayed in Fig. 4. The spread is of the order of a factor 1.5, which is comparable with the spread for the detected sources that were discussed in the previous section. Five sources were detected, although their expected flux density was below three times the local noise level (see Table 3). The flux density at 325 MHz could not be estimated for three other pulsars, since no reliable flux density data at other frequencies were available.
The contours of four pulsars are confused by nearby radio sources. These sources are discussed in the following and are shown in Fig. 7.
*PSR B0655+64:* the WENSS source and the pulsar position are 4.3$`\sigma `$ apart. The estimated flux density of the pulsar at 325 MHz is 7 $`\pm `$ 2 mJy, but the WENSS source is 23 $`\pm `$ 5 mJy. Also, the NVSS (see Sect. 7) shows a radio source at the WENSS position and clearly away from the pulsar position. Its flux density at 1400 MHz is about 5.6 mJy, while the pulsar flux density is expected to be 0.3 $`\pm `$ 0.1 mJy. The pulsar has a proper motion, but it is small and directed towards negative declinations.
*PSR B1112+50:* The WENSS source is bright (135 mJy) and has an accurate position. The separation between the fitted WENSS position and the known pulsar position is 4.6$`\sigma `$. Extrapolation of the pulsar spectrum results in an estimated flux density at 325 MHz of 16 $`\pm `$ 3 mJy, much less than that of the WENSS source. Kaplan et al. Kaplan et al. (1998) and Han & Tian Han & Tian (1999) searched the NVSS for pulsar counterparts and also noted that the pulsar is confused by a strong NVSS source 12″ away.
*PSR B1951+32:* The coincident source in the WENSS catalog is marked as extended and the emission is dominated by the supernova remnant CTB80. The pulsar is associated with this remnant (Strom 1987, Kulkarni et al. 1988). The WENSS peak flux density is 983 mJy, which is about a factor 70 stronger than the expected pulsar flux density.
*PSR B2306+55:* It can be clearly seen in the map that this source has two components, of which the weaker one is probably the counterpart to the pulsar. This component is not listed in the WENSS source list. The estimated pulsar flux density is about 30 $`\pm `$ 3 mJy. The bright component of the WENSS source is 125 $`\pm `$ 9 mJy, separated 28 $`\sigma `$ from the pulsar position. The second component has a flux density of approximately 24 mJy.
In all these four cases I conclude that the source in the WENSS catalog and the pulsar are unrelated. Galama et al. Galama et al. (1997) reached the same conclusion for PSR B0655+64.
## 5 Non-detected pulsars
Table 4 lists the pulsars that have no counterpart in the WENSS. In 14 cases the expected pulsar flux density is higher than three times the local noise level. Still, the pulsar was not detected. In two cases no reliable pulsar flux density estimate at 325 MHz is available. In three other cases the estimate is based on 400 MHz observations. In case of PSRs B0841+80 and B1839+36A this was done, because there was no spectral information available. The spectrum of PSR J1012+5307 might also have a low frequency turnover (see its spectrum as plotted by Kramer et al. 1999) . Its flux density is known to vary by up to a factor four from its mean value of 30 mJy Nicastro et al. (1995).
Five (and possibly eight) pulsars are detected, which were expected to be not detectable. The number of non-detected sources that were expected to have a flux density greater than the detection limit, should be roughly the same as the number of unexpected detections. The difference may be due to Poisson fluctuations in the (small) number of pulsars in this study.
## 6 Sources in the NVSS
The NRAO VLA Sky Survey (NVSS, Condon et al. 1998) is a survey of the sky above declination $``$40° at 1400 MHz. Its final resolution is 45″ and is comparible with the WENSS resolution of 54″ (in right ascension). The bandwidth was effectively 42 MHz. The noise in the final NVSS maps is about 0.45 mJy. These maps are constructed by taking the average of a number of snapshot observations. Each point on the sky is observed in about three snapshots.
Seventeen of the twenty-five pulsars with a WENSS counterpart are also detected in the NVSS (see Kaplan et al. Kaplan et al. (1998), who did not include PSR B0138+59 and Han & Tian Han & Tian (1999)). Two pulsars (PSR B0353+52 and PSR J1518+4904) with a marginal WENSS counterpart coincide with a NVSS source. Kaplan et al. Kaplan et al. (1998) also lists the WENSS sources which coincide with their NVSS sources. Besides the three pulsars already mentioned, PSR B0809+74 is also not in their list, though it matches their coincidence criterion. Eight pulsars with a WENSS counterpart are not detected in the NVSS. This can be caused by scintillation or these pulsars must have a steep spectrum (spectral index between about $``$1.6 and $``$2.0). Two pulsars (PSRs J0218+4232 and B2319+60) may have a very steep spectrum (spectral index less than $``$2.6 and $``$2.3, respectively).
From the WENSS and NVSS flux density values the spectral index can be calculated. These can be compared with the spectral indices as determined from a large number of dedicated pulsar observations, as by LYLG. The latter indices will suffer less from scintillation effects. The uncertainty due to scintillation is included in the errors given by LYLG.
From Table 5 it is clear that there are some large differences between the spectral index as determined by the WENSS and NVSS flux densities and the long term averaged spectral index. The differences are displayed in Fig. 8. This histogram has a mean of $``$0.14 and a standard deviation of 0.48. There is a large spread and a WENSS-NVSS source flux density comparison is not a very precise way to determine a spectral index. This will limit the success of a pulsar candidate selection based on their WENSS-NVSS spectral index.
## 7 Discussion
The deviations between the flux density values of pulsars and the correlated WENSS and NVSS sources are mainly caused by scintillation effects. Small-scale inhomogeneities in the interstellar medium affect the travel path of radio waves and can amplify or weaken them. For reviews, see Rickett Rickett (1990) or Narayan Narayan (1992). Walker Walker (1998) summarizes the involved equations and dependencies.
In the case of strong scintillation, the phase changes due to the scattering in a certain region are larger than the changes in phase due to normal geometry. The scintillation is called weak in the opposite case. The boundary between the two is dependent on the frequency of the radio waves and the distance to the pulsar. At 325 and 1400 MHz almost all pulsars are in the strong scintillation regime.
Two types of strong scintillation exist: diffractive and refractive. They differ in their typical timescale $`\tau `$, frequency bandwidth $`\mathrm{\Delta }\nu `$ and size of the resulting flux density variation. This strength of the scintillation is usually quantified by the modulation index, i.e. the rms fractional flux density variation.
In the case of strong scintillation waves from multiple locations in the scattering region interfere constructively (or destructively). Both the typical timescale and the frequency bandwidth are small. The modulation index equals one. The timescale is dependent on the relative velocities of the pulsar, the interstellar medium and the Earth. For the WENSS observations of pulsars typical timescales are 1 to 10 minutes and typical bandwidths are 10 Hz to 500 kHz. This means that for almost all pulsars any variations due to diffractive scintillation are averaged out over the 5 MHz WENSS bandwidth and when the 6 $`\times `$ 18 short observations spread over 6 $`\times `$ 12 hours are combined.
Only for PSR B0809+74 are the diffractive scintillation timescale and bandwidth large enough that some effect remains. From the equations given by Walker Walker (1998) one finds a timescale of 12 minutes and a bandwidth of 500 kHz. The actual value of the typical scintillation timescale and bandwidth are even higher, since several authors have already shown that the Taylor-Cordes distance model Taylor & Cordes (1993) gives too small predictions for this pulsar (e.g. Rickett et al. 2000). Diffractive scintillation can explain the WENSS image of this pulsar (Fig. 5). As 40 minutes is the time between two observations of a field in one mosaic, flux density variations on that time scale cause spokes in the map. Such spokes can also be seen in the complete map of PSR B0329+54 (Fig. 9).
Refractive scintillation is caused by the focussing effect of a large scattering region. The timescales and bandwidths involved are much larger than in the case of diffractive scintillation. The modulation index is also smaller. The scintillation bandwidth is of the order of the observing frequency. Refractive time scales for the pulsars detected in the WENSS vary from a couple of days to several years and the expected modulation index from 0.05 to 0.3.
If the refractive time scale is less than the time between observations of the same mosaic (couple of days to several years), any flux density variation will be averaged out when the mosaics are combined. However, stong flux density variations between 12 hour sessions cause a ring at the first grating ring of the synthesised beam. Since the mutual distances between the dishes are multiples of 72 m, the ring will have radius of 72 m / ($`c/`$325 MHz) radians, i.e. 44′ in right ascension and 44′ $`\times \mathrm{cosec}\delta `$ in declination. The second and higher grating rings are not visible, since data far from the field center gets a low weighting factor when the final image is created.
It is hard to give good estimates for the expected diffractive and refractive scintillation timescales. They depend on the often poorly known pulsar velocity. For large pulsar velocities (compared to the velocity of the interstellar medium, being about 50 km/s) the dependence is as one over the square root of this velocity. Since some pulsar velocities might be up to several hundred kilometers per second, this cannot be neglected. I divided the pulsars that are detected in the WENSS in two groups, based on their expected refractive modulation index if their velocity is neglected. Both groups had similar relative deviations between their measured and expected flux densities.
The observed modulation index is about 0.4, much larger than the expected value of about 0.2. The WENSS pulsar flux densities vary more due to refractive scintillation than predicted by the equations. This has been observed before by several authors and is attributed to the assumption that the turbulence in the interstellar medium has a Kolmogorov spectrum (e.g. Blandford & Narayan, 1985).
The NVSS flux densities are even more affected by scintillation effects. Each point in the NVSS maps is an average of about three snapshots. Two of these three are taken right after each other (snapshot series are taken at constant declination and increasing right ascension, see figure 7 in Condon et al. 1998) and very little averaging takes place. The expected refractive modulation index at 1400 MHz is larger than at the WENSS frequency. The expected diffractive frequency bandwidth is also larger at the NVSS frequency, even relative to the increased total bandwidth used in the NVSS. One therefore expects the differences between the pulsar flux densities in the NVSS and the flux densities reported by LYLG to be larger than the differences reported in this study. This is indeed observed: the modulation index of the WENSS sources in Table 2, excluding PSRs J0218+4232, B0329+54 and B2021+51 is 0.40, the modulation index of the NVSS sources in table 1 of Han & Tian Han & Tian (1999) is 0.60.
De Breuck et al. De Breuck et al. (2000) show that the total spectral index distribution of WENSS sources that are correlated with a NVSS source ($`S_{325}>50`$ mJy and $`b>15^{}`$) has a mean of $``$0.80 and a standard deviation of 0.24. Pulsar have a mean spectral index of $``$1.6 (see LYLG). An increase of 0.5 due to scintillation will move the pulsar spectral index well into the distribution of normal sources, which are much more frequent. The spectrum of some point-like quasars will also be affected by scintillation and some of them will seem to have a much steeper spectrum than they really have. This effect should be taken into consideration if pulsar candidates are selected on the basis of their spectral index derived from the WENSS and NVSS.
###### Acknowledgements.
I thank F. Verbunt and A. G. de Bruyn for discussion and comments. I thank B. W. Stappers for comments on the manuscript. I am supported by The Netherlands Research School for Astronomy (NOVA), a national association of astronomy departments at the Universities of Amsterdam, Groningen, Leiden and Utrecht.
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# Dynamics of vortex tangle without mutual friction in superfluid 4He
## I INTRODUCTION
Superfluid <sup>4</sup>He (Helium II) behaves like an irrotational ideal fluid, whose characteristic phenomena can be explained well by the Landau two-fluid model. However, superflow becomes dissipative (superfluid turbulence) above some critical velocity. The concept of superfluid turbulence was introduced by Feynman who stated that the superfluid turbulent state consists of a disordered set of quantized vortices, called vortex tangle(VT). Reminding the inertial range of the classical-fluid turbulence, Feynman proposed that VT undergoes the following cascade process. At zero temperature, a large distorted vortex loop breaks up to smaller loops through reconnections, and the cascade process continues self-similarly down to the order of the interatomic scale. At finite temperatures, however, normal fluid collides with vortices and takes energy from them.
This idea was developed further by Vinen. In order to describe an amplification of a temperature difference at the ends of a capillary retaining thermal counterflow, Gorter and Mellink introduced some additional interactions between the normal fluid and superfluid. Through experimental studies of the second-sound attenuation, Vinen considered this Gorter-Mellink mutual friction in relation to the macroscopic dynamics of the VT. Assuming homogeneous superfluid turbulence, Vinen obtained an evolution equation for the vortex line density (VLD) $`L(t)`$, what we call the Vinen’s equation
$$\frac{dL}{dt}=\alpha |𝒗_{\mathrm{ns}}|L^{3/2}\chi _2\frac{\kappa }{2\pi }L^2,$$
(1)
where $`\alpha `$ and $`\chi _2`$ are parameters dependent on temperature and $`𝒗_{\mathrm{ns}}`$ is the relative velocity between the normal flow and superflow, $`\kappa `$ the quantized circulation. This Vinen’s theory could describe well a large number of observations of mostly stationary cases.
However the nonlinear and nonlocal dynamics of vortices had long delayed the progress in further microscopic understanding of the VT. It was Schwarz that broke through. His most important contribution was that the direct numerical simulation of vortex dynamics connected with the scaling analysis enabled us to calculate such physical quantities as the VLD, some anisotropic parameters, the mutual friction force, etc. The observable quantities obtained by Schwarz’s theory agree well with the experimental results of the steady state of the VT. This research field pioneered by Schwarz has revealed many problems of vortex dynamics, such as the flow properties in channels , sideband instability of Kelvin waves, vortex array in rotating superfluid, vortex pinning.
The mutual friction plays an important role in the above vortex dynamics. The stationary state of the VT Schwarz obtained is self-sustaining, and realized by the competition between the excitation and dissipation due to the mutual friction subject to the $`𝒗_{\mathrm{ns}}`$ field, as described in the next section. Hence the system free from the mutual friction cannot sustain the stationary VT.
Compared with the steady state, there have been less studies of the transient behavior of the VT. Although the transient behavior generally refers to both the growth and decay process, this paper considers only the decay of the VT after the driving velocity is suddenly reduced to zero. The early measurements by Vinen and the later ones observed a decay of the VT which was consistent with the Vinen’s equation (1) with only the decay term, although Schwarz and Rozen coupled the Vinen’s equation with the hydrodynamical equations of the normal flow and the superflow in order to explain a slow decay following an initial rapid decay they observed. Apart from these experiments on thermal counterflow, the decay of vorticity in turbulence generated by towing a grid was studied recently. This turbulence is expected to be homogeneous and isotropic. The experimental results may be understood by the picture that the mutual friction can be so strong that the normal fluid and the superfluid lock together, behaving effectively like a single fluid. The experimental results are compared with the change in the turbulent energy spectrum which includes the Kolmogorov law.
Both these numerical and experimental results are much affected by the mutual friction. However, recently, Davis et al. observed that vortices did decay even at mK temperatures where the normal fluid density became vanishingly small and, as a consequence, the mutual friction did not work effectively. The vortices were created by a vibrating grid, and detected by their trapping of negative ions. The first important point is that the vortices actually decay at such low temperatures. The second is that the decay rate becomes independent of temperature below $`T70`$mK. It is unclear how the vortices decay. This experimental work, which is just preliminary at present, can develop a new research field of superfluid or vortex dynamics at mK temperatures; it can reveal some essence that may be covered with the normal fluid at higher temperatures.
Motivated by this experimental work, we study numerically the vortex dynamics without the mutual friction. The calculation under the localized induction approximation(LIA) is made for the dense VT, while the full Biot-Savart calculation for the more dilute vortices. The absence of the mutual friction makes the vortices kinked, which promotes vortex reconnections. Consequently small vortices are cut off from a large one through the reconnections. The resulting vortices also follow the self-similar process to break up to smaller ones. Although our formulation cannot describe the final destiny of the minimum vortex, the decay of the VT is found to be connected with this cascade process, which is just the cascade process at zero temperature Feynman proposed.
The contents of this paper are the following. Section II describes the equations of motion of vortices and the method of numerical calculation. Section III studies the dynamics of dilute vortices under the full Biot-Savart law both without and with solid boundaries; this calculation reveals the essence of the cascade process. The dynamics of the dense VT under the LIA is discussed in Sec. IV. The obtained results are compared with the solution of the Vinen’s equation in Sec. V. The agreement is good, which supports the picture of the cascade process. The decay of the VT subject to the mutual friction is discussed too. Section VI is devoted to conclusions and discussions.
## II EQUATIONS OF MOTION AND NUMERICAL SIMULATION
A quantized vortex is represented by a filament passing through the fluid and has a definite direction corresponding to its vorticity. Except for the thin core region, the superflow velocity field has a classically well-defined meaning and can be described by ideal fluid dynamics. The velocity produced at a point $`𝒓`$ by a filament is given by the Biot-Savart expression :
$$𝒗_{s,\omega }=\frac{\kappa }{4\pi }_{}\frac{(𝒔_1𝒓)\times d𝒔_1}{|𝒔_1𝒓|^3},$$
(2)
where $`\kappa `$ is the quantized circulation. The filament is represented by the parametric form $`𝒔=𝒔(\xi ,t)`$, $`𝒔_1`$ refers to a point on the filament and the integration is taken along the filament. The Helmholtz’s theorem for a perfect fluid states that the vortex moves with the superfluid velocity at the point. Attempting to calculate the velocity $`𝒗_{s,\omega }`$ at a point $`𝒓=𝒔`$ on the filament makes the integral diverge as $`𝒔_1𝒔`$. To avoid it, we divide the velocity $`\dot{𝒔}`$ of the filament at the point $`𝒔`$ into two components :
$$\dot{𝒔}=\frac{\kappa }{4\pi }𝒔^{}\times 𝒔^{\prime \prime }\mathrm{ln}\left(\frac{2(\mathrm{}_+\mathrm{}_{})^{1/2}}{e^{1/4}a_0}\right)+\frac{\kappa }{4\pi }_{}^{^{}}\frac{(𝒔_1𝒓)\times d𝒔_1}{|𝒔_1𝒓|^3}.$$
(3)
The first term shows the localized induction field arising from a curved line element acting on itself, and $`\mathrm{}_+`$ and $`\mathrm{}_{}`$ are the lengths of the two adjacent line elements that hold the point $`𝒔`$ between, and the prime denotes differentiation with respect to the arc length $`\xi `$. The mutual perpendicular vectors $`𝒔^{}`$, $`𝒔^{\prime \prime }`$ and $`𝒔^{}\times 𝒔^{\prime \prime }`$ point along the tangent, the principal normal and the binormal at the point $`𝒔`$, respectively, and their magnitudes are 1, $`R^1`$ and $`R^1`$, where $`R`$ is the local radius of curvature. The parameter $`a_0`$ is a cutoff parameter corresponding to a core radius. Thus the first term tends to move the local point $`𝒔`$ with a velocity inversely proportional to $`R`$, along the binormal direction. The second term represents the nonlocal field obtained by carrying out the integral of Eq. (2) along the rest of the filament. The approximation that describes the vortex dynamics neglecting the nonlocal terms and replacing Eq. (3) by
$$\dot{𝒔}=\beta 𝒔^{}\times 𝒔^{\prime \prime }$$
(4)
is called the localized induction approximation(LIA). Here the coefficient $`\beta `$ is defined by
$$\beta =\frac{\kappa }{4\pi }\mathrm{ln}\left(\frac{c<R>}{a_0}\right),$$
(5)
where $`c`$ is a constant of order 1 and $`(\mathrm{}_+\mathrm{}_{})^{1/2}`$ is replaced by the characteristic radius $`<R>`$.
When boundaries are present, the boundary-induced field $`𝒗_{s,b}`$ is added to $`𝒗_{s,\omega }`$ so that the boundary condition $`(𝒗_{s,\omega }+𝒗_{s,b})\widehat{𝒏}=0`$ can be satisfied. If the boundaries are specular plane surfaces, $`𝒗_{s,b}`$ is just the field by an image vortex made by reflecting the vortex into the plane and reversing its direction of the vorticity. Some other applied field $`𝒗_{s,a}`$, if present, is added, which results in the total velocity $`\dot{𝒔}_0`$ of the vortex filament without dissipation:
$`\dot{𝒔}_0=`$ $`{\displaystyle \frac{\kappa }{4\pi }}𝒔^{}\times 𝒔^{\prime \prime }\mathrm{ln}\left({\displaystyle \frac{2(\mathrm{}_+\mathrm{}_{})^{1/2}}{e^{1/4}a_0}}\right)`$ (6)
$`+`$ $`{\displaystyle \frac{\kappa }{4\pi }}{\displaystyle _{}^{^{}}}{\displaystyle \frac{(𝒔_1𝒓)\times d𝒔_1}{|𝒔_1𝒓|^3}}+𝒗_{s,b}(𝒔)+𝒗_{s,a}(𝒔).`$ (7)
At finite temperatures the mutual friction due to the interaction between the vortex core and the normal fluid flow $`𝒗_n`$ is taken into account. Then the velocity of a point $`𝒔`$ is given by
$$\dot{𝒔}=\dot{𝒔}_0+\alpha 𝒔^{}\times (𝒗_n\dot{𝒔}_0)\alpha ^{}𝒔^{}\times [𝒔^{}\times (𝒗_n\dot{𝒔}_0)],$$
(8)
where $`\alpha `$ and $`\alpha ^{}`$ are the temperature-dependent friction coefficients, and $`\dot{𝒔}_0`$ is calculated from Eq.(7). All calculations in this work are made for $`\alpha ^{}=0`$.
As discussed by Barenghi and Samuels , this formulation is essentially kinematic in the sense that the driving flows $`𝒗_n`$ and $`𝒗_{s,a}`$ are constant, that is, they only act on the vortex dynamics but are never affected by it. When the dynamics of the driving flows is concerned, it should be coupled selfconsistently to the vortex dynamics. However, since this work studies the system without the normal fluid and the driving superflow, this formulation will be useful to describe correctly the vortex dynamics, except for the phenomena that is concerned with the vortex core region, such as vortex reconnection, nucleation and annihilation.
Studying the vortex dynamics without the mutual friction needs to understand qualitatively the role of the mutual friction. Let us assume the LIA and neglect the term with $`\alpha ^{}`$. Then Eqs. (7) and (8) are reduced to
$$\dot{𝒔}=\beta 𝒔^{}\times 𝒔^{\prime \prime }+𝒗_{s,a}+\alpha 𝒔^{}\times (𝒗_n𝒗_{s,a}\beta 𝒔^{}\times 𝒔^{\prime \prime }).$$
(9)
If the mutual friction is absent, the dynamics due to only the self-induced velocity conserves the total line length of vortices. Under the above mutual friction, one can easily find that when the applied relative flow $`𝒗_n𝒗_{s,a}`$ blows against the local self-induced velocity $`\beta 𝒔^{}\times 𝒔^{\prime \prime }`$, the mutual friction always shrinks the vortex line locally. On the other hand, the relative flow along the self-induced velocity yields a critical radius of curvature
$$R_c\frac{\beta }{|𝒗_n𝒗_{s,a}|}.$$
(10)
When the local radius $`R`$ at a point on a vortex is smaller than $`R_c`$, the vortex will shrink locally, while the vortex of $`R>R_c`$ balloons out. Thus it should be noted that the mutual friction plays the dual role of the growth and decay of vortex line length. This dual role of the mutual friction sustains the steady state of the VT subject to the applied flow, where the highly curved structure whose local radius of curvature is less than $`R_c`$ will be smoothed out. If this applied field is absent, $`R_c`$ becomes infinite so that an arbitrary curved configuration of vortex lines shrinks away.
Here we will describe shortly the dynamical scaling discussed by Swanson and Donnelly, and Schwarz , which is necessary for understanding the cascade process of the VT dynamics. Using the LIA and absorbing the factor $`\beta `$ into reduced time $`t_0=\beta t`$ and velocity $`𝒗_0=𝒗/\beta `$, Eq. (8) becomes
$`{\displaystyle \frac{𝒔}{t_0}}=`$ $`𝒔^{}\times 𝒔^{\prime \prime }+𝒗_{s,0}+\alpha 𝒔^{}\times (𝒗_{n,0}𝒗_{s,0}𝒔^{}\times 𝒔^{\prime \prime })`$ (11)
$``$ $`\alpha ^{}𝒔^{}\times [𝒔^{}\times (𝒗_{n,0}𝒗_{s,0}𝒔^{}\times 𝒔^{\prime \prime })].`$ (12)
This equation is invariant under the scale transformation:
$`𝒔=\lambda 𝒔^{},\xi =\lambda \xi ^{},t_0=\lambda ^2t_0^{},`$ (13)
$`𝒗_{n,0}=\lambda ^1𝒗_{n,0}^{},𝒗_{s,0}=\lambda ^1𝒗_{s,0}^{}.`$ (14)
Accordingly, if all space coordinates of a system are reduced by a factor $`\lambda (<1)`$, the dynamics of the new system will look like the same as that of the old one, except that the velocity increases by $`\lambda ^1`$ and the time passes more rapidly by $`\lambda ^2`$. In other words, a small vortex loop whose configuration is similar to a large one but size is reduced by $`\lambda `$ follows the similar motion whose time scale shortens by $`\lambda ^2`$ compared with the large one.
Some important quantities which are useful for characterizing the VT will be introduced. The vortex line density(VLD) is
$$L=\frac{1}{\mathrm{\Omega }}_{}𝑑\xi ,$$
(15)
where the integral is made along all vortices in the sample volume $`\mathrm{\Omega }`$. Even though the VT may be homogeneous, it need not generally isotropic. The anisotropy of the VT which is made under the counterflow $`𝒗_{ns}`$ is represented by the dimensionless parameters
$`I_{}={\displaystyle \frac{1}{\mathrm{\Omega }L}}{\displaystyle _{}}[1(𝒔^{}\widehat{𝒓}_{})^2]𝑑\xi ,`$ (17)
$`I_{}={\displaystyle \frac{1}{\mathrm{\Omega }L}}{\displaystyle _{}}[1(𝒔^{}\widehat{𝒓}_{})^2]𝑑\xi ,`$ (18)
$`I_{\mathrm{}}\widehat{𝒓}_{}={\displaystyle \frac{1}{\mathrm{\Omega }L^{3/2}}}{\displaystyle _{}}𝒔^{}\times 𝒔\mathrm{"}𝑑\xi .`$ (19)
Here $`\widehat{𝒓}_{}`$ and $`\widehat{𝒓}_{}`$ stand for unit vectors parallel and perpendicular to the $`𝒗_{ns}`$ direction. The symmetry generally yields the relation $`I_{}/2+I_{}=1`$. If the VT is isotropic, the average of these measures are $`\overline{I}_{}=\overline{I}_{}=2/3`$ and $`\overline{I}_{\mathrm{}}=0`$.
The method of the numerical calculations is similar to that of Schwarz and described in our previous paper . A vortex filament is represented by a single string of points. The vortices configuration of a moment determines the velocity field in the fluid, thus moving the points on vortex filaments by Eqs. (7) and (8). Both local and nonlocal terms are represented by means of line elements connecting two adjacent points. As discussed in Ref. , the explicit forward integration of the local term may be numerically unstable. To prevent the difficulty, a modified hopscotch algorithm is adopted. As the vortex configuration develops and, particularly, two vortices approach each other, the length of a line element can change. Then it is necessary to add or remove points properly so that the local resolution does not lose(an adaptive meshing routine). Through the cascade process described in Sec. III, a large vortex can break up many times, eventually to a small one whose size is less than the space resolution, i.e., the distance between neighboring points on the filament. Of course the numerical calculation generally cannot follow the dynamics beyond its space resolution. Thus such vortices are eliminated numerically; the physical justification of this cut-off procedure will be discussed in Sec. IV.
How to deal with vortex reconnection is very important in the simulation of the VT. The numerical study of the incompressible Navie-Stokes fluid showed that the close interaction of two vortices leads to their reconnection, chiefly because of the viscous diffusion of the vorticity. Koplik and Levine solved directly the Gross-Pitaevskii equation to show the two close quantized vortices reconnect even in a inviscid fluid. Of course our numerical method for vortex filaments cannot represent the reconnection process itself. However Schwarz and the authors simulated the vortex dynamics near the reconnection using the full Biot-Savart law. When two vortices approach each other, let us define a critical distance
$$\mathrm{\Delta }2R/\mathrm{ln}(c<R>/a_0),$$
(20)
at which the nonlocal field from the other becomes comparable to its own local-induced field. Two vortices approaching within $`\mathrm{\Delta }`$ cause local twists on each other so that they become antiparallel at the closest place, even though they are not antiparallel initially. Then local cusps connecting these two develop, which will lead to reconnection. After the reconnection, two vortices run away rapidly from each other owing to their self-induced velocity. Considering both the full Biot-Savart calculation and the results of Ref.(), it will be reasonable to assume that two close filaments would reconnect. This assumption has an important meaning beyond a numerical expedient. The numerical simulation of the dense VT forces us to use the LIA, because the full Biot-Savart calculation requires much computing time. The LIA is expected to be a good approximation (to order 10$`\%`$) provided the inter-vortex spacing is enough large. However, when two vortices approach each other more closely than $`\mathrm{\Delta }`$, the nonlocal field becomes not negligible in reality. All the effects coming from such nonlocal field may be thought to be renormalized artificially by making the vortices reconnect. In the numerical simulation of the VT, Schwarz assumed that vortices which pass within $`\mathrm{\Delta }`$ are reconnected with unit probability. He noticed that the details of when and how the vortices are reconnected have no significant influence on the behavior of the VT, while the judgment by this $`\mathrm{\Delta }`$ can make unphysical reconnections. For example, two almost straight vortices must reconnect even if they are very apart, because their large radius $`R`$ of curvature results in the large $`\mathrm{\Delta }`$. The full Biot-Savart calculation shows that two vortices that once approach within $`\mathrm{\Delta }`$ can get away without reconnecting. Hence, in contrast to the method of Schwarz, this work reconnects the vortices which pass within not $`\mathrm{\Delta }`$ but the space resolution $`\mathrm{\Delta }\xi `$, for both the LIA and the full biot-Savart calculations. The concrete procedure is the following. Every vortex initially consists of a string of points at regular intervals of $`\mathrm{\Delta }\xi `$. The subsequent vortex motion can change the intervals of two adjacent points, yet the above adaptive meshing routine keeps each interval almost $`\mathrm{\Delta }\xi `$. When a point on a vortex approaches another point on another vortex more closely than the fixed space resolution $`\mathrm{\Delta }\xi `$, we join these two points and reconnect the vortices. Before and after the reconnection, the local line length may increase or decrease by a small quantity less than $`\mathrm{\Delta }\xi `$. This procedure is best for the filament reconnection under the full Biot-Savart calculation. The dependence of the LIA dynamics on $`\mathrm{\Delta }\xi `$ will be discussed in Sec. IV.
The numerical space resolution $`\mathrm{\Delta }\xi `$ and the time resolution $`\mathrm{\Delta }t`$ will be described for each calculation. For example, the dense tangle in a 1cm<sup>3</sup> cube shown in Fig. 9 (a) is calculated using $`\mathrm{\Delta }\xi =1.83\times 10^2`$cm, $`\mathrm{\Delta }t=1.0\times 10^3`$sec., $`N16,000`$points. Then, as described in Sec. IV, the VLD is conserved properly under the LIA, except for at each moment of reconnection.
## III DECAY OF DILUTE VORTICES
This section will investigate the dynamics of dilute vortices by the full Biot-Savart law described by Eq. (7).
We will begin with the collision of a straight vortex line and a moving ring in order to investigate what happens after the reconnection. Figure 1 shows the motion without the mutual friction. Toward the reconnection, the ring and the line twist themselves so that they become locally antiparallel at the closest place(Fig. 1 (a)). After the reconnection(Fig. 1 (b) and (c)), the resulting local cusps propagate along the vortices, exciting vortex waves. As shown in Fig. 2, the dynamics with the mutual friction($`\alpha =0.1`$) is similar, but there is a noticeable difference; the vortices are relatively smooth because of that smoothing effect of the mutual friction. For comparison, we calculated the dynamics under the LIA without the mutual friction. Although the twist due to the nonlocal interaction is absent, the behavior is similar to that of Fig. 1. It should be noted that the total line length under the LIA without the mutual friction is properly conserved within the numerical error except for at the moment of reconnection, while it is just lengthened by the nonlocal interaction in Fig. 1.
A typical scenario that vortex loop follows is shown in Fig. 3, which is a part of the process of Fig. 4. Two vortex loops approach each other to reconnect, thus becoming one loop. The reconnection excites vortex waves along the loop and makes it kinked. The kinked parts reconnect with the loop itself they belong to, thereby dividing into smaller loops. Then we are afraid that these kinks may arise from bad numerical methods, which can be denied by the following reasons. First, the calculation is made by enough mesh points even when there appear kinks. For example, even the left vortex in Fig. 3 (a) is represented by about 60 points. Secondly, as described in the last paragraph, we confirm that the total line length is conserved in the dynamics under the LIA without the mutual friction. Thirdly, a circular vortex ring is found to move at the expected speed without making kinks, which was proposed by Schwarz as a method that checks the numerical scheme.
Considering the above results, we will study the dynamics of dilute vortices with and without the mutual friction. The computation sample is taken to be a cube of size 1cm. The calculation is made by the space resolution $`\mathrm{\Delta }\xi =1.83\times 10^2`$cm and the time resolution $`\mathrm{\Delta }t=1.0\times 10^3`$sec. The initial configuration consists of four identical vortex rings placed symmetrically. We will study first the system subject to the periodic boundary conditions in all directions, that is, any vortex leaving the volume appears to reenter it from the opposite face, and next that surrounded by smooth, rigid walls.
Figure 4 shows the dynamics in the absence of the mutual friction. Four rings move toward the center of the cube by their self-induced velocity to make the first reconnection(a); the four rings resulting afterthat move outside oppositely(b). During the motion, they become kinked because of that mechanism described previously, and cut off their small kinked parts by reconnection. The periodic boundary conditions make the vortices collide repeatedly((c) and (d)), so that this self-similar process continues down to the scale of the space resolution below which the vortices are supposed to be eliminated numerically. This can be considered as the degenerate cascade process that follows the cascade decay process of the dense tangle investigated in the next section. Figure 5 shows the decay of the VLD $`L(t)`$ in the process of Fig. 4. When two vortices approach each other, the nonlocal interaction can stretch them, which sometimes causes just a little increase in $`L(t)`$. However the superior cascade process decreases the VLD as a whole. The effect of the mutual friction is shown in Fig. 6. The difference is apparent. The mutual friction smoothes and shrinks the vortex lines before lots of reconnection.
Figure 7 shows the dynamics with boundaries, starting from the same initial conditions. Although the early behavior (a) is similar to that of Fig. 4, all vortices collide with the boundaries and get attached there (b), afterthat behaving differently. Running along the walls (c) and colliding with the faces of the cube, they become kinked and broken up through the cascade process, ending in a degenerate state (d). As shown in Fig. 5, the VLD with the boundaries decays faster than that without boundaries. Under the periodic boundary conditions, the vortices collide only when they happen to meet each other in the volume. In the presence of solid boundaries, however, the vortex which runs along one boundary surface of the cube collides with its image vortex whenever it comes across another face. Thus the presence of the boundaries causes more reconnections and promotes the cascade process, which reduces VLD faster than the case of periodic boundary condition. We find that the system whose size of the cube is enlarged by a factor delays the decay of the VLD by the same factor, which supports strongly this scenario.
## IV DECAY OF THE VORTEX TANGLE
This section studies the free decay of the dense VT without mutual friction under the LIA. The decay of dilute vortices described in the last section follows this decay of the VT.
Throughout this section, the computation sample is taken to be a cube of size 1cm. The calculation is made by the space resolution $`\mathrm{\Delta }\xi =1.83\times 10^2`$cm and the time resolution $`\mathrm{\Delta }t=1.0\times 10^3`$sec. The one set of faces is subject to periodic boundary conditions. The other two sets of faces are treated as smooth, rigid boundaries, in which case vortices approaching the faces reconnect to them and their ends can move smoothly along the wall. The reason why we do not adopt the periodic boundary conditions in all directions is that then an artificial mixing process is necessary for obtaining an isotropic VT.
How to prepare the initial VT for free decay follows the method used by Schwarz. An initial state of six vortex rings is allowed to develop under a pure driving normal flow $`𝒗_n=v_n\widehat{𝒛}`$, where $`\widehat{𝒛}`$ is parallel to the direction along which the periodic boundary condition is used. This process should be made through the dynamics with the mutual friction($`\alpha =0.1`$), because the vortices free from the mutual friction never grow to a tangle as shown by Eq. (9). Although Schwarz continued the calculation until the vortices grew up to a steady self-sustaining state, we will take a growing VT at a moment to prepare a initial state for the simulation of the free decay. Figure 8(a) shows a example of the transient VT, which is anisotropic reflecting the anisotropy of the system. Turning off suddenly both the applied flow and the mutual friction transforms this VT into that of Fig. 8(b) after some time steps; this VT is nearly isotropic taking $`I_{}0.7`$; the little deviation from the isotropic value $`I_{}=2/3`$ may be attributed to the anisotropic boundary conditions.
The comparison of Fig. 8(a) and (b) shows a marked difference. The VT with the mutual friction consists of relatively smooth vortex lines, while the VT without it is very kinked owing to the lack of the smoothing effect of the mutual friction. Here it is necessary to check the accuracy of the numerical calculation. The LIA must conserve the VLD $`L(t)`$, whereas each numerical procedure of reconnection can change the local line length by a small quantity less than $`\mathrm{\Delta }\xi `$ before and after the event. We can monitor every reconnection in the VT dynamics, thus confirming that $`L(t)`$ is conserved completely within the numerical error except for at each moment of reconnection. Then we find that our calculation is enough accurate.
Figure 9 shows the decay of the VT without mutual friction. It is apparent that the tangle is becoming dilute. During this process, as shown in Fig. 10, $`L(t)`$ is actually reduced, with keeping the VT nearly isotropic with $`I_{}0.7`$. Since this system is free from the mutual friction, the only mechanism for the VT decay is that cut-off procedure which eliminates the small vortices whose size is less than the numerical space resolution. However it should be noted that the continuous reduction of $`L(t)`$ results in the presence of the stationary cascade process wherein large vortices break up to smaller ones through reconnections. This is because, if such cascade process is absent, even though the system is subject to that cut-off procedure, the VT only decays a little instantaneously and the continuous decay is never sustained. Only the cascade process that keeps supplying the small vortices can reduce the VT constantly.
Figure 11 compares the decay of $`L(t)`$ for the original space resolution $`\mathrm{\Delta }\xi `$ and its quarter $`\mathrm{\Delta }\xi /4`$; the latter calculation is made by the finer time resolution $`\mathrm{\Delta }t/16`$. The decay rate is found to be almost independent of the space resolution. Although more coarse space resolution would affect the decay rate, ours turn out to be enough fine to describe the cascade process.
What does this independence of the space resolution mean? If the original resolution $`\mathrm{\Delta }\xi `$ is improved to its quarter, the vortices of the size from $`\mathrm{\Delta }\xi `$ and to $`\mathrm{\Delta }\xi /4`$, which are supposed to vanish for the resolution $`\mathrm{\Delta }\xi `$, should still survive for the renewed one $`\mathrm{\Delta }\xi /4`$. Investigating the size distribution of vortices shows that the line length of the vortices of the size between $`\mathrm{\Delta }\xi `$ and $`\mathrm{\Delta }\xi /4`$ is not negligible compared with the total line length. Nevertheless the decay of $`L(t)`$ little depends on the space resolution, which is understood by the dynamical scaling described in Sec. II. A small vortex whose size is reduced by a factor $`\lambda `$ follows the dynamics whose time scale is shortened by $`\lambda ^2`$. Accordingly the small surviving vortices between $`\mathrm{\Delta }\xi `$ and $`\mathrm{\Delta }\xi /4`$ follow the rapid cascade dynamics to reach the cut-off scale $`\mathrm{\Delta }\xi /4`$, which proceeds much faster than the overall decay of $`L(t)`$ that includes the slow dynamics of large vortices too. Since it is difficult to improve the space resolution furthermore because of the computational constraints, we made the cut-off scale coarse oppositely keeping the spare resolution $`\mathrm{\Delta }\xi `$, in order to check how the decay rate is affected. When the cut-off scale is increased to$`2\mathrm{\Delta }\xi `$, $`3\mathrm{\Delta }\xi `$ and $`4\mathrm{\Delta }\xi `$, the decay rate of $`L(t)`$ is found to be almost the same as that with the cut-off scale $`\mathrm{\Delta }\xi `$, though more reduction of small vortices leads to larger fluctuation of $`L(t)`$. Accordingly, the decay rate is independent of the space resolution and the cut-off scale as far as we investigate in this work. This means that the overall decay rate of the VLD is determined principally by not small vortices but large ones whose size is comparable to the average line spacing.
It is important to know how this behavior depends on the scale of the system. Section II describes that the vortex dynamics under the LIA is subject to the dynamical scaling. Exactly speaking, this dynamical scaling is approximate, because the logarithmic term that depends on the characteristic radius $`<R>`$ through $`\beta `$ is neglected(Eq. (5)). The logarithmic dependence is so weak that the dynamical scaling is expected to be realized well, which should be confirmed numerically. We made the calculation for the systems with the different scaling factors $`\lambda =1,10^1,10^2`$. The dynamical scaling states that the VLD satisfies the relation $`L(\lambda )=\lambda ^2L(\lambda =1)`$, which was found to be well realized in the decay of the VT. Hence the VT dynamics is subject to the dynamical scaling within very high accuracy, thus being considered to be self-similar.
It is possible to classify the kinds of reconnection in the VT dynamics. The vortex reconnection is divided topologically into three classes, as shown in Fig. 12. The first refers to the process whereby two vortices reconnect to two vortices, which is most usual. The second is the process which divides one vortex into two vortices (the split type); the cascade process is driven by this kind of reconnection. Third is the process whereby two vortices are combined to one vortex against the cascade process (the combination type). Table 1 shows the number of reconnection events for each period in the VT dynamics of Fig. 9. The column ”total” refers to the total event number of all reconnections , and the columns ”split” and ”comb.” represent the event number of the above split and combination type, respectively. Most of reconnections belong to the first class. The reconnection of the second split type occupies about 17$`\%`$ of the total reconnections, being superior to that of the third combination type of about 10$`\%`$. It is found that the reconnection of the split type actually promotes the cascade process, against the reverse process due to that of the combination type.
The cascade process is revealed further by investigating the size distribution of vortices. Figure 13 shows the change of the size distribution in the VT dynamics of Fig. 9. Each figure shows the number $`n(x)`$ of vortices as a function of their length $`x`$. The system size $`a`$(=1cm) and the space resolution $`\mathrm{\Delta }\xi (=1.83\times 10^2`$cm), i.e., the cut-off length are the characteristic scales in this system. The vortices longer than $`a`$ are originally few, and most vortices are concentrated in the scale range $`[\mathrm{\Delta }\xi ,a]`$. As the cascade process progresses, every vortex generally divide into smaller ones through the split type reconnections, although some combination type reconnections may occur. As a result, the vortices larger than $`a`$ become fewer, and the vortices between $`\mathrm{\Delta }\xi `$ and $`a`$ are decreased in number too because they become smaller than $`\mathrm{\Delta }\xi `$ and be eliminated. However the contribution to the VLD is just different. Figure 14 shows the contribution to the VLD from the vortices in the size range $`[\mathrm{\Delta }\xi ,a]`$, $`[a,4a]`$, $`[4a]`$, respectively. The contribution from three ranges are comparable. The VLD of the large vortices fluctuates because they are few. The smooth VLD due to the vortices in the range $`[\mathrm{\Delta }\xi ,a]`$ seems to be similar to the overall $`L(t)`$ of Fig. 10. In the late stage($`t50`$s) of the dynamics, the large vortices become fewer, so that the contribution of the vortices between $`\mathrm{\Delta }\xi `$ and $`a`$ to the overall VLD is increased relatively.
The final destiny of small vortices through the cascade process may be interpreted several ways. First, the vortices whose size is eventually reduced to the order of the interatomic distance no longer sustain the vortex state, probably changing into such short-wavelength excitation as roton whose energy is comparable to that of the vortex. Secondly, the vortices can vanish at a small scale by radiating phonons, which is discussed recently by Vinen(See Sec. VI). Both mechanisms remove the small vortices from the system. Since both mechanisms work only at a small scale, some process that transfers energy from a large scale to smaller scales is necessary for the decay of the VT; this is just the cascade process. Thirdly, in a real system, the small vortices may collide with the vessel walls as studied in Sec. III. Since only the vortices in the bulk are observed experimentally, the reconnection with the walls may reduce the observed VLD effectively.
## V COMPARISON WITH THE VINEN’S EQUATION
This section compares our numerical results with the solution of the Vinen’s equation to show the good agreement between them.
The derivation of the Vinen’s equation will be reviewed briefly. Considering that cascade process at zero temperature proposed by Feynman , Vinen suggested that the homogeneous turbulence in the superflow without any normal fluid develops in a manner analogous to that of turbulence of high Reynolds number in an ordinary fluid. The vortices are supposed to be approximately evenly spaced with an average separation $`\mathrm{}=L^{1/2}`$. Then the energy of the vortices spreads from the eddies of wave number $`1/\mathrm{}`$ into a wide range of wave numbers, which means the self-similar VT sustained by the cascade process. The overall decay of the energy density will be governed by the chracteristic velocity $`v_s=\kappa /2\pi \mathrm{}`$ and the time constant $`\mathrm{}/v_s`$ of the eddies of the size $`\mathrm{}`$, so that
$$\frac{dv_s^2}{dt}=\chi _2\frac{v_s^2}{\mathrm{}/v_s}=\chi _2\frac{v_s^3}{\mathrm{}},$$
(21)
where $`\chi _2`$ is a parameter. Rewriting this by $`L`$, we obtain
$$\frac{dL}{dt}=\chi _2\frac{\kappa }{2\pi }L^2.$$
(22)
This is the Vinen’s equation that describes the decay of the VLD $`L(t)`$, and its solution is given by
$$\frac{1}{L}=\frac{1}{L_0}+\chi _2\frac{\kappa }{2\pi }t,$$
(23)
where $`L_0`$ is the VLD at $`t=0`$. At finite temperatures, the presence of the normal fluid may affect the cascade process. However, since the addition of the normal fluid introduces no new dimensional parameters into the vortex dynamics, the form of Eq.(22) cannot be altered and $`\chi _2`$ becomes a function of the temperature. The values of $`\chi _2`$ observed at finite temperatures are shown in Fig. 15. The symbols $``$ denotes the values observed when a heat current is suddenly switched on, while $`\mathrm{}`$ the values when a heat current is turned off. In any case, two kinds of $`\chi _2`$ reflects the complicated behavior of the normal fluid.
Figure 16 shows the comparison of our numerical results and the solution of the Vinen’s equation. The solid line refers to our result for the VT decay of Fig. 9, while three other lines denote Eq. (23) with the parameters $`\chi _2=0.5,0.3,0.2`$. Then we find that our result agrees excellently with the solution of $`\chi _2=0.3`$. There are two meanings for this. First, the decay of the numerical VT is well described by the Vinen’s equation. As stated in the last paragraph, the Vinen’s equation is based closely on the cascade process. Hence their agreement supports that the cascade process occurs really in the numerical simulation. Secondly, as seen from Fig. 15, the two kinds of data $``$ and $`\mathrm{}`$ are extrapolated towards zero temperature, then seeming to reach reasonably to $`\chi _20.3`$; the value obtained numerically may be consistent with those observed at finite temperatures.
In order to study how the mutual friction affects the cascade process, we calculate the decay of the VT with the mutual friction under the static normal fluid. As noted by Barenghi and Samuels , such phenomena might as well be calculated not kinematically but by a self-consistent approach which takes into account the back reaction of the VT onto the normal fluid. However, since the decay of an approximately isotropic and homogeneous VT may not induce some overall flow in the static normal fluid, this work, for simplicity, calculates kinematically the problem subject to the static normal fluid. Similar to the above calculation, we compare the numerical decay of the VT at finite temperatures with Eq. (23) with a fitting parameter $`\chi _2`$. The obtained dependence of $`\chi _2`$ on the mutual friction coefficient $`\alpha `$ is also shown in Fig. 15. When the temperatures are relatively low ($`T=`$0.91K, 1.07K and 1.26K), the solution with a proper value of $`\chi _2`$ can describe well the numerical result. However, as the temperature increases ($`T=`$1.6K), the numerical results become to deviate from Eq. (23). This seems to be reasonable. The decay term of the Vinen’s equation was derived originally based on the idea of the homogeneous turbulence. At low temperatures, the mutual friction is too small to disturb the inertial range, while the mutual friction at high temperatures shrinks not only small vortices but also large ones, thus disturbing the inertial range and deviating the numerical result from Eq. (23).
## VI CONCLUSIONS AND DISCUSSIONS
Motivated by the recent experimental work by Davis et.al. , we studied numerically the dynamics of the VT without the mutual friction. The absence of the mutual friction means that the usual well-known mechanism does not work for its free decay, so that we do not know why the VT decays. Throughout this paper, we conclude that the self-similar cascade process whereby large vortex loops break up to smaller ones proceeds in the VT, being closely concerned with the decay of the VT. This cascade process, which may be covered with the mutual friction at high temperatures, is just the one at zero temperature Feynman proposed , although the eventual destiny of the minimum vortex ring is beyond this formulation. The full Biot-Savart calculation is made for dilute vortices, while the LIA calculation for the dense VT. The former reveals the scenario: the reconnection of the vortices excites vortex waves on them and makes the vortex lines kinked, which would be suppressed in the presence of the mutual friction. The kinked parts reconnect with the body loop they belong to, breaking up to small loops. The LIA calculation shows that the cascade process proceeds in the VT, keeps making the small vortices below the space resolution and reduces the VLD $`L(t)`$. Although the small vortices below the space resolution are eliminated numerically, it should be emphasized that the VT never decays without the cascade process. The decay of $`L(t)`$ obtained numerically is consistent with the solution of the Vinen’s equation. The calculation that takes account of the mutual friction shows that both the modified cascade process and the vortex shrinkage due to the mutual friction proceeds together in the VT at a finite temperature.
Here we will describe the recent work by Vinen. In relation to the experimental work of the grid turbulence , Vinen discussed the dissipation of the VT at zero temperature. The dissipation can occur only by the emission of sound waves (phonons) by an oscillating vortex. The vortex oscillation of the average vortex spacing $`\mathrm{}=L^{1/2}`$ has the characteristic velocity $`v_{\mathrm{}}\kappa /\mathrm{}`$ and the characteristic time $`\tau _{\mathrm{}}\mathrm{}^2/\kappa `$. Estimating the dipole and quadrupole radiation from a Kelvin wave finds that such oscillation can cause only the very slow decay of the VT compared with $`\tau _{\mathrm{}}`$. Hence Vinen considered the excitation of the Kelvin wave whose wavelength is much smaller than $`\mathrm{}`$. In a classical viscous fluid, there is a flow of energy from components of the velocity field with small wave numbers to components with large wave numbers, energy being dissipated by viscosity near the Kolmogorov wave number. The superfluid system will have the energy cascade process of the Kelvin waves, whereby the energy is transformed to Kelvin waves with wave numbers greater than $`\mathrm{}^1`$ and eventually dissipated at a wave number $`\stackrel{~}{k}_2`$ by sound radiation. Based on this picture, Vinen reformulated the Vinen’s equation and obtained
$$\stackrel{~}{k}_2\mathrm{}=\left(\frac{C\mathrm{}}{A^{1/2}\kappa }\right)^{1/2}$$
(24)
for the case of dipole radiation, where $`C`$ is the speed of sound and $`A`$ is a constant. It should be noted that this Vinen’s Kelvin wave cascade process corresponds to our cascade process which is shown by the direct simulation of the vortex dynamics. The difference is that, although Vinen considered only the Kelvin wave, our cascade process includes not only the excitation of vortex waves but also the breakup of large loops to smaller ones through reconnection, which was assumed to be negligible by Vinen but is found to be present by our simulation. Whether the excitation of vortex waves or the breakup of vortex loops, the structure of small wave number will be produced continuously. We will estimate Eq. (24) for our simulation of the decay of the dense VT. As shown in Fig. 10, $`L`$ is supposed to be 400 cm<sup>-2</sup>, so that $`\mathrm{}=L^{1/2}=1/20`$cm. Taking $`C2\times 10^4`$cm/s for liquid helium and $`\kappa 10^3`$ cm<sup>2</sup>/s and assuming the unknown constant $`A`$ is the order of 1, Eq. (24) yields $`\stackrel{~}{k}_2\mathrm{}10^3`$, i.e., $`\stackrel{~}{k}_22\times 10^4`$ cm <sup>-1</sup>. Since the characteristic length $`\stackrel{~}{k}_2^15\times 10^5`$cm for sound radiation is enough smaller than our numerical space resolution $`\mathrm{\Delta }\xi `$, our cut-off procedure may be considered to be used for the effect of the sound radiation, assuming the cascade process continues self-similarly also from $`\mathrm{\Delta }\xi `$ to $`\stackrel{~}{k}_2^1`$.
We have to comment on how the nonlocal interaction acts on the VT. In a VT, the local field is usually superior to the nonlocal field. As stated in Sec. III, however, when two vortices approach each other, the nonlocal interaction can stretch them partly. The full Biot-Savart calculation in Sec. III shows that in dilute vortices the cascade process is superior to the stretch due to the nonlocal interaction. In a dense VT, these two processes can compete with each other; which is superior may depend on the VLD or the size distribution of vortices. Although the full Biot-Savart calculation for a dense VT is much CPU expensive and difficult, we start the calculation and obtain some preliminary results showing that the decay due to the cascade process still proceeds. The detail will be reported shortly.
Our results are compared with the recent experiment by Davis et.al. The observed $`T`$-independent decay below 70mK strongly suggests that the phonon gas plays no role, because the phonon density falls as $`T^3`$ in this range, and there must be an unknown intrinsic process in this superfluid system. We believe that our cascade process is closely connected with the $`T`$-independent decay. Davis et.al. observed the time costant of the decay was the order of 10 sec. The time constant depends on the amplitude of the VLD, but we do not know exactly the homogeneity of the VT and the amplitude of the VLD in the experiments. Accordingly it is difficult to compare our results quantitatively with the experimental data at present.
Such sound radiation can heat the fluid, which is recently discussed by Samuels and Barenghi. They estimated thermodynamically how much the temperature of the fluid increases when the kinetic energy of the VT is transformed to compressive energy, i.e., phonons. Since the traditional second-sound technique fails in the very low temperatures, the observation of the vortex heating is useful for investigating this system.
Nore et.al. studied the dynamics of the VT without any friction, by the direct numerical simulation of the Gross-Pitaevskii equation. They show that the total energy of the VT is partly transformed to compressive energy, and the energy spectrum can follow the Kolmogorov law. The dynamics they studied seems to include the cascade process of this work, but its detail is not clear.
Finally we will comment on the eddy viscosity. The superfluid turbulent state in a capillary flow induces excess temperature and pressure differences between both ends of the capillary, more than those in the laminar flow state. The excess temperature difference is understood by the mutual friction, while the excess pressure difference is described phenomenologically by the eddy viscosity. The eddy viscosity works for superfluid and reduces its total momentum, but its origin has not been necessarily revealed. The eddy viscosity which is thought to be an intrinsic mechanism in superfluid may be related with this cascade process.
###### Acknowledgements.
We acknowledge W.F. Vinen and P.V.E. McClintock for useful discussions. One of the author(S.N.) thanks Osaka City University(OCU) for giving an opportunity to visit OCU and Russian Foundation of Basic Research (grant N 99-02-16942) for supporting that field.
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# Charge occupancy of two interacting electrons on artificial molecules – exact results
## I Introduction
Much of the early work on quantum dots concentrated on large dots, which may contain many electrons, in the limit of very weak coupling between the dots and the leads, where one can employ the simple Coulomb blockade picture. In that limit, the energy cost of adding a single electron to the dot is of order of the excess energy for the added charge. For a total of two electrons, this is of the order of the Coulomb repulsion between the two electrons. One expects that for a good contact between the dots and the leads this energy cost will be reduced. However, the study of a general coupling strength presents a great challenge. The case of an ‘open’ dot, connected by quantum point contacts to a two-dimensional electron gas, has been analyzed using bozonization techniques and mapping of the Hamiltonian onto the two- and four-channel Kondo problem.
It has only recently become possible to also study small quantum dots, which have a small number of states and contain a small number of electrons. Such a small dot, connected to external leads, is similar to a donor in a doped semiconductor: both may be modeled as an ‘impurity’ connected to external leads. A set of such quantum dots, or an artificial molecule, can then be modeled by a set of such ‘impurities’. In what follows, we sometimes interchange the terms ‘dot’ and ‘impurity’. It is usually assumed that the electrons interact with each other only when they are on the same quantum dot, and behave as free electrons when they are on the leads. In what follows we therefore assume a contact interaction, which exists only on the dots, and present exact results for the case of two electrons. Given the difficulties in solving the general problem, such analytical results (even for the most simple configruations) are helpful. They are particularly useful for nanostructures, where one might design controlled experiments.
We have recently reported on several exact results for two interacting electrons on a general number of dots $`𝒩`$, which are modeled as ‘impurities’ which have single electronic states: we presented a general scheme for finding the eigenenergies, and presented some results for the spectra of a single dot. We have also discussed the exact two-electron current through a single dot. Here we generalize these results, with emphasis on the charge accumulated on each quantum dot and on its relationship with the Coulomb blockade picture. We then devote most of this paper to discuss the more complex case of a double dot. Such a double dot, with one state per dot, has recently been proposed as a possible candidate for the two-qubit entanglement required for quantum computation. The case of two coupled quantum dots is also amenable to experiments.
In our earlier work , we showed that the spectrum and the wave functions of the two interacting electrons can be obtained in terms of the energy spectrum and the wave functions of the single-electron Hamiltonian. We reproduce these results in section II in a slightly different method, and use them to obtain new results for the charge occupancies on the quantum dots. The next two sections are devoted to the study of specific configurations: a single dot, and a system made up of two dots, separated by a distance $`R`$. The single-electron spectra of these two configurations, required for the study of the two electron one, are discussed in the Appendix.
## II Two interacting electrons – general scheme
As has been demonstrated in Ref. , the knowledge of the spectrum of the single-electron Hamiltonian is sufficient for deducing the spectrum and the wave functions of two interacting electrons, for any number $`𝒩`$ of dots. Basically, we start with the Hamiltonian
$$=_{\mathrm{se}}+_{\mathrm{int}}.$$
(1)
The spin-independent single-electron part $`_{\mathrm{se}}`$ involves site energies $`ϵ_\mathrm{i}`$ on the dots ($`\mathrm{i}=1,2,\mathrm{},𝒩`$) and zero on the lead sites, and also nearest neighbor hopping matrix elements $`t_{n,m}`$ which assume special values near the dots. This part is diagonalized by the eignenergies $`\{ϵ_a\}`$ and the corresponding eigenfunctions $`\{\varphi _a(n)\}`$.
For simplicity, we assume that the two electrons interact only when they are both on the same dot i, with interaction energy $`U(\mathrm{i})`$ (though the method of solution can be extended for other types of interactions):
$$_{\mathrm{int}}=\underset{\mathrm{i}}{}U(\mathrm{i})c_\mathrm{i}^{}c_\mathrm{i}c_\mathrm{i}^{}c_\mathrm{i}.$$
(2)
Using the single-electron eigenstates, the two-electron Hamiltonian takes the form
$`={\displaystyle \underset{a\sigma }{}}ϵ_ac_{a\sigma }^{}c_{a\sigma }+{\displaystyle \underset{\mathrm{i}}{}}{\displaystyle \underset{abcd}{}}U_{acbd}(\mathrm{i})c_a^{}c_bc_c^{}c_d,`$ (3)
$`U_{acbd}(\mathrm{i})=U(\mathrm{i})\varphi _a^{}(\mathrm{i})\varphi _c^{}(\mathrm{i})\varphi _b(\mathrm{i})\varphi _d(\mathrm{i}).`$ (4)
Here, $`c_{a\sigma }^{}_n\varphi _a(n)c_{n\sigma }^{}`$ creates an electron in the state $`a`$ with spin $`\sigma `$.
For such a contact electron-electron interaction, one is interested only in the singlet state of the two electrons (the energies of the two electrons in the triplet state are simply given by the non-interacting sums $`ϵ_a+ϵ_b`$). We hence write for the two-electron singlet wave function
$`|\mathrm{\Psi }={\displaystyle \underset{ab}{}}X_{ab}(E)c_a^{}c_b^{}|0,`$ (5)
where $`|0`$ is the vacuum and $`X_{ab}=X_{ba}`$. The Schrödinger equation
$`|\mathrm{\Psi }=E|\mathrm{\Psi }`$ (6)
then yields
$`{\displaystyle \underset{ab}{}}\left(Eϵ_aϵ_b\right)X_{ab}(E)c_a^{}c_b^{}|0`$ (7)
$`=`$ $`{\displaystyle \underset{\mathrm{i}}{}}{\displaystyle \underset{aba^{}b^{}}{}}X_{ab}(E)U_{a^{}b^{}ab}(\mathrm{i})c_a^{}^{}c_b^{}^{}|0.`$ (8)
Multiplying this equation from the left by $`0|c_{b^{\prime \prime }}c_{a^{\prime \prime }}`$ gives
$`X_{ab}(E)={\displaystyle \underset{\mathrm{i}}{}}{\displaystyle \underset{a^{}b^{}}{}}{\displaystyle \frac{U_{aba^{}b^{}}(\mathrm{i})X_{a^{}b^{}}(E)}{Eϵ_aϵ_b}}.`$ (9)
The simple form of the matrix elements of the contact interaction \[see Eq. (4)\] allows us to rewrite Eq. (9) as a set of $`𝒩`$ linear equations: Defining the quantities
$`A_\mathrm{i}(E)={\displaystyle \underset{ab}{}}\varphi _a(\mathrm{i})\varphi _b(\mathrm{i})X_{ab}(E)`$ (10)
(which represent the amplitudes of $`|\mathrm{\Psi }`$ for the singlet state with both electrons on site i, denoted by $`|\mathrm{i},\mathrm{i}`$), one arrives at
$`A_\mathrm{i}(E)={\displaystyle \underset{\mathrm{i}^{}}{}}U(\mathrm{i}^{})G_E(\mathrm{i},\mathrm{i};\mathrm{i}^{},\mathrm{i}^{})A_\mathrm{i}^{}(E),`$ (11)
in which $`G_E(\mathrm{i},\mathrm{i};\mathrm{i}^{},\mathrm{i}^{})`$ is the two-particle Green’s function of two non-interacting electrons,
$`G_E(n_1,n_2;n_1^{},n_2^{})={\displaystyle \underset{ab}{}}{\displaystyle \frac{\varphi _a(n_1)\varphi _b(n_2)\varphi _a^{}(n_1^{})\varphi _b^{}(n_2^{})}{Eϵ_aϵ_b}},`$ (12)
calculated at the impurity locations. The determinant of Eqs. (11) gives the eigenenergies $`\{E\}`$ of the two interacting electrons, and in particular determines the ground state energy, $`E_G`$. The coefficients $`X(E)`$ are then obtained from Eq. (9), which can be rewritten as
$`X_{ab}(E)={\displaystyle \underset{\mathrm{i}}{}}{\displaystyle \frac{U(\mathrm{i})\varphi _a^{}(\mathrm{i})\varphi _b^{}(\mathrm{i})A_\mathrm{i}(E)}{Eϵ_aϵ_b}}.`$ (13)
Substituting this result into Eq. (5), it is easy to check that
$$|\mathrm{\Psi }=\underset{\mathrm{i}}{}U(\mathrm{i})A_\mathrm{i}(E)\underset{n_1,n_2}{}G_E(n_1,n_2;\mathrm{i},\mathrm{i})|n_1,n_2,$$
(14)
$`|n_1,n_2`$ is a singlet state with one electron at site $`n_1`$ and the other at site $`n_2`$. Note that Eq. (11) determines the $`A_\mathrm{i}`$’s only up to a multiplicative constant. This constant should be determined by the normalization of $`|\mathrm{\Psi }`$, i. e. from the condition $`_{ab}|X_{ab}(E)|^2=1`$. Such solutions will be discussed in some detail below.
The electronic states on a quantum dot are commonly probed by varying the gate voltages on the dots, represented here by the $`ϵ_\mathrm{i}`$’s, and measuring the conductance. Alternatively, one may probe the total charge on the dots, by measuring the capacitance when the gate voltage is changed. The total charge on the dots, in the state $`|\mathrm{\Psi }`$, is given by (in units of the electron charge, $`e`$)
$`P`$ $`=`$ $`{\displaystyle \underset{\sigma }{}}{\displaystyle \underset{\mathrm{i}=1}{\overset{𝒩}{}}}\mathrm{\Psi }|c_{\mathrm{i}\sigma }^{}c_{\mathrm{i}\sigma }|\mathrm{\Psi }`$ (16)
$`\mathrm{\Psi }|{\displaystyle \underset{\mathrm{i}=1}{\overset{𝒩}{}}}{\displaystyle \underset{ab}{}}\varphi _a^{}(\mathrm{i})\varphi _b(\mathrm{i}){\displaystyle \underset{\sigma }{}}c_{a\sigma }^{}c_{b\sigma }|\mathrm{\Psi }.`$
Using Eq. (5), we obtain
$`P=2{\displaystyle \underset{\mathrm{i}}{}}{\displaystyle \underset{abc}{}}\varphi _a^{}(\mathrm{i})\varphi _b(\mathrm{i})X_{ac}^{}(E)X_{bc}(E).`$ (17)
Alternatively, we note that
$$P=\underset{\mathrm{i}}{}\frac{E}{ϵ_\mathrm{i}}.$$
(18)
This follows from first-order perturbation theory: writing $`E=\mathrm{\Psi }||\mathrm{\Psi }`$, the derivative with respect to $`ϵ_\mathrm{i}`$ becomes $`\mathrm{\Psi }|c_{\mathrm{i}\sigma }^{}c_{\mathrm{i}\sigma }|\mathrm{\Psi }`$. In the following sections, we will use these results to discuss the ground state of simple quantum dot systems occupied by two electrons.
An important technical advantage of the representation of the two-electron spectrum in terms of the two-particle Green’s function of the non-interacting system, is that the latter can be expressed in terms of the single-particle Green’s function. The spectral representation of the single-particle Green’s function, $`g_\omega (n;n^{})`$, is
$`g_\omega (n;n^{})={\displaystyle \underset{a}{}}{\displaystyle \frac{\varphi _a(n)\varphi _a^{}(n^{})}{\omega +i\zeta ϵ_a}},`$ (19)
where $`\zeta 0^+`$. In conjunction with Eq. (12), one finds
$`G_E(n_1,n_2;n_1^{},n_2^{})`$ (20)
$`=`$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega g_{E\omega }(n_1;n_1^{})\mathrm{}g_\omega (n_2;n_2^{})`$ (21)
$`=`$ $`{\displaystyle \frac{i}{2\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\omega g_{E\omega }(n_1;n_1^{})g_\omega (n_2;n_2^{}),`$ (22)
where the last equality follows from the Kramers-Kronig relations. The relation (22) is very useful in the detailed calculations of the ground state properties, because the single-particle Green’s functions are relatively easy to find. We give in Appendix A the details of the single-particle Green’s functions required for the dot configurations investigated in this paper.
## III a single dot on a one-dimensional wire
We model a single dot on a one-dimensional wire by a single impurity, located at site 0, which has the on-site energy $`ϵ_0`$, and is coupled to two ideal one-dimensional leads, with the amplitude $`t_0`$ for tunneling between the impurity and its nearest neighbors on the leads. Both these parameters are experimentally accesssible: $`ϵ_0`$ models the plunger gate voltage on the dot, while $`t_0`$ is related to the transmittance (or conductance) of the barriers between each dot and the leads, which can be varied by changing the gate voltages on these barriers. The corresponding amplitudes between sites inside the leads are equal to $`t`$. Although some aspects of the solution of this problem were discussed in Ref. , we present here an alternative derivation, which is more adapted to the calculation of the occupation $`P`$ and to the more complicated case of two dots.
As shown in Appendix A, the single-particle Hamiltonian of such a system has none, one or two bound states, depending on the values of the “gate voltage” $`ϵ_0`$ and the “hybridization” $`t_0^2\gamma `$ (energies are measured in units of $`t`$). For the sake of concreteness, we concentrate on the region where there is only one bound state below the band, of energy $`ϵ_\beta <2`$. This occurs for $`ϵ_0<ϵ_{00}2(\gamma 1)`$ (see Appendix A).
In the simple case of a single ‘impurity’, Eq. (11) reduces to a single equation, and the eigenenergies $`\{E\}`$ of the two interacting electrons are given by the solutions of
$`{\displaystyle \frac{1}{U}}=G_E(0,0;0,0)={\displaystyle \underset{ab}{}}{\displaystyle \frac{|\varphi _a(0)|^2|\varphi _b(0)|^2}{Eϵ_aϵ_b}}.`$ (23)
It is easy to deduce the beahvior of $`G_E`$ as function of the two-electron energy $`E`$ from this equation. $`G_E`$ is negative for $`E<2ϵ_\beta `$, decreasing from 0 to $`\mathrm{}`$ as $`E`$ increases from $`\mathrm{}`$ towards $`2ϵ_\beta `$. As $`E`$ crosses this value, it jumps to $`+\mathrm{}`$ and then decreases. The value $`E=2+ϵ_\beta `$ marks the beginning of the two-particle continuous band states: one electron is bound and the other is in the continuum.
As discussed in Ref. , $`G_E`$ is finite at $`E=2+ϵ_\beta `$ in the thermodynamic limit of infinite leads, due to the vanishing of the band state wavefunction $`\varphi _k`$ (with energy $`ϵ_k=2\mathrm{cos}k`$) at the impurity site $`\mathrm{i}=0`$ for $`k=0`$. This value of $`G`$ determines whether there is or there is not a bound state of the two interacting electrons: When $`1/U<G_{2+ϵ_\beta }`$, then Eq. (23) has no solution for $`E<2+ϵ_\beta `$, and there is no doubly occupied bound state below the band. One of the electrons is then in a band state. The behavior of $`G`$ at $`E=2+ϵ_\beta `$, as function of $`ϵ_0`$, is plotted in Fig. 1. Since $`G_{E=2+ϵ_\beta }`$ has a maximum, $`G_{\mathrm{max}}`$, there is always a doubly occupied bound state (or an “insulator”) for $`U<1/G_{\mathrm{max}}`$. For larger $`U`$, the equation $`G_{2+ϵ_\beta }=1/U`$ has two solutions, $`ϵ_{0,}`$ and $`ϵ_{0,+}`$ (which depend on $`U`$). For $`ϵ_{0,}<ϵ_0<ϵ_{0,+}`$ one has no doubly occupied bound state, and the ground state of the two electrons lies in the continuum (i. e. represents a “metal”). Then, as the on-site energy $`ϵ_0`$ becomes more attractive, the bound state of the two electrons re-appears. As seen from Fig. 1, $`ϵ_{0,}`$ diverges to $`\mathrm{}`$ when $`U\mathrm{}`$, and this re-entrance then disappears. The region between $`ϵ_{0,}`$ and $`ϵ_{0,+}`$ becomes narrower as $`\gamma `$ increases, and for finite $`U`$ it always disappears above some critical hybridization $`\gamma _c`$ (which diverges to $`\mathrm{}`$ as $`U\mathrm{}`$).
These “insulator to metal” transitions of the two-electron ground state, from being bound to being in the continuum and back, are reflected in the occupancy $`P`$ of the dot in the ground state \[see Eq. (18)\]. To find $`P`$, we re-write the equation for the ground energy $`E_G`$ in the form
$$\frac{1}{U}=\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}d\omega (D_{E_G\omega })^1\mathrm{}(D_\omega )^1,$$
(24)
where we have used Eq. (22), and the results for $`g_\omega `$ relevant for this geometry \[Eq. (A10)\]. The imaginary part appearing in this expression has a delta-function contribution coming from the single-electron bound energy, and the contribution arising from the band states. Separating these two, we have
$`{\displaystyle \frac{1}{U}}`$ $`=`$ $`{\displaystyle \frac{r(ϵ_\beta )}{D_{E_Gϵ_\beta }}}`$ (25)
$``$ $`{\displaystyle \frac{1}{\pi }}{\displaystyle _2^2}d\omega (D_{E_G\omega })^1\mathrm{}(D_\omega )^1,`$ (26)
where $`r(\omega )=(D_\omega /\omega )^1`$ is the residue at the bound energy pole $`\omega =ϵ_\beta `$. It is now straightforward to differentiate all terms in (26) with respect to $`ϵ_0`$, and obtain $`P=E_G/ϵ_0`$.
We have solved Eq. (26) for $`E_G`$, by calculating the integral numerically. We have then computed the variation of the occupancy $`P`$ as function of the gate voltage $`ϵ_0`$, and the results are shown in Fig. 2, for a comparatively large hybridization, and in Fig. 3, for a small hybridization $`\gamma `$. Generally, $`P`$ starts at 0 at $`ϵ_0=ϵ_{00}=2(\gamma 1)`$, when the single electron bound state just moves below the band (with the inverse localization length $`\kappa _\beta =0`$, i. e. zero weight on the impurity). $`P`$ then grows as $`ϵ_0`$ decreases. As seen from Fig. 1, $`ϵ_{0,+}`$ is quite close to $`ϵ_{00}=2(\gamma 1)`$, due to the steepness of $`G`$ \[$`G`$ diverges to $`\mathrm{}`$ as $`ϵ_02(\gamma 1)`$\]. Therefore, the first transition into the “metallic” phase, at $`ϵ_{0,+}`$, occurs when the localization lengths of the two bound electrons are still quite large, and their weights on the impurity (and thus also $`P`$) are relatively small. The calculation of $`P`$ in this regime is not easy, due to numerical problems related to the above mentioned steepness. In any case, $`P`$ reaches values close to 1 somewhere inside the “metallic” phase, i. e. for $`ϵ_{0,}<ϵ_0<ϵ_{0,+}`$. We note that in the first “insulating” phase, which appears at $`ϵ_{0,+}<ϵ_0<2(\gamma 1)`$, both electrons are bound on very shallow states, hence the small value of $`P`$. Thus, it is not enough to know $`P`$ in order to determine the transport nature of the system. As $`ϵ_0`$ crosses below $`ϵ_{0,}`$, into the second “insulating” phase, $`P`$ gradually increases towards 2, reflecting the strongly bound state of the two electrons. This gradual increase becomes steeper as the hybridization $`\gamma `$ becomes smaller, and the width of the “metallic” single electron occupancy regime (of order $`ϵ_{0,+}ϵ_{0,}`$) increases with increasing $`U`$. Both of these facts are in qualitative accordance with the Coulomb blockade picture (where usually the derivative of $`P`$ with respect to the gate voltage has peaks whose width increases with the hybridization and whose inter-peak distance increases with $`U`$). In fact, the distance between the $`N`$’th and the $`(N1)`$’th peaks is usually interpreted as the energy cost of adding the $`N`$’th electron. However, it is usually very difficult to obtain quantitative estimates for these quantities in that picture. Furthermore, the similarity of our results to the simple Coulomb blockade picture is completely lost as $`\gamma `$ increases towards and beyond $`\gamma _c`$: the width of the “metallic” regime then shrinks, and the there is a continuous gradual increase of $`P`$ from 0 to 2.
Returning to Eq. (14), we now observe that for a single impurity, the two-electron state is given by
$$|\mathrm{\Psi }=UA\underset{n_1,n_2}{}G_E(n_1,n_2;0,0)|n_1,n_2,$$
(27)
where $`A`$ is found from the normalization. One can now use Eq. (22) and the single electron Green’s functions $`g_\omega (n;0)`$ to obtain $`|\mathrm{\Psi }`$. For the bound ground state, the results show an exponential decay of the amplitudes as either electron moves away from the impurity.
## IV Two dots on a one-dimensional wire
Two dots are modeled by two ‘impurities’, connected to each other and to the outside by ideal linear leads. The presence of two impurities gives rise to up to four single-particle bound states. For simplicity, we consider two identical impurities, each having the same on-site energy $`ϵ_0`$, which are located at sites $`\mathrm{}`$ and $`r`$, and are separated by a distance $`R`$ ($`R2`$). Confining ourselves again to the configuration where the bound states appear only below the band, the first bound state appears when $`ϵ_0`$ is smaller than $`ϵ_{00}=2(\gamma 1)`$, while the second appears only for $`R>R_c`$, where
$`R_c=2\gamma /(2\gamma 2ϵ_0).`$ (28)
At fixed $`R`$, there exists a single bound state below the band only in the narrow regime
$$2(\gamma 1\gamma /R)<ϵ_0<2(\gamma 1).$$
(29)
We also restrict ourselves to the regime with $`ϵ_0<2(1\gamma )`$, so that there are no bound states above the band (see Appendix A).
Assuming the on-site Coulomb interaction to be identical on the two impurities, $`U(\mathrm{})=U(r)U`$, Eqs. (11) yield
$`A_{\mathrm{}}(E)`$ $`=`$ $`UG_E(\mathrm{d})A_{\mathrm{}}(E)+UG_E(\mathrm{nd})A_r(E),`$ (30)
$`A_r(E)`$ $`=`$ $`UG_E(\mathrm{nd})A_{\mathrm{}}(E)+UG_E(\mathrm{d})A_r(E),`$ (31)
where the labels d and nd stand for the diagonal and the nondiagonal elements of the matrix. By symmetry,
$`G_E(\mathrm{d})`$ $``$ $`G_E(\mathrm{},\mathrm{};\mathrm{},\mathrm{})=G_E(r,r;r,r),`$ (32)
$`G_E(\mathrm{nd})`$ $``$ $`G_E(\mathrm{},\mathrm{};r,r)=G_E(r,r;\mathrm{},\mathrm{}).`$ (33)
The eigenenergies of the two interacting electrons are given by the solutions of the two equations
$`{\displaystyle \frac{1}{U}}=G_E(\mathrm{d})\pm G_E(\mathrm{nd})G_E^\pm ,`$ (34)
and the corresponding solutions obey
$$A_{\mathrm{}}^\pm (E)=\pm A_r^\pm (E).$$
(35)
It is instructive to rewrite these equations in terms of the single-electron wave functions, using Eq. (12). In the symmetric molecule case, one can divide the solutions into even and odd single-electron wave functions, with $`\varphi _a(\mathrm{})=\pm \varphi _a(r)`$. From Eq. (12) it now follows that
$`G_E^\pm ={\displaystyle \underset{ab}{}}{\displaystyle \frac{[\varphi _a(\mathrm{})\varphi _b(\mathrm{})\pm \varphi _a(r)\varphi _b(r)]\varphi _a^{}(\mathrm{})\varphi _b^{}(\mathrm{})}{Eϵ_aϵ_b}}.`$ (36)
Thus, it is clear that $`G_E^+`$ contains only pairs of states where both $`a`$ and $`b`$ are even or odd, while $`G_E^{}`$ contains only mixed combinations, where one state is even and the other is odd. It thus follows that the solutions of the two-electron problem divide into two separate families: the solutions of $`G_E^+=1/U`$ involve only even-even and odd-odd single electron states, while those of $`G_E^{}=1/U`$ involve only even-odd states: the coefficients in Eq. (5) will split into two separate families, associated with the different solutions of Eqs. (34). This can be easily seen by substituting Eq. (35) into Eq. (14).
To discuss the two-electron energies, we need to analyze the $`E`$-dependence of $`G_E^\pm `$. This depends on $`R`$: For $`R>R_c`$, there exist two single-electron bound states, the even $`\varphi _{\beta +}`$ and the odd $`\varphi _\beta `$. Thus, $`G_E^+`$ decreases from $`\mathrm{}`$ to $`\mathrm{}`$ as $`E`$ increases from $`2ϵ_{\beta +}`$ towards $`2ϵ_\beta `$. Therefore, in this case the equation $`G_E^+=1/U`$ always has a discrete solution, with $`E`$ between $`2ϵ_{\beta +}`$ and $`2ϵ_\beta `$. In the same case, $`G_E^{}`$ decreases from $`\mathrm{}`$ towards a finite value as $`E`$ increases from $`ϵ_{\beta +}+ϵ_\beta `$ towards the bottom of the continuum $`2+ϵ_{\beta +}`$ (which contains both even and odd states). The new lowest even-odd state may thus be either “insulating” or “metallic”, depending on the sign of $`1/UG_{2+ϵ_{\beta +}}^{}`$. However, the energy of this even-odd state is always above the lowest triplet energy, which is equal to the non-interacting value $`ϵ_{\beta +}+ϵ_\beta `$. From our numerical calculations we observe that $`G_E^+`$ is negative at this lowest triplet energy. Therefore, the lowest solution of $`G_E^+=1/U`$ is the ground state of the two-electron problem, which is thus a singlet. In a way, this might have been expected: breaking the system into two parts, by removing a bond in the middle between the two dots, one ends up with separate “atomic” states on each side, each coupled to its own lead. Each side can then contain one electron with either spin up or spin down. However, switching on the hopping $`t_h`$ between the two sides would lower the energy of the singlet state, similarly to the antiferromagnetic ground state of the Hubbard model; to lowest order in $`t_h`$, the “exchange” difference between the triplet and singlet states is of order $`t_h^2/U`$. It is interesting to note that in our case there exists a finite difference between the singlet and the triplet even in the limit $`U\mathrm{}`$, since we find that $`G_{ϵ_{\beta +}+ϵ_\beta }^+`$ is strictly negative. It would be interesting to study generalizations of our simple model, e. g. including interdot Coulomb and exchange interactions, which would allow an interchange of the singlet and triplet ground states.
The only chance to find a “metallic” ground state is thus for $`R<R_c`$, when there exists only one single-electron bound state below the band. This limits the possible range of parameters to that in Eq. (29). In this regime, $`G_E^{}`$ yields no doubly bound state, and $`G_E^+`$ yields one only if $`G_{2+ϵ_{\beta +}}^+<1/U`$. Note that this regime becomes narrower (in terms of $`ϵ_0`$) as $`R`$ increases. It is therefore interesting to find the borderline in the parameter space, at which $`G_{2+ϵ_{\beta +}}^+=1/U`$. Inside this broderline, the ground energy of the two electrons is in the continuum, i. e. “metallic”.
To calculate $`G_E^\pm `$, we use Eq. (22) and the results (A21) and (A24) of Appendix A:
$`G_E^+=G_E(\mathrm{d})+G_E(\mathrm{nd})`$ (37)
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑\omega \left(\frac{1}{D_{E\omega }^{}}\mathrm{}\frac{1}{D_\omega ^{}}+\frac{1}{D_{E\omega }^+}\mathrm{}\frac{1}{D_\omega ^+}\right)},`$ (39)
$`G_E^{}=G_E(\mathrm{d})G_E(\mathrm{nd})`$
$`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑\omega \left(\frac{1}{D_{E\omega }^{}}\mathrm{}\frac{1}{D_\omega ^+}+\frac{1}{D_{E\omega }^+}\mathrm{}\frac{1}{D_\omega ^{}}\right)},`$ (40)
where $`D_\omega ^{}`$ are given by Eqs. (A24).
We next separate the contributions of the bound energies from the integrals, to find
$`G_E^+={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{r^+(ϵ_{\beta +})}{D_{Eϵ_{\beta +}}^+}}+\mathrm{\Theta }(RR_c){\displaystyle \frac{r^{}(ϵ_\beta )}{D_{Eϵ_\beta }^{}}}\right)`$ (41)
$`{\displaystyle \frac{1}{2\pi }}{\displaystyle _2^2}𝑑\omega \left({\displaystyle \frac{1}{D_{E\omega }^{}}}\mathrm{}{\displaystyle \frac{1}{D_\omega ^{}}}+{\displaystyle \frac{1}{D_{E\omega }^+}}\mathrm{}{\displaystyle \frac{1}{D_\omega ^+}}\right),`$ (42)
$`G_E^{}={\displaystyle \frac{1}{2}}\left({\displaystyle \frac{r^+(ϵ_{\beta +})}{D_{Eϵ_{\beta +}}^{}}}+\mathrm{\Theta }(RR_c){\displaystyle \frac{r^{}(ϵ_\beta )}{D_{Eϵ_\beta }^+}}\right)`$ (43)
$`{\displaystyle \frac{1}{2\pi }}{\displaystyle _2^2}𝑑\omega \left({\displaystyle \frac{1}{D_{E\omega }^{}}}\mathrm{}{\displaystyle \frac{1}{D_\omega ^+}}+{\displaystyle \frac{1}{D_{E\omega }^+}}\mathrm{}{\displaystyle \frac{1}{D_\omega ^{}}}\right),`$ (44)
where $`r^\pm (\omega )=(D_\omega ^\pm /\omega )^1`$ are the residues at the poles.
We have used Eq. (44) to solve the equation $`G_{2+ϵ_\beta }^+=0`$, which yields the borderline of the “metallic” regime in the limit $`U\mathrm{}`$. The result for $`R=2`$ is depicted by the dotted line in Fig. 4. The area enclosed inside this line represents the “metal”, where $`G_{2+ϵ_{\beta +}}^+>0`$. For smaller $`U`$ and for larger $`R`$’s this area shrinks further.
Figure 4 highlights a major difference between the single-dot and the double-dot cases. In the former, the width of the “metallic” regime (in terms of $`ϵ_0`$) was equal to $`ϵ_{0,+}ϵ_{0,}`$, and at fixed $`\gamma `$ it increased with $`U`$, diverging to $`\mathrm{}`$ for $`U\mathrm{}`$. Although this width was not equal to $`U`$, as assumed in the simple Coulomb blockade picture, it still resembled the qualitative features of that picture. In contrast, in the double-dot case this width is bounded by Eq. (29), and this bound is independent of $`U`$. Therefore, the width of the singly occupied “metallic” regime remains finite and small even when $`U\mathrm{}`$. Basically, this happens because in the double-dot case, there exist two single-electron bound states. The level-repulsion between these bound states then prevents the two-electron ground state from merging into the coninuum. In fact, we expect similar bounds on the Coulomb-blockade-like energy even for a single quantum dot, whenever the dot has more than a single bound state.
As $`R`$ increases, the two single-electron ground energies of the double dot become closer to each other and to the single-dot bound state energy. Since $`2ϵ_{\beta +}<E_G<2ϵ_\beta `$, it follows that the “insulating” ground energy of the two interacting electrons is almost independent of $`U`$. Moreover, as $`|ϵ_0|`$ increases, the two single-particle bound energies approach one another, and practically we have only one, doubly-degenerate, single-particle bound energy. This behavior is shown in Fig. 5, for a very ‘open’ dot. Similar effects arise when the hybridization is reduced: $`ϵ_{\beta +}`$ and $`ϵ_\beta `$ also become indistinguishable. It hence follows that independently of $`U`$, the ground state energy of the two interacting electrons is $`E_G2ϵ_{\beta +}2ϵ_\beta `$. In such a situation, the charge accumulated on the quantum dot, $`P`$, will just follow the weight of the single-particle localized wave functions on the impurities. These have a ‘smooth’ behavior as function of $`ϵ_0`$ (see Fig. 6). Consequently, the “Coulomb blockade” type behavior, which is obtained for the single-impurity dot (Figs. 2 and 3) is washed out.
Finally, we comment on the two-electron wave function in the ground state. At small $`R`$, when there exists only one single-electron bound state, this wave function is dominated by the even-even state in which both electrons occupy the state $`\varphi _{\beta +}`$. To leading order in $`\gamma `$ this term represents the Heitler-London molecular state, proportional to $`|\mathrm{},\mathrm{}+|r,r+|\mathrm{},r+|r,\mathrm{}`$, with equal weights to the electrons being on the same impurity or each electron being on a different impurity. For large $`R`$, $`E_G`$ is close to both $`2ϵ_{\beta +}`$ and to $`2ϵ_\beta `$, and therefore the coefficients $`X_{ab}(E_G)`$ will be dominated by the terms with $`a,b=\beta +,\beta +`$ or $`a,b=\beta +,\beta `$, with roughly equal magnitudes for these two terms \[see Eq. (13)\]. Combining these two terms, it is easy to see that $`|\mathrm{\Psi }`$ is then dominated by the atomic orbitals $`|\mathrm{},r+|r,\mathrm{}`$. Thus, our exact results interpolate nicely between these two leading approximations, which are common in chemistry textbooks.
## V conclusions
Our exact solution of the two-electron problem does not have the Fermi gases on the leads; it is limited to the ‘canonical ensemble’, with a fixed small number of electrons in the system. We do believe however that this exact solution for this simple case is still useful, in that it throws light on issues which sometimes remain unclear in approximate (and much more complicated) treatments of the real problem.
Our solution may also be directly connected with e. g. the ionization of donors into the conduction band, as function of the system parameters. Simple models for this problem may involve $`𝒩`$ coupled one electron donors, or a single dielectronic donor. Our calculation shows that as the distance between a pair of single-electron donors decreases, then the on-site interaction $`U`$ helps this ionization. This effect may well be an important ingredient for the real metal-insulator transition in some semiconductors. Specifically, if each donor in the semiconductor has exactly one electron attached to it, then as the density of donors increases we expect pair of donors to combine into “molecules” which allow only one bound electron. The remaining electrons will move to the band, and the system will become “metallic”.
The fact that even an onsite $`U`$ can have effects which are so different from the naive Coulomb blockade model, should also be of interest. This is especially so for the double quatum dot case. As mentioned in the introduction, the spectra of such systems are in principle addressable by transport and capacitance experiments. Usually, one looks at the contribution of the resonant states (lying in the continuum) of these dots to the conductivity, but the bound states will also contribute to the off-resonance transmission. The dependence of the average occupancy on the gate voltage is of interest both theoretically and experimentally.
Generalizations of this treatment to more realistic situations, even for two electrons, are relatively easy to achieve. For example, the inclusion of an interdot interaction $`V`$ does not affect the need to solve only $`𝒩`$ linear equations. As stated, such interactions may cause an interchange of the singlet and triplet ground states , and will certainly affect the magnetic exchange interactions between the electrons on different ‘impurities’.
###### Acknowledgements.
This paper is dedicated to Franz Wegner, on the occasion of his 60th birthday. All authors acknowledge the hospitality of the Centre for Advanced Studies of the Norwegian Academy, where parts of this work were done. This research is also supported by grants from the Israeli Science Foundation and from the Israeli Ministry of Science and the French Ministry of Research and Technology.
## A The single-particle Green’s function
As discussed in the text, the two-particle Green’s function can be expressed in terms of the single-particle one. Here we derive the latter.
The single-electron tight-binding type Hamiltonian, $`_{\mathrm{se}}`$, is given by
$`_{\mathrm{se}}`$ $`=`$ $`{\displaystyle \underset{n}{}}ϵ_nc_n^{}c_n`$ (A1)
$`+`$ $`{\displaystyle \underset{n}{}}(t_{n,n+1}c_n^{}c_{n+1}+t_{n1,n}c_n^{}c_{n1}),`$ (A2)
where we ignore the spin indices, since the single particle Hamiltonian is spin-independent. Writing the single-particle Green’s function, $`g`$, in the form
$`g(n,n^{};t)=i\mathrm{\Theta }(t)[c_n^{}(t),c_n^{}]_+,`$ (A3)
we find, for the Fourier transform $`g_\omega `$, the equation
$`\omega g_\omega (n,n^{})=\delta _{n,n^{}}+ϵ_ng_\omega (n,n^{})`$ (A4)
$`+`$ $`t_{n,n+1}g_\omega (n+1,n^{})+t_{n1,n}g_\omega (n1,n^{}).`$ (A5)
This equation is straightforwardly solved for the configurations described in the text.
### 1 The single impurity case
For a single impurity with one-dimensional leads we have $`ϵ_n=0`$, for $`n0`$, and $`ϵ_0`$ is the on-site energy of the impurity, $`t_{n,n\pm 1}t`$ for $`n0,\pm 1`$, and $`t_{0,\pm 1}t_0`$. It is sufficient to consider Eq. (A5) for $`n^{}=0`$. Then for any $`n0`$ or $`\pm 1`$ that equation gives
$`\omega g_\omega (\pm n,0)=tg_\omega (\pm (n+1),0)tg_\omega (\pm (n1),0).`$ (A6)
It is easy to convince oneself that, in the limit of infinite leads, the Green’s function does not depend on the details of the boundary conditions. For $`n0`$ one can therefore assume the solution
$$g_\omega (n,0)=C_\omega a_\omega ^{|n|},$$
(A7)
and find from Eq. (A6) that
$`{\displaystyle \frac{\omega }{t}}`$ $`=`$ $`a_\omega 1/a_\omega .`$ (A8)
Thus, for $`|\omega /t|<2`$ we can denote $`a_\omega =e^{ik_\omega }`$, with $`\omega /t=2\mathrm{cos}k_\omega `$. For $`\omega /t)\stackrel{>}{<}\pm 2`$ we denote $`a_\omega =\pm e^{\kappa _\omega }`$, and $`\omega /t=2\mathrm{cosh}\kappa _\omega `$.
The equation for $`n=\pm 1`$ now yields $`C_\omega =(t_0/t)g_\omega (0,0)`$, and finally the equation for $`n=0`$ yields
$`g_\omega (0,0)`$ $`=`$ $`1/D_\omega ,`$ (A9)
$`D_\omega `$ $`=`$ $`\omega ϵ_0+2(t_0^2/t)a_\omega .`$ (A10)
From Eq. (19), the poles of $`g_\omega (0,0)`$ give the eigenvalues of the single electron problem, while the corresponding residues give the probability that an electron in a given state in on the impurity. The boundary between having or not having a bound state is easily found by setting $`a_\omega =\pm 1`$ and $`\omega /t=2`$ in the equation $`D_\omega =0`$ (with the upper sign refering to a state below the band). Measuring energies in units of $`t`$, and denoting
$`\gamma =(t_0/t)^2,`$ (A11)
we find that for $`ϵ_0<2(\gamma 1)`$ there exists a bound state below the band (with $`0<a_\omega =e^{\kappa _\beta }1`$), with a localization length $`1/\kappa _\beta `$ \[which diverges to $`\mathrm{}`$ at $`ϵ_0=2(\gamma 1)`$\] and energy $`ϵ_\beta =2\mathrm{c}\mathrm{o}\mathrm{s}\mathrm{h}(\kappa _\beta )`$. For $`ϵ_0>2(1\gamma )`$ there appears a bound state above the band, with a localization length $`1/\kappa _\alpha `$ and energy $`ϵ_\alpha =2\mathrm{c}\mathrm{o}\mathrm{s}\mathrm{h}(\kappa _\alpha )`$. The two localization lengths are given by
$`e^{\kappa _{\alpha ,\beta }}=\pm {\displaystyle \frac{ϵ_0}{2}}+\sqrt{\left({\displaystyle \frac{ϵ_0}{2}}\right)^21+2\gamma }.`$ (A12)
The weights of the localized wave function on the impurity (i.e., the residues of $`g_\omega `$ at the bound energies) are accordingly
$`|\varphi _{\alpha ,\beta }(0)|^2{\displaystyle \frac{ϵ_{\alpha ,\beta }}{ϵ_0}}=\left[1+{\displaystyle \frac{2\gamma }{e^{2\kappa _{\alpha ,\beta }}1}}\right]^1.`$ (A13)
A similar analysis gives the band of extended states, with energies $`ϵ_k=2\mathrm{cos}k`$.
### 2 The two impurity case
Here we consider a system with two impurities, which are separated by a distance $`R`$ ($`R2`$). We denote the locations of the two impurities by $`\mathrm{}`$ and $`r`$, and assume $`t_{n,n\pm 1}t`$ for $`n\mathrm{}`$ or $`r`$, $`t_{n,n\pm 1}t_0`$ for $`n=\mathrm{}`$ or $`r`$, $`ϵ_n=0`$ for $`n\mathrm{}`$ or $`r`$, and $`ϵ_{\mathrm{},r}ϵ_0`$.
Again, it is sufficient to consider $`g_\omega (n,n^{})`$ with $`n^{}=\mathrm{},r`$. Referring to Eq. (A5), we assume a solution of the form
$`g_\omega (n,\mathrm{})`$ $`=`$ $`C_\omega ^<(\mathrm{})a_\omega ^{|n\mathrm{}|},n<\mathrm{},`$ (A14)
$`g_\omega (n,\mathrm{})`$ $`=`$ $`C_\omega ^>(\mathrm{})a_\omega ^{|nr|},n>r,`$ (A15)
$`g_\omega (n,\mathrm{})`$ $`=`$ $`A_\omega (\mathrm{})a_\omega ^n+B_\omega (\mathrm{})a_\omega ^n,\mathrm{}<n<r,`$ (A16)
where $`a_\omega `$ is the solution of Eq. (A8). Writing Eq. (A5) for $`n=\mathrm{}1`$ and $`n=r+1`$ gives
$`C_\omega ^<(\mathrm{})`$ $`=`$ $`{\displaystyle \frac{t_0}{t}}g_\omega (\mathrm{},\mathrm{}),C_\omega ^>(\mathrm{})={\displaystyle \frac{t_0}{t}}g_\omega (r,\mathrm{}).`$ (A17)
The other two coefficients, $`A_\omega (\mathrm{})`$ and $`B_\omega (\mathrm{})`$ are found by using the equation for $`n=\mathrm{}+1`$ and $`n=r1`$. Then
$`A_\omega (\mathrm{})`$ $`=`$ $`{\displaystyle \frac{t_0}{t}}{\displaystyle \frac{g_\omega (\mathrm{},\mathrm{})a_\omega ^rg_\omega (r,\mathrm{})a_\omega ^{\mathrm{}}}{a_\omega ^Ra_\omega ^R}},`$ (A18)
$`B_\omega (\mathrm{})`$ $`=`$ $`{\displaystyle \frac{t_0}{t}}{\displaystyle \frac{g_\omega (r,\mathrm{})a_\omega ^{\mathrm{}}g_\omega (\mathrm{},\mathrm{})a_\omega ^r}{a_\omega ^Ra_\omega ^R}}.`$ (A19)
Finally we write Eq. (A5) for $`n=\mathrm{},r`$ and obtain
$`g_\omega (r,r)=g_\omega (\mathrm{},\mathrm{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{D_\omega ^+}}+{\displaystyle \frac{1}{D_\omega ^{}}}\right),`$ (A20)
$`g_\omega (\mathrm{},r)=g_\omega (r,\mathrm{})`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{D_\omega ^+}}{\displaystyle \frac{1}{D_\omega ^{}}}\right),`$ (A21)
where
$`D_\omega ^{}`$ $`=`$ $`\omega ϵ_0+{\displaystyle \frac{t_0^2}{t}}a_\omega `$ (A22)
$`+`$ $`{\displaystyle \frac{t_0^2}{t}}{\displaystyle \frac{a_\omega ^{R+1}a_\omega ^{R1}}{a_\omega ^Ra_\omega ^R}}\pm {\displaystyle \frac{t_0^2}{t}}{\displaystyle \frac{a_\omega a_\omega ^1}{a_\omega ^Ra_\omega ^R}}`$ (A23)
$`=`$ $`\omega ϵ_0+\gamma \left(a_\omega +{\displaystyle \frac{a_\omega ^{\frac{R}{2}1}a_\omega ^{1\frac{R}{2}}}{a_\omega ^{\frac{R}{2}}a_\omega ^{\frac{R}{2}}}}\right).`$ (A24)
The single-particle bound energies are determined by the poles of the Green’s functions, i.e., when $`D_\omega ^\pm `$ vanishes. Let us for simplicity confine ourselves to bound states below the band, with $`0<a_\omega 1`$. Then $`D_\omega ^+`$ produces a bound state with energy $`ϵ_{\beta +}`$ as long as $`ϵ_0<2(\gamma 1)`$, as is the case for the single impurity configuration. However, the second bound state, $`ϵ_\beta `$, coming from $`D_\omega ^{}`$, appears only at more negative $`ϵ_0`$, or (for fixed $`ϵ_0`$) when the distance between the impurities, $`R`$, is large enough: $`ϵ_0<2(\gamma 1\gamma /R)`$ (solve $`D_\omega ^{}=0`$ with $`a_\omega 1`$ and $`\omega =2`$). As $`R`$ tends to $`\mathrm{}`$, the two bound energies are approaching the same value, that of the bound energy of the single impurity system. We exemplify this behavior in Fig. 7, for $`\gamma =0.4`$ and $`ϵ_0=1.5`$; the state with the higher energy appears only for $`R>R_c\gamma /(\gamma 1ϵ_0/2)=2.666\mathrm{}`$. The calculations presented in this paper are also restricted to $`ϵ_0<2(1\gamma )`$, so that there exist no bound states above the band.
Generally, all the eigenstates of the problem divide into two subsets. Those which arise from $`D_\omega ^\pm =0`$ obey $`g_\omega (\mathrm{},r)/g_\omega (r,r)=\pm 1`$, and hence represent even (odd) solutions which obey
$$g_\omega ^\pm (\mathrm{}m,\mathrm{})=\pm g_\omega ^\pm (r+m,\mathrm{}).$$
(A25)
In particular, one can associate these subsets with the even and odd single-electron wave functions, $`\varphi _a(\mathrm{}m)=\pm \varphi _a(r+m)`$. The two bound states below the band, $`\varphi _{\beta \pm }(n)`$, thus correspond to the “bonding” and “antibonding” states of molecular chemistry. The antibonding energy $`ϵ_\beta `$ joins the band for inter-impurity distances below $`R_c`$.
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# 1 Introduction
## 1 Introduction
Study of string theory on $`AdS_3`$ background has been a subject of great importance mainly for the following two reasons: Firstly, it is a non-trivial example of completely solvable string theories on a curved background with the Lorentzian signature . The second reason, which is comparably newer than the first one, is the possibility of understanding the $`AdS/CFT`$-duality at a stringy level .
Although the string theory on this background (with NS $`B`$-field) is believed to be described by a simple conformal field theory, $`SL(2;\text{R})`$-WZW model, there are still subtleties. Especially two different theoretical grounds have been proposed with respect to the set up of the Hilbert space of quantum states;
1. The Hilbert space is defined to be a representation space of current algebra $`\widehat{SL}(2;\text{R})`$.
2. The Hilbert space is defined to be the Fock space of free fields which should be identified with the string coordinates in some suitable parametrizations of the group manifold $`SL(2;\text{R})`$.
The first ground is based on the standard prescription of two dimensional conformal field theory ($`CFT_2`$). It is well-known that the WZW model for a compact group actually has such a Hilbert space: the physical Hilbert space should be made up of a finite number of the unitary (integrable) representations of current algebra. Since we now have a non-compact group $`SL(2;\text{R})`$, the situation becomes rather non-trivial. If $`kk^{}2`$ ($`k^{}`$ means the level of $`\widehat{SL}(2;\text{R})`$) is equal to a negative rational value, we have a finite number of the “admissible” representations<sup>1</sup><sup>1</sup>1The admissible representations are not necessarily unitary representations of $`\widehat{SL}(2;\text{R})`$ with $`k=q/p`$. But these define the good-natured conformal blocks in the similar manner as the familiar $`(p,q)`$ minimal model with $`c=1{\displaystyle \frac{6(pq)^2}{pq}}`$ . which contain a rich structure with many singular vectors . However, string theory on $`AdS_3`$ background corresponds to the cases of positive $`k`$, in which we cannot have any unitary representation of $`\widehat{SL}(2;\text{R})`$ (except for the trivial representation). There are infinite number of non-unitary representations some of which have a few singular vectors (generically, no singular vectors). The best we can expect is that the BRST condition successfully eliminates all negative norm states from the physical Hilbert space. Many discussions have been given concerning the no-ghost theorem along this line . With respect to the discrete series the no-ghost theorem was proved under the assumption of a truncation of quantum number “$`j`$” parametrizing the second Casimir .
In this traditional approach to $`CFT_2`$ from the representation theory of current algebra, primary states are naively characterized by the two quantum numbers $`j`$, $`m`$. However, as was carefully discussed by Bars , there is a subtle point if we recall that we are working in a $`\sigma `$-model with a three dimensional non-compact target space. One should keep it in mind that under the limit of weak curvature; $`k+\mathrm{}`$, $`AdS_3`$ space (the universal covering of $`AdS_3`$, strictly speaking) may be replaced by a flat background $`\text{R}^{2,1}`$. In this sense it may be more natural that we have three conserved momenta characterizing the physical states, and the second definition of Hilbert space is likely to be more appropriate, as was claimed in the works .
More recently, based on the representation theory of $`\widehat{SL}(2;\text{R})`$, Maldacena and Ooguri claimed that one should take the Hilbert space enlarged by the spectral flow and proved the extended no-ghost theorem. On physical ground it may be plausible that the third momentum we mentioned above is related to such an enlargement of Hilbert space. One of the main purposes of this paper is to clarify the relation between them, namely, the role of spectral flow in the argument of and that of the third momentum of zero-modes (discrete light-cone momentum in the context of this paper).
Another aim of this work is to manifest further the role of “space-time Virasoro algebra” introduced in . It is inspired by the asymptotic isometries of Brown-Henneaux and is understood to describe the conformal symmetry of the long string sector . The generators of this algebra are most conveniently realized as operators acting on the Fock space of free fields (the Wakimoto’s $`\phi `$, $`\beta `$, $`\gamma `$ system in the usual treatment). This is one of the reasons why we take the free field realization rather than the abstract representation theory of current algebra.
This paper is organized as follows;
In section 2 we start with reformulating the bosonic and superstring theories on $`AdS_3`$ background by a free field realization. With the help of some field redefinitions we show that the $`AdS_3`$ string theory can be described by a string theory on a linear dilaton background (along the transverse direction) and with a light-like compactification. We further demonstrate that the space-time conformal algebra given in has a quite simple form analogous to the DDF operators in this framework of “discrete light-cone Liouville theory”.
In section 3 we analyze the physical spectrum in our framework. As that of a free field system we will reproduce the spectrum proposed in : Only the principal series is allowed due to the unitarity. We also comment on the outcome of turning on the Liouville potential term, which should be a marginal perturbation. Such an interaction term breaks the translational invariance along the radial direction (in other words, makes the “screening” of the extra zero-mode momentum from the view points of the $`SL(2;\text{R})`$ current algebra), and hence the “bound string states” possessing the imaginary momentum along this direction may appear in the physical spectrum. The space-time Virasoro operators play a role as the spectrum generating algebra, and we will observe that one must include the representations of $`\widehat{SL}(2;\text{R})`$ which are broader than those given in in order to make the full space-time Virasoro algebra act successfully on the physical Hilbert space.
In section 4 we further investigate the spectrum for superstring cases that give rise to space-time $`N=2,4`$ SCFTs. We present the complete set of on-shell chiral primaries. We will find that there are infinite number of on-shell chiral primary states with the different light-cone momenta, and the spectral flows act naturally among them. They become the chiral primaries also in the sense of space-time and have the space-time $`U(1)_R`$ charges in agreement with the expectation of $`AdS_3/CFT_2`$-duality. As a byproduct we also clarify the relationship with the symmetric orbifold CFT describing the multiple long strings discussed in .
We will summarize the main results of our analyses and present some discussions in section 5.
## 2 Reformulation of $`AdS_3`$ String Theory as the Discrete Light-cone Liouville Theory
### 2.1 Bosonic String on $`AdS_3\times 𝒩`$
Through this paper we shall consider the universal covering of the $`AdS_3`$ space with the Lorentzian signature so that the time direction is non-compact. We start with the following world-sheet Lagrangian for the (quantum) bosonic string on $`AdS_3`$ with $`NS`$ B-filed
$$=\phi \overline{}\phi \sqrt{\frac{2}{k}}R^{(2)}\phi +\beta \overline{}\gamma +\overline{\beta }\overline{\gamma }\beta \overline{\beta }\mathrm{exp}\left(\sqrt{\frac{2}{k}}\phi \right),$$
(2.1)
where $`R^{(2)}`$ denotes the curvature on the world-sheet. Throughout this paper we shall only focus on the physics at the near boundary region $`\phi +\mathrm{}`$, in which we can consider the interaction term (“screening charge term”) $`\beta \overline{\beta }\mathrm{exp}\left(\sqrt{{\displaystyle \frac{2}{k}}}\phi \right)`$ as a small perturbation. By dropping this term simply we obtain the free conformal field theory
$$T=\frac{1}{2}\phi \phi \frac{1}{\sqrt{2k}}^2\phi +\beta \gamma ,$$
(2.2)
$$\phi (z)\phi (0)\mathrm{ln}z,\beta (z)\gamma (0)\frac{1}{z}.$$
(2.3)
We will later discuss the effect of restoring this interaction term $`\beta \overline{\beta }\mathrm{exp}\left(\sqrt{{\displaystyle \frac{2}{k}}}\phi \right)`$ on the physical spectrum.
The $`SL(2;\text{R})`$ symmetry in this free system is described by the Wakimoto representation
$`j^{}`$ $`=`$ $`\beta `$
$`j^3`$ $`=`$ $`\beta \gamma +\sqrt{{\displaystyle \frac{k}{2}}}\phi `$ (2.4)
$`j^+`$ $`=`$ $`\beta \gamma ^2+\sqrt{2k}\gamma \phi +(k+2)\gamma ,`$
which generates the $`\widehat{SL}(2;\text{R})`$ current algebra of level $`k+2`$
$$\{\begin{array}{ccc}j^3(z)j^3(0)\hfill & \hfill & \frac{(k+2)/2}{z^2}\hfill \\ j^3(z)j^\pm (0)\hfill & \hfill & \frac{\pm 1}{z}j^\pm (0)\hfill \\ j^+(z)j^{}(0)\hfill & \hfill & \frac{k+2}{z^2}\frac{2}{z}j^3(0).\hfill \end{array}$$
(2.5)
By using the standard bosonization of $`\beta ,\gamma `$
$$\beta =ive^{uiv},\gamma =e^{u+iv},u(z)u(0)\mathrm{ln}z,v(z)v(0)\mathrm{ln}z,$$
(2.6)
we can rewrite the currents (2.4)
$`j^{}`$ $`=`$ $`ive^{uiv}`$
$`j^3`$ $`=`$ $`u+\sqrt{{\displaystyle \frac{k}{2}}}\phi `$ (2.7)
$`j^+`$ $`=`$ $`e^{u+iv}(k(u+iv)+\sqrt{2k}\phi +iv).`$
Moreover it is convenient to introduce the following new variables;
$`Y^0`$ $`:=`$ $`\sqrt{{\displaystyle \frac{2}{k+2}}}iu\sqrt{{\displaystyle \frac{k}{k+2}}}i\phi `$
$`Y^1`$ $`:=`$ $`\sqrt{{\displaystyle \frac{k+2}{2}}}v+{\displaystyle \frac{k}{\sqrt{2(k+2)}}}iu+\sqrt{{\displaystyle \frac{k}{k+2}}}i\phi `$ (2.8)
$`\rho `$ $`:=`$ $`\sqrt{{\displaystyle \frac{k}{2}}}(u+iv)+\phi .`$
Since this field redefinition is an $`SO(2,1)`$-rotation, we simply have
$$Y^0(z)Y^0(0)\mathrm{ln}z,Y^1(z)Y^1(0)\mathrm{ln}z,\rho (z)\rho (0)\mathrm{ln}z,$$
(2.9)
and any other combinations have no singular OPEs.
In these variables the $`\widehat{SL}(2;\text{R})_{k+2}`$ currents are given by
$`j^3`$ $`=`$ $`\sqrt{{\displaystyle \frac{k+2}{2}}}iY^0`$
$`j^\pm `$ $`=`$ $`\left(\sqrt{{\displaystyle \frac{k+2}{2}}}iY^1\pm \sqrt{{\displaystyle \frac{k}{2}}}\rho \right)e^{\sqrt{\frac{2}{k+2}}i(Y^0+Y^1)},`$ (2.10)
and also the stress tensor is rewritten as
$$T=\frac{1}{2}(Y^0)^2\frac{1}{2}(Y^1)^2\frac{1}{2}(\rho )^2\frac{1}{\sqrt{2k}}^2\rho ,$$
(2.11)
which of course has the correct central charge $`c=3+{\displaystyle \frac{6}{k}}`$.
In this way we have found that the bosonic string theory on $`AdS_3`$ can be realized by two free bosons $`Y^0`$, $`Y^1`$ (with no background charge) and a “Liouville mode” $`\rho `$ with the background charge $`Q\sqrt{{\displaystyle \frac{2}{k}}}`$. The essentially same realizations of $`\widehat{SL}(2;\text{R})`$ were used in several works . It was suggested in that the fields $`Y^0,Y^1,\rho `$ roughly corresponds to the global coordinates of $`AdS_3`$ space, and in the similar sense one might suppose that the Wakimoto fields $`\phi ,\beta ,\gamma `$ correspond to the Poincaré coordinates.
The definitions of Hermitian conjugations of these free fields are standard
$$(Y^i(z))^{}=Y^i(1/z),(\rho (z))^{}=\rho (1/z)Q\mathrm{ln}z.$$
(2.12)
(Recall that the Liouville field $`\rho `$ has a background charge.) One can easily verify that the Hermitian conjugations of the current generators take the usual forms $`j_n^3=j_n^3`$, $`j_n^\pm =j_n^{}`$ under (2.12). This is an advantage of this free field realization (2.10) compared with the Wakimoto representation in which the rules of Hermitian conjugation are not simple.
Notice also that the OPE of $`Y^0`$ with itself has the wrong sign (2.9) and thus $`Y^0`$ should correspond to the time-like coordinate. This fact is consistent with the usual interpretation; $`j_0^3\text{energy}`$.
There is a subtle point with respect to the realizations of currents (2.10): $`j^\pm (z)`$ are not necessarily defined as local operators on the whole Fock space of $`Y^0`$, $`Y^1`$, $`\rho `$. To overcome this difficulty it is natural to assume the following light-like compactification. Let us introduce the light-cone coordinates
$$Y^\pm =\frac{1}{\sqrt{2}}(Y^0\pm Y^1),$$
(2.13)
and assume the periodic identification
$$Y^{}Y^{}+2\pi R,R=\frac{2}{\sqrt{k+2}}.$$
(2.14)
Such a prescription is known as the name “discrete light-cone quantization”, and have been applied to the studies of M(atrix) theory with finite $`N`$ .
In this case, the conjugate momentum of $`Y^{}`$ is quantized as
$$\frac{Y^+}{\tau }P^++\overline{P}^+=\frac{2n}{R},n𝐙,$$
(2.15)
and the winding mode of $`Y^{}`$ should take
$$\frac{Y^{}}{\sigma }P^{}\overline{P}^{}=mR,m𝐙.$$
(2.16)
Since $`Y^+`$ remains non-compact, there is no winding mode along this direction
$$\frac{Y^+}{\sigma }P^+\overline{P}^+=0,$$
(2.17)
and thus we obtain $`P^+=\overline{P}^+={\displaystyle \frac{n}{R}}`$ $`(n\text{Z})`$.
By using these facts, we can concretely write the “tachyon” vertex operators $`𝒱_{j,m,\overline{m},p}(z,\overline{z})V_{j,m,p}(z)\overline{V}_{j,\overline{m},p}(\overline{z})`$, where the left mover is defined by
$$V_{j,m,p}=e^{\left(\frac{\sqrt{k+2}}{2}p\frac{2m}{\sqrt{k+2}}\right)iY^++\frac{\sqrt{k+2}}{2}piY^{}\sqrt{\frac{2}{k}}j\rho },$$
(2.18)
and the winding condition (2.16) means that $`m\overline{m}\text{Z}`$. The corresponding Fock vacuum $`|j,m,p`$ has the following properties;
$`j_0^3|j,m,p`$ $`=`$ $`(m{\displaystyle \frac{k+2}{2}}p)|j,m,p,`$ (2.19)
$`j_p^\pm |j,m,p`$ $`=`$ $`(m\pm j)|j,m\pm 1,p,`$ (2.20)
$`j_{p+n}^\pm |j,m,p`$ $`=`$ $`0,(n1),`$ (2.21)
$`L_0|j,m,p`$ $`=`$ $`\left({\displaystyle \frac{1}{k}}j(j1)+mp{\displaystyle \frac{k+2}{4}}p^2\right)|j,m,p.`$ (2.22)
Namely, $`j`$, $`m`$ mean the quantum numbers appearing in the usual $`\widehat{SL}(2;\text{R})`$ theory and $`p`$ corresponds to the label of “flowed representation” of . In fact, the spectral flow in the context of is defined as the following transformations
$$\{\begin{array}{c}j^3(z)j^3(z)+\frac{k+2}{2}\frac{p}{z}\hfill \\ j^\pm (z)z^pj^\pm (z).\hfill \end{array}$$
(2.23)
In the system of $`Y^0`$, $`Y^1`$, $`\rho `$, this is simply the momentum shift
$$Y^0Y^0p\sqrt{\frac{k+2}{2}}i\mathrm{ln}z,$$
(2.24)
and $`Y^1`$, $`\rho `$ remain unchanged.
The global $`SL(2;\text{R})`$ algebra $`\{j_0^3,j_0^\pm \}`$ is manifestly BRST invariant. We can immediately extend this algebra to the “space-time Virasoro algebra”
$$\{\begin{array}{ccc}_0\hfill & =\hfill & j_0^3\sqrt{\frac{k+2}{2}}iY^0\hfill \\ _n\hfill & =\hfill & \left(\sqrt{\frac{k+2}{2}}iY^1n\sqrt{\frac{k}{2}}\rho \right)e^{\frac{2n}{\sqrt{k+2}}iY^+}(n0),\hfill \end{array}$$
(2.25)
which actually generates the Virasoro algebra on the Fock space over the vacuum $`|j,m,p`$
$$[_n,_m]=(nm)_{n+m}+\frac{c}{12}(n^3n)\delta _{n+m,0},$$
(2.26)
where $`c=6(k+2)p`$.
We here make a few comments: Firstly, the Virasoro operators $`_n`$ are well-defined as local operators on the whole Fock space compatible with the light-like compactification (2.14). Secondly, it is natural to regard $`_n`$ ($`n0`$) as analogs of the DDF operators along the $`\rho `$-direction. In fact, $`_n`$ ($`n0`$) is no other than the unique solution for the BRST condition among the operators having the form $`{\displaystyle (A\rho +BY^++CY^{})e^{\frac{2n}{\sqrt{k+2}}iY^+}},A0`$ (up to BRST exact terms and some overall constant, of course). In the next section we will make use of such DDF like operators in order to construct the complete set of the physical states. As the last comment, we should point out that $`_n`$ are BRST equivalent to the space-time Virasoro operators constructed in . It is straightforward to confirm that the quantum number $`p`$ precisely coincides with the “winding of $`\gamma `$”; $`{\displaystyle \gamma ^1𝑑\gamma }=p`$.
### 2.2 Superstring on $`AdS_3\times S^1\times 𝒩/U(1)`$
Let us try to extend our previous results to the superstring cases. We start with the general superstring vacua $`AdS_3\times S^1\times 𝒩/U(1)`$ studied in , which are compatible with the world-sheet $`N=2`$ SUSY. The most familiar example $`AdS_3\times S^3\times M^4`$ ($`M^4=T^4`$ or $`K3`$) is nothing but a special example of these backgrounds, and we can readily apply the results in this subsection to that case too.
First of all, to fix the notations we summarize the world-sheet properties of this superstring model;
* $`AdS_3`$ sector ($`j^A,\psi ^A`$)
To extend to the superstring case, we introduce free fermions in the adjoint representation
$`\psi ^3(z)\psi ^3(0){\displaystyle \frac{1}{z}}`$ , $`\psi ^+(z)\psi ^{}(0){\displaystyle \frac{2}{z}},`$
$`\psi ^\pm `$ $`=`$ $`\psi ^1\pm i\psi ^2.`$ (2.27)
The total $`\widehat{SL}(2;\text{R})`$ currents are given by
$$J^A=j^A+j_F^A=j^A\frac{i}{2}ϵ_{BC}^A\psi ^B\psi ^C,A,B,C=1,2,3,$$
(2.28)
where the fermionic currents $`j_F^A`$ have the level $`2`$. The fermionic currents $`j_F^A`$ can be written by free fermions as
$$j_F^\pm =\pm \psi ^\pm \psi ^3,j_F^3=\frac{1}{2}\psi ^+\psi ^{}.$$
(2.29)
This sector has an $`N=1`$ superconformal symmetry given by
$`G_{SL(2,𝐑)}`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{k}}}({\displaystyle \frac{1}{2}}\psi ^+j^{}+{\displaystyle \frac{1}{2}}\psi ^{}j^+\psi ^3j^3{\displaystyle \frac{1}{2}}\psi ^+\psi ^{}\psi ^3),`$ (2.30)
and the central charge is
$$c=\frac{3(k+2)}{k}+\frac{3}{2}.$$
(2.31)
* $`S^1`$ sector ($`Y,\chi `$)
We have a scalar field $`Y`$ parametrizing $`S^1`$
$$Y(z)Y(0)\mathrm{ln}z,$$
(2.32)
and its fermionic partner $`\chi `$
$$\chi (z)\chi (0)\frac{1}{z}.$$
(2.33)
This sector has the simplest $`N=1`$ superconformal symmetry
$$G_{S^1}=\chi iY,$$
(2.34)
with the central charge
$$c=\frac{3}{2}.$$
(2.35)
* $`𝒩/U(1)`$ sector
We require that this sector has an $`N=2`$ superconformal symmetry described by the currents
$$T_{𝒩/U(1)},G_{𝒩/U(1)}^\pm ,J_{𝒩/U(1)}.$$
(2.36)
The relation between $`N=2`$ and $`N=1`$ superconformal current is
$$G_{𝒩/U(1)}=G_{𝒩/U(1)}^++G_{𝒩/U(1)}^{}.$$
(2.37)
Because of the criticality condition, the central charge of this sector should be equal to
$$c=9\frac{6}{k},$$
(2.38)
and the $`U(1)_R`$ current satisfies
$$J_{𝒩/U(1)}(z)J_{𝒩/U(1)}(0)\frac{3\frac{2}{k}}{z^2}.$$
(2.39)
We can realize the $`N=2`$ superconformal symmetry on the world-sheet in this system. We choose the $`U(1)_R`$ current as
$$J_R=J_{R1}+J_{R2}+J_{𝒩/U(1)},$$
(2.40)
where
$`J_{R1}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\psi ^+\psi ^{}+{\displaystyle \frac{2}{k}}J^3`$ (2.41)
$`J_{R2}`$ $`=`$ $`\chi \psi ^3.`$ (2.42)
According to the charge of this current the $`N=1`$ superconformal current splits into two terms
$`G`$ $`=`$ $`G_{SL(2,𝐑)}+G_{S^1}+G_{𝒩/U(1)}^++G_{𝒩/U(1)}^{}`$ (2.43)
$``$ $`G^++G^{},`$
where
$`G^\pm `$ $`=`$ $`G_1^\pm +G_2^\pm +G_{𝒩/U(1)}^\pm `$ (2.44)
$`G_1^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2k}}}\psi ^\pm j^{}`$ (2.45)
$`G_2^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\chi \psi ^3)\left({\displaystyle \frac{1}{\sqrt{2}}}iY\pm {\displaystyle \frac{1}{\sqrt{k}}}J^3\right).`$ (2.46)
The energy-momentum tensor is also decomposed as follows;
$`T`$ $`=`$ $`T_1+T_2+T_{𝒩/U(1)}`$ (2.47)
$`T_1`$ $`=`$ $`{\displaystyle \frac{1}{k}}(j^Aj_A+J^3J^3){\displaystyle \frac{1}{4}}(\psi ^+\psi ^{}\psi ^+\psi ^{})`$ (2.48)
$`T_2`$ $`=`$ $`{\displaystyle \frac{1}{k}}J^3J^3{\displaystyle \frac{1}{2}}(Y)^2{\displaystyle \frac{1}{2}}\chi \chi +{\displaystyle \frac{1}{2}}\psi ^3\psi ^3.`$ (2.49)
It may be worthwhile to mention that the superconformal generators $`\{T_i,G_i^\pm ,J_{Ri}\}`$ (anti-) commute among the different sectors. Furthermore $`\{T_1,G_1^\pm ,J_{R1}\}`$ has the same expression as that of the Kazama-Suzuki coset model for $`SL(2;\text{R})/U(1)`$ .
The BRST charge $`Q_{BRST}`$ is defined in the standard manner
$$Q_{BRST}=\left[c\left(T\frac{1}{2}(\varphi )^2^2\varphi \eta \xi +cb\right)+\eta e^\varphi Gb\eta \eta e^{2\varphi }\right],$$
(2.50)
where $`\varphi `$, $`\eta `$, $`\xi `$ are the familiar bosonized superghosts .
Now let us try to reformulate this superstring model as the discrete light-cone Liouville theory as in the case of bosonic string. Our goal is the $`N=2`$ Liouville theory with the light-like compactification; $`\text{R}_+\times S_{}^1\times \text{R}_\rho \times S^1\times 𝒩/U(1)`$. To this aim we need to perform further field redefinitions.
As a preliminary we bosonize the fermions $`\psi ^\pm `$
$$\psi ^\pm =\sqrt{2}e^{\pm iH_1},$$
(2.51)
where $`H_1(z)H_1(0)\mathrm{ln}z`$, and the radius of compact boson $`H_1`$ should be 1. Let $`Y_0,Y_1`$ be as given in (2.8), and define
$`X^0`$ $`:=`$ $`\sqrt{{\displaystyle \frac{k+2}{k}}}Y^0+\sqrt{{\displaystyle \frac{2}{k}}}H_1`$
$`X^1`$ $`:=`$ $`{\displaystyle \frac{2}{\sqrt{k(k+2)}}}Y^0+\sqrt{{\displaystyle \frac{k}{k+2}}}Y^1\sqrt{{\displaystyle \frac{2}{k}}}H_1`$ (2.52)
$`H_1^{}`$ $`:=`$ $`\sqrt{{\displaystyle \frac{2}{k+2}}}(Y^0+Y^1)+H_1.`$
Since this is again an $`SO(2,1)`$ rotation, we have the OPEs
$$X^0(z)X^0(0)\mathrm{ln}z,X^1(z)X^1(0)\mathrm{ln}z,H_1^{}(z)H_1^{}(0)\mathrm{ln}z,$$
(2.53)
and all the non-diagonal OPEs vanish. We also rewrite
$$X^2:=Y,\mathrm{\Psi }^2:=\chi ,\mathrm{\Psi }^0:=\psi ^3,\mathrm{\Psi }^\pm (\frac{1}{\sqrt{2}}(\mathrm{\Psi }^1\pm i\mathrm{\Psi }^\rho )):=e^{\pm iH_1^{}}.$$
(2.54)
After all, we have changed the system of
$$\{\phi ,\beta ,\gamma ,Y,\psi ^\pm ,\psi ^3,\chi \}$$
(2.55)
into the system of new free fields
$$\{\rho ,X^0,X^1,X^2,\mathrm{\Psi }^\pm ,\mathrm{\Psi }^0,\mathrm{\Psi }^2\}.$$
(2.56)
In these new variables the energy-momentum tensors (2.48), (2.49) are rewritten as
$`T_1`$ $`=`$ $`{\displaystyle \frac{1}{2}}(X^1)^2{\displaystyle \frac{1}{2}}(\rho )^2{\displaystyle \frac{1}{\sqrt{2k}}}^2\rho {\displaystyle \frac{1}{2}}(\mathrm{\Psi }^+\mathrm{\Psi }^{}\mathrm{\Psi }^+\mathrm{\Psi }^{})`$ (2.57)
$`T_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}(X^0)^2{\displaystyle \frac{1}{2}}(X^2)^2+{\displaystyle \frac{1}{2}}\mathrm{\Psi }^0\mathrm{\Psi }^0{\displaystyle \frac{1}{2}}\mathrm{\Psi }^2\mathrm{\Psi }^2.`$ (2.58)
The $`U(1)_R`$ currents (2.41), (2.42) become
$`J_{R1}`$ $`=`$ $`\mathrm{\Psi }^+\mathrm{\Psi }^{}QiX^1`$ (2.59)
$`J_{R2}`$ $`=`$ $`\mathrm{\Psi }^0\mathrm{\Psi }^2.`$ (2.60)
$`Q`$ is the background charge of Liouville mode $`\rho `$ and in this case $`Q=\sqrt{{\displaystyle \frac{2}{k}}}`$.
The $`N=2`$ superconformal currents (2.45), (2.46) now take the following forms
$`G_1^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\mathrm{\Psi }^\pm (iX^1\pm \rho ){\displaystyle \frac{Q}{\sqrt{2}}}\mathrm{\Psi }^\pm `$ (2.61)
$`G_2^\pm `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(\mathrm{\Psi }^0\mathrm{\Psi }^2)\times {\displaystyle \frac{1}{\sqrt{2}}}i(X^0\pm X^2).`$ (2.62)
It is also convenient to rewrite the total super current
$$G=\mathrm{\Psi }^0iX^0+\mathrm{\Psi }^1iX^1+\mathrm{\Psi }^2iX^2+\mathrm{\Psi }^\rho i\rho +Qi\mathrm{\Psi }^\rho .$$
(2.63)
Notice that $`\{T_1,G_1^\pm ,J_{R1}\}`$ (2.48), (2.45), (2.41) have been now transformed into the expressions of superconformal algebra in the $`N=2`$ Liouville theory as we mentioned before. The essential part of this field redefinition is the identification between $`SL(2;\text{R})/U(1)`$ Kazama-Suzuki model and the $`N=2`$ Liouville theory (see the appendix B of , and also refer ) and it was claimed in that these two theories are related by a T-duality.
As in the bosonic case, we introduce the light-cone coordinates
$$X^\pm =\frac{1}{\sqrt{2}}(X^0\pm X^1),$$
(2.64)
and assume the compactifications
$$X^{}X^{}+\frac{4\pi }{\sqrt{k}},H_1^{}H_1^{}+2\pi .$$
(2.65)
These are indeed consistent with the previous compactifications $`Y^{}Y^{}+{\displaystyle \frac{4\pi }{\sqrt{k+2}}}`$, $`H_1H_1+2\pi `$, because we can obtain from (2.52)
$$\{\begin{array}{c}X^{}=\frac{2}{\sqrt{k(k+2)}}Y^++\sqrt{\frac{k+2}{k}}Y^{}+\frac{2}{\sqrt{k}}H_1\hfill \\ H_1^{}=\frac{2}{\sqrt{k+2}}Y^++H_1.\hfill \end{array}$$
(2.66)
We likewise introduce the tachyon vertices compatible with this light-like compactification
$$V_{j,m,p}=e^{\left(\frac{\sqrt{k}}{2}p\frac{2m}{\sqrt{k}}\right)iX^++\frac{\sqrt{k}}{2}piX^{}\sqrt{\frac{2}{k}}j\rho },$$
(2.67)
and the corresponding Fock vacua satisfying
$`J_0^3|j,m,p`$ $`=`$ $`(m{\displaystyle \frac{k}{2}}p)|j,m,p,`$ (2.68)
$`J_p^\pm |j,m,p`$ $`=`$ $`(m\pm j)|j,m\pm 1,p,`$ (2.69)
$`J_{p+n}^\pm |j,m,p`$ $`=`$ $`0,(n1),`$ (2.70)
$`L_0|j,m,p`$ $`=`$ $`\left({\displaystyle \frac{1}{k}}j(j1)+mp{\displaystyle \frac{k}{4}}p^2\right)|j,m,p.`$ (2.71)
The total $`SL(2;\text{R})`$ currents are also rewritten in the new coordinates;
$`J^3`$ $`=`$ $`\sqrt{{\displaystyle \frac{k}{2}}}iX^0`$ (2.72)
$`J^\pm `$ $`=`$ $`\left(\sqrt{{\displaystyle \frac{k}{2}}}iX^1\pm \sqrt{{\displaystyle \frac{k}{2}}}\rho \mathrm{\Psi }^+\mathrm{\Psi }^{}\pm \sqrt{2}\mathrm{\Psi }^0\mathrm{\Psi }^\pm \right)e^{\frac{2}{\sqrt{k}}iX^+}.`$ (2.73)
To close this section, we present the space-time superconformal algebra in our new variables, which again has the forms reminiscent of the DDF operators.
First, the space-time Virasoro algebra is given (in the $`(1)`$-picture) by
$$\{\begin{array}{c}_0=\sqrt{\frac{k}{2}}e^\varphi \mathrm{\Psi }^0\hfill \\ _n=\sqrt{\frac{k}{2}}e^\varphi e^{n\frac{2}{\sqrt{k}}iX^+}(\mathrm{\Psi }^1+ni\mathrm{\Psi }^\rho )(n0).\hfill \end{array}$$
(2.74)
We again mention that $`_n(n0)`$ is the unique solution of the BRST constraint among the operators of the form $`{\displaystyle e^\varphi e^{n\frac{2}{\sqrt{k}}iX^+}(A\mathrm{\Psi }^\rho +B\mathrm{\Psi }^0+C\mathrm{\Psi }^1)}`$, $`A0`$.
The space-time $`U(1)_R`$ current is given by
$$𝒥_n=\sqrt{2k}e^\varphi e^{n\frac{2}{\sqrt{k}}iX^+}\mathrm{\Psi }^2.$$
(2.75)
To construct the space-time super currents we must introduce the spin fields. According to , we bosonize the “deformed $`U(1)_R`$ current” on the world-sheet;
$`J_R^{}`$ $`:=`$ $`J_RQi(X^0+X^2)`$ (2.76)
$``$ $`iH_1^{}iH_2i\sqrt{3}H_3Qi(X^0+X^1),`$
where we set
$$\begin{array}{c}iH_1^{}=\mathrm{\Psi }^+\mathrm{\Psi }^{}(\text{as defined before})\hfill \\ iH_2=\mathrm{\Psi }^0\mathrm{\Psi }^2\hfill \\ i\sqrt{3}H_3=J_{𝒩/U(1)}QiX^2.\hfill \end{array}$$
(2.77)
(The combined current $`J_{𝒩/U(1)}QiX^2`$ actually has the Schwinger term $`{\displaystyle \frac{3}{z^2}}`$.) The practical reason why we do so is as follows: $`J_R`$ has non-trivial OPEs with the vertex operators such as $`e^{n\frac{2}{\sqrt{k}}iX^+}`$, and thus the DDF like operators including the spin fields (see (2.81)) made up of $`J_R`$ do not nicely behave under the BRST transformation. In contrast, we can rather simply solve the BRST condition for the vertices associated to the current $`J_R^{}`$, since it has no singular OPE with $`e^{n\frac{2}{\sqrt{k}}iX^+}`$.
Now the spin fields should take the form (up to cocycle factors)
$$e^{\frac{i}{2}(ϵ_1H_1^{}+ϵ_2H_2+\sqrt{3}ϵ_3H_3)}e^{ϵ_4\frac{i}{2}Q(X^0+X^1)}(ϵ_i=\pm 1).$$
(2.78)
The GSO projection leaves only a half of them satisfying
$$\underset{i=1}{\overset{3}{}}ϵ_i=1,$$
(2.79)
and we use the notation
$$S^{ϵ_3ϵ_1}=e^{\frac{i}{2}(ϵ_1H_1^{}+ϵ_2H_2+\sqrt{3}ϵ_3H_3)}$$
(2.80)
to express the spin fields allowed by the GSO condition. We can explicitly verify that the fermion vertices of the type $`{\displaystyle e^{\frac{\varphi }{2}}e^{\frac{i}{2}ϵ_4Q(X^0+X^1)}S^{ϵ_3ϵ_1}}`$ are BRST invariant, if and only if $`ϵ_4=ϵ_1`$, namely the vertices of the type $`{\displaystyle e^{\frac{\varphi }{2}}e^{\frac{2}{\sqrt{k}}\frac{ϵ_1}{2}iX^+}S^{ϵ_3ϵ_1}}`$. They define the space-time $`N=4`$ SUSY (the global part of the space-time $`N=2`$ superconformal symmetry in NS sector). We can work out more general fermion vertices of the type
$$e^{\frac{\varphi }{2}}e^{r\frac{2}{\sqrt{k}}iX^+}\underset{ϵ_3,ϵ_1}{}A_{ϵ_3,ϵ_1}S^{ϵ_3ϵ_1},(r\frac{1}{2}+\text{Z}).$$
(2.81)
The BRST condition uniquely (up to BRST exact terms and an overall normalization) determines the coefficients $`A_{ϵ_3,ϵ_1}`$ and we finally obtain the physical vertices
$$𝒢_r^\pm =k^{1/4}e^{\frac{\varphi }{2}}e^{r\frac{2}{\sqrt{k}}iX^+}\left[(r+\frac{1}{2})S^{\pm +}(r\frac{1}{2})S^\pm \right],(r\frac{1}{2}+\text{Z}).$$
(2.82)
They generate the $`N=2`$ superconformal algebra (in NS sector) together with $`_n`$, $`𝒥_n`$, and the central charge is equal to $`6kp`$ on the Fock space over the vacuum $`|,,p`$. In fact, it is a straightforward calculation to check that these space-time superconformal generators (2.74), (2.75), (2.82) coincide with the ones constructed in up to BRST exact terms.
## 3 Analyses on Spectra of Physical States
In this section we shall investigate the spectrum of physical states in our reformulation of $`AdS_3`$ string theory. The unitarity of physical Hilbert space is an important problem. We first analyze the physical spectrum as that of a free field theory, namely without taking the Liouville potential term into account, and later discuss the effect of turning on this term.
### 3.1 Spectrum as Free Field Theory
First, we analyze the spectrum of bosonic string on $`AdS_3\times 𝒩`$. Let us recall that the Fock vacuum is defined from the tachyon vertex
$$|j,m,p=\underset{z0}{lim}V_{j,m,p}(z)|0,$$
(3.1)
and we denote the corresponding Fock space as $`_{j,m,p}`$ from now on. We also define “out state” as
$$j,m,p|=\underset{z\mathrm{}}{lim}0|V_{j,m,p}(z)z^{2h_{j,m,p}},$$
(3.2)
where
$$h_{j,m,p}=\frac{1}{k}j(j1)+mp\frac{k+2}{4}p^2.$$
(3.3)
Using the momentum conservation and taking account of the existence of background charge along the $`\rho `$-direction, we obtain
$$1j,m,p|j,m,p0,$$
(3.4)
and the other combinations vanish. Notice that we must use the following Hermitian conjugation
$$1j,m,p|=\left(|j,m,p\right)^{}$$
(3.5)
to discuss the unitarity.
We shall here neglect the Liouville potential term (that is, the “screening charge term” $`{\displaystyle \beta \overline{\beta }e^{\sqrt{\frac{2}{k}}\phi }}{\displaystyle Y^1\overline{}Y^1e^{\sqrt{\frac{2}{k}}\rho }}`$ in the $`\sigma `$-model action (2.1)). This means that the Hilbert space of physical states should be defined as the BRST cohomology on the Fock spaces of the free fields $`Y^0`$, $`Y^1`$, $`\rho `$ properly tensored by the Hilbert space of $`𝒩`$ sector. Deciding the physical spectrum is a rather simple problem as in the usual string theory on the flat Minkowski space (at least as long as we only take the primary states in the $`𝒩`$ sector in constructing the physical states). However, there is one non-trivial point: the existence of background charge in the $`\rho `$-direction. As is well-known, the BRST constraint eliminates two longitudinal degrees of freedom in the case of Minkowski background: one is eliminated by the BRST condition itself and another becomes BRST exact. From a physical point of view this aspect is closely related to the fact that the first excited states (the graviton states) become mass-less, and thus one of the polarization vectors must be light-like. On the other hand, in our case of linear dilaton theory we have a mass gap originating from the background charge of $`\rho `$ and so the first excited states become massive. This implies that one of the two longitudinal modes does not decouple from the physical Hilbert space, since all the polarization vectors are space-like.
To make this point clearer, let us consider a simple example. We are here given one transverse oscillator $`i\rho ={\displaystyle \underset{n}{}}{\displaystyle \frac{\alpha _n^\rho }{z^{n+1}}}`$ and two longitudinal oscillators $`iY^\pm =\pm {\displaystyle \underset{n}{}}{\displaystyle \frac{\alpha _n^\pm }{z^{n+1}}}`$. The simplest candidate of the first excited states is
$$\alpha _1^\rho |p^+,p^{},p^\rho c_1|0_{\text{gh}},$$
(3.6)
where the on-shell condition is given by
$$p^+p^{}+\frac{1}{2}(p^\rho )^2+\frac{1}{4k}=0.$$
(3.7)
(The relation of the momenta with our previous convention is given by $`p^+={\displaystyle \frac{\sqrt{k+2}}{2}}p`$, $`p^{}={\displaystyle \frac{\sqrt{k+2}}{2}}p{\displaystyle \frac{2m}{\sqrt{k+2}}}`$, $`p^\rho =i\sqrt{{\displaystyle \frac{2}{k}}}\left(j{\displaystyle \frac{1}{2}}\right)`$.) Now we assume $`p^+p^{}0`$. The BRST transformation of this candidate (3.6) yields a non-vanishing term due to the background charge. We must cancel it by mixing the longitudinal modes to recover the BRST invariance. After some simple calculation we find two independent solutions
$$\left(\alpha _1^\rho +i\frac{Q}{p^\pm }\alpha _1^\pm \right)|p^+,p^{},p^\rho c_1|0_{\text{gh}}.$$
(3.8)
In the usual free string theory the general physical states are created by making the transverse DDF operators act on suitable Fock vacuum (“allowed states”). The above simple observation suggests that, in our case of $`AdS_3`$, we must make use of the two independent DDF operators that are not purely transversal. Some candidates for the suitable DDF operators are already given by Satoh <sup>2</sup><sup>2</sup>2In the other candidates for the DDF operators are also proposed. However, they include the ghost number current explicitly in their expressions and are not BRST invariant. Y.S. should express his great thanks to Dr. Satoh for the discussion about this point.
$$B_n^\pm :=i\left(\rho +\frac{Q}{2}\mathrm{ln}Y^\pm \right)e^{\frac{n}{p^\pm }iY^\pm },$$
(3.9)
which are BRST invariant and satisfy the commutation relation of a free boson
$$\begin{array}{c}[B_m^\pm ,B_n^\pm ]=m\delta _{m+n,0},\hfill \\ B_m^\pm |p^+,p^{},p^\rho =0,(m1)\hfill \end{array}$$
(3.10)
Moreover, $`B_n^\pm `$ act on the Fock space defined by $`|p^+,p^{},p^\rho `$ as the operators $`\alpha _n^\rho +\mathrm{}`$, as is expected, and it is easy to check that $`B_1^\pm `$ actually give rise to the first excited states discussed above (3.8). So the reader might suppose that we can naively use $`B_n^+`$, $`B_n^{}`$ to construct the complete set of physical states. But this is not the case. $`B_n^+`$ and $`B_m^{}`$ are not mutually local and thus we can never use them at the same time. The best we can do is to use only one of them, $`B_m^+`$, which is compatible with the light-like compactification (2.14). We rewrite it as
$$𝒜_n^{(p)}=i\left(\rho +\frac{Q}{2}\mathrm{ln}Y^+\right)e^{\frac{n}{p}\frac{2}{\sqrt{k+2}}iY^+},[𝒜_m^{(p)},𝒜_n^{(p)}]=m\delta _{m+n,0}.$$
(3.11)
which are defined as local operators on $`_{,,p}`$ ($`p0`$).
Now, the question we must solve is as follows: What is the missing DDF operator that can compensate (3.11)? As we already suggested, the answer is the space-time Virasoro operators. Let $`_n`$ be the space-time Virasoro operators defined in (2.25). $`_{\frac{n}{p}}`$ are well-defined as local operators on the Fock space $`_{j,m,p}`$ and behave as $`\alpha _n^\rho +\mathrm{}`$. Therefore they are the candidates of the missing DDF operators. Alternatively we shall define
$$_n^{(p)}:=p_{\frac{n}{p}}\frac{k+2}{4}(p^21)\delta _{n,0},$$
(3.12)
which are shown to generate a Virasoro algebra with the central charge $`c=6(k+2)`$ (irrespective of the value $`p`$). Furthermore, $`_n^{(p)}`$ are mutually local with $`𝒜_m^{(p)}`$ and satisfy the commutation relation
$$[_n^{(p)},𝒜_m^{(p)}]=m𝒜_{m+n}^{(p)}+i\alpha n^2\delta _{n+m,0},\alpha \sqrt{\frac{k}{2}}(1\frac{1}{k}).$$
(3.13)
It is also convenient to introduce improved Virasoro operators
$$\stackrel{~}{}_n^{(p)}:=_n^{(p)}\frac{1}{2}\underset{m}{}:𝒜_{(n+m)}^{(p)}𝒜_m^{(p)}:i\alpha n𝒜_n^{(p)}\frac{1}{2}\alpha ^2\delta _{n,0},$$
(3.14)
which are defined so that they commute with $`\{𝒜_m^{(p)}\}`$ and generate the Virasoro algebra with $`c=23{\displaystyle \frac{6}{k}}c_𝒩`$. This value of the central charge is quite expected. One can show that the DDF operators $`\stackrel{~}{}_n^{(p)}`$ correspond to the energy-momentum tensor of $`𝒩`$ sector in the light-cone gauge.
In this way, we have found that the physical Hilbert space should be spanned by the states having the forms
$$\stackrel{~}{}_{n_1}^{(p)}\stackrel{~}{}_{n_2}^{(p)}\mathrm{}𝒜_{m_1}^{(p)}𝒜_{m_2}^{(p)}|j,m,p\mathrm{},(n_i1,m_i1).$$
(3.15)
In order to complete our discussion we must confirm the linear independence of the states of the type (3.15). Happily, this is very easy to prove in our case. The Virasoro algebra $`\{\stackrel{~}{}_n^{(p)}\}`$ has the central charge greater than 1 for sufficiently large $`k`$, and thus the Kac determinant does not vanish, as is well-known in the representation theory of Virasoro algebra.
We can now present the complete list of physical states. This spectrum is specified by the momenta of the Fock space $`_{j,m,p}\overline{}_{j,\overline{m},p}`$ previously defined. The light-like compactification (2.14) implies that $`p\text{Z}`$, and also $`\stackrel{~}{}_0^{(p)}\overline{\stackrel{~}{}_0^{(p)}}p(_0\overline{}_0)p\text{Z}`$ (the “level matching condition”). Since $`𝒜_0^{(p)}=𝒜_0^{(p)}`$ holds and we have
$$𝒜_0^{(p)}|j,m,p=i\sqrt{\frac{2}{k}}\left(j\frac{1}{2}\right)|j,m,p,$$
(3.16)
the value of $`j`$ allowed by the unitarity is $`j={\displaystyle \frac{1}{2}}+is`$ ($`s\text{R}`$), at least when $`p0`$. It corresponds to the principal continuous series of unitary representation of $`SL(2;\text{R})`$. Also in the case of $`p=0`$ we can show that only the principal series is permitted from the unitarity, as we will observe below.
To avoid unnecessary complexity we shall only take the primary states in the $`𝒩`$ sector in constructing the physical states, and assume the conformal weights $`h_𝒩`$ of these primary states are non-negative. It is not difficult to construct more general physical states including the descendant states in the $`𝒩`$ sector, if we are concretely given the unitary CFT model describing this sector.
We discuss the $`p=0`$ and $`p0`$ cases separately.
1. $`p=0`$
In this case, the DDF operators of the types $`𝒜_n^{(p)}`$, $`_n^{(p)}`$ are ill-defined. But we must require the space-time Virasoro operators $`\{_n\}`$ (with central charge 0) should define an unitary representation, since they are well-defined as local operators even in this sector. $`_n`$ simply maps a Fock vacuum to another Fock vacuum and so the representation space is given by $`_{r\text{Z}}\text{C}|j,m+r,0`$ with arbitrary fixed values of $`m\text{R}`$, $`j`$. We can obtain
$$\begin{array}{c}1j,mr,p|_n_n|j,m+r,p\hfill \\ =\left\{\left(m+r+\frac{n}{2}\right)^2n^2\left(j\frac{1}{2}\right)^2\right\}1j,mrn,0|j,m+r+n,0.\hfill \end{array}$$
(3.17)
Since the conformal weight (in the sense of world-sheet CFT) must be real at least, $`j`$ should take the values $`j\text{R}`$ or $`j={\displaystyle \frac{1}{2}}+is`$ $`(s\text{R})`$. If $`j{\displaystyle \frac{1}{2}}i\text{R}`$ (principal series), the coefficient of R.H.S in (3.17) is always positive and we have an unitary representation of $`\{_n\}`$. On the other hand, if $`j\text{R}`$, we can always find $`r\text{Z}`$ for which this coefficient becomes negative as long as we choose $`n`$ to be sufficiently large. This means that the cases of $`j\text{R}`$ cannot be unitary representations of $`\{_n\}`$, and hence we must rule out these sectors.
The general physical states with $`p=0`$ are written as
$`|{\displaystyle \frac{1}{2}}+is,m,0|h_𝒩c_1|0_{\text{gh}}`$
$``$ $`|{\displaystyle \frac{1}{2}}+is,\overline{m},0|h_𝒩\overline{c}_1|0_{\text{gh}}`$
$`m,\overline{m}\text{R},m\overline{m}\text{Z},`$
where $`|h_𝒩`$ is the primary state with conformal weight $`h_𝒩`$ in the $`𝒩`$ sector and $`|0_{\text{gh}}`$ is the vacuum of the ghost system. The on-shell condition is given by
$$\frac{1}{k}\left(s^2+\frac{1}{4}\right)+h_𝒩=1.$$
(3.19)
If $`k>1/4`$, we can always solve the on-shell condition (3.19) for $`h_𝒩=0`$. These physical states are tachyons whose mass-squared are lower than the Breitenlohner-Freedman bound . Such an instability in bosonic string theory is not surprising, and we later observe that the GSO projection successfully eliminates these tachyonic states in the superstring case.
There is one comment: If we took account of the unitarity of the representation only of $`SL(2;\text{R})`$ (that is, $`\{_n\}`$, $`n=0,\pm 1`$), many representations would survive in the sector $`j\text{R}`$: the discrete series $`𝒟_j^\pm `$ and the exceptional series $`_{j,\alpha }`$, as is well-known and many readers might expect. It is crucial in the above argument to take the full Virasoro algebra $`\{_n\}`$ in place of the $`SL(2;\text{R})`$ subalgebra.
2. $`p0`$
As we already discussed, in this sector $`j`$ must be equal to $`j={\displaystyle \frac{1}{2}}+is`$ $`(s\text{R})`$, and the physical Hilbert space is generated by the actions of the DDF operators $`\{𝒜_n^{(p)}\}`$, $`\{\stackrel{~}{}_n^{(p)}\}`$ $`(n\text{Z}_1)`$ on the on-shell Fock vacua. We must discuss the positivity of the norm of such physical states. Obviously $`𝒜_n^{(p)}`$ create only positive norm states, and do not lead to any constraint. However, the Virasoro generators $`\{\stackrel{~}{}_n^{(p)}\}`$ give rise to a non-trivial constraint for unitarity. Since this Virasoro algebra has the central charge $`c>1`$, the condition for the unitarity means that the $`\stackrel{~}{}_0^{(p)}`$-eigenvalue of the Fock vacuum $`|{\displaystyle \frac{1}{2}}+is,m,p`$ is non-negative. (We again assume $`k`$ is sufficiently large.) It is easy to show that this unitarity condition is equivalent to a simple inequality $`h_𝒩0`$ thanks to the on-shell condition
$$\frac{1}{k}\left(s^2+\frac{1}{4}\right)+mp\frac{k+2}{4}p^2+h_𝒩=1.$$
(3.20)
This equivalence is not surprising, since $`\stackrel{~}{}_0^{(p)}`$ corresponds to the stress tensor for $`𝒩`$ sector in the light-cone gauge. In this way we conclude that the no-ghost theorem for this sector is trivially satisfied as long as the internal CFT $`𝒩`$ is unitary. This result is consistent with those of , although we here take a different convention of free field representation: Our convention diagonalizes the time-like current $`j^3`$ (corresponding to the energy operator). On the other hand, those given in diagonalize one of the space-like currents. We remark that the light-cone momentum $`p`$ plays a role similar to that of the extra zero-mode momentum emphasized in .
One can find that the $`_0`$-eigenvalue (not $`\stackrel{~}{}_0^{(p)}`$) of the on-shell Fock vacuum, which corresponds to the space-time energy, is bounded below
$$_0\frac{h_𝒩}{p}+\frac{(k1)^2}{4kp}+\frac{k+2}{4}\left(p\frac{1}{p}\right)\frac{h_𝒩1}{p}+\frac{k+2}{4}p,$$
(3.21)
(for a sufficiently large value $`k`$). This means that this sector corresponds to the long string states in the sense of and belongs to the continuous spectrum above the threshold energy $`{\displaystyle \frac{k}{4}}p`$ discussed in .
In summary, the general physical states are written as
$`\stackrel{~}{}_{n_1}^{(p)}\stackrel{~}{}_{n_2}^{(p)}\mathrm{}𝒜_{m_1}^{(p)}𝒜_{m_2}^{(p)}\mathrm{}|{\displaystyle \frac{1}{2}}+is,m,p|h_𝒩c_1|0_{\text{gh}}`$ (3.22)
$``$ $`\overline{\stackrel{~}{}^{(p)}}_{\overline{n}_1}\overline{\stackrel{~}{}^{(p)}}_{\overline{n}_2}\mathrm{}\overline{𝒜^{(p)}}_{\overline{m}_1}\overline{𝒜^{(p)}}_{\overline{m}_2}\mathrm{}|{\displaystyle \frac{1}{2}}+is,\overline{m},p|\overline{h}_𝒩\overline{c}_1|0_{\text{gh}}`$
$`n_1,n_2,\mathrm{}1,m_1,m_2,\mathrm{}1`$
$`\overline{n}_1,\overline{n}_2,\mathrm{}1,\overline{m}_1,\overline{m}_2,\mathrm{}1,`$
where the on-shell conditions are
$`{\displaystyle \frac{1}{k}}\left(s^2+{\displaystyle \frac{1}{4}}\right)+mp{\displaystyle \frac{k+2}{4}}p^2+h_𝒩`$ $`=`$ $`1`$
$`{\displaystyle \frac{1}{k}}\left(s^2+{\displaystyle \frac{1}{4}}\right)+\overline{m}p{\displaystyle \frac{k+2}{4}}p^2+\overline{h}_𝒩`$ $`=`$ $`1,`$ (3.23)
and the “level matching condition” is given as
$$\underset{i}{}n_i+\underset{j}{}m_jmp=\underset{i}{}\overline{n}_i+\underset{j}{}\overline{m}_j\overline{m}p(\text{mod}p).$$
(3.24)
The superstring case $`AdS_3\times S^1\times 𝒩/U(1)`$ is similarly analyzed. The unitarity of the physical Hilbert space is derived from the unitarity in the $`N=2`$ SCFT describing $`𝒩/U(1)`$ sector. We here only discuss how tachyonic states in the $`p=0`$ sector are eliminated by the GSO projection.
The tachyon vertex operators have slightly different expressions as compared with the bosonic case
$$V_{j,m,p}=e^{\left(\frac{\sqrt{k}}{2}p\frac{2m}{\sqrt{k}}\right)iX^++\frac{\sqrt{k}}{2}piX^{}\sqrt{\frac{2}{k}}j\rho }.$$
(3.25)
Together with the vertex for the $`S^1`$ direction $`e^{i\sqrt{\frac{2}{k}}qX^2}`$ we construct the Fock vacuum $`|j,m,p,q`$ such that
$`J_0^3|j,m,p,q`$ $`=`$ $`(m{\displaystyle \frac{k}{2}}p)|j,m,p,q,`$ (3.26)
$`J_p^\pm |j,m,p,q`$ $`=`$ $`(m\pm j)|j,m\pm 1,p,q,`$ (3.27)
$`J_{p+n}^\pm |j,m,p,q`$ $`=`$ $`0,(n1),`$ (3.28)
$`L_0|j,m,p,q`$ $`=`$ $`\left({\displaystyle \frac{1}{k}}j(j1)+mp{\displaystyle \frac{k}{4}}p^2+{\displaystyle \frac{q^2}{k}}\right)|j,m,p,q,`$ (3.29)
$`J_{R0}^{}|j,m,p,q`$ $`=`$ $`\left({\displaystyle \frac{2q}{k}}+p\right)|j,m,p,q.`$ (3.30)
In the $`p=0`$ sector the argument similar to the bosonic case leads to the following physical states;
$`|{\displaystyle \frac{1}{2}}+is,m,0,q|h_𝒩,q_𝒩ce^\varphi |0_{\text{gh}}`$
$``$ $`|{\displaystyle \frac{1}{2}}+is,\overline{m},0,\overline{q}|\overline{h}_𝒩,\overline{q}_𝒩\overline{c}e^{\overline{\varphi }}|0_{\text{gh}}`$
$`m,\overline{m}\text{R},m\overline{m}\text{Z},`$
with the on-shell conditions
$$\begin{array}{c}\frac{1}{k}\left(s^2+\frac{1}{4}\right)+\frac{q^2}{k}+h_𝒩=\frac{1}{2}\hfill \\ \frac{1}{k}\left(s^2+\frac{1}{4}\right)+\frac{\overline{q}^2}{k}+\overline{h}_𝒩=\frac{1}{2}.\hfill \end{array}$$
(3.32)
Naively we can solve the on-shell conditions as in the bosonic case and they are tachyonic states except for $`s=0`$. However, we can show that the GSO projection eliminates such tachyonic states, as is expected. From the on-shell conditions (3.32) and the assumption $`h_𝒩{\displaystyle \frac{1}{2}}|q_𝒩|`$, which is derived from the unitarity in the $`𝒩/U(1)`$-sector, we obtain the inequality
$$\frac{1}{k}(\frac{1}{4}+s^2)+\frac{q^2}{k}+\frac{|q_𝒩|}{2}\frac{1}{2}.$$
(3.33)
We should define the GSO condition with respect to the deformed $`U(1)_R`$ current $`J_R^{}`$ and it reads as $`{\displaystyle \frac{2q}{k}}+q_𝒩=2l+1`$ ($`l\text{Z}`$). First we assume $`q_𝒩0`$. If $`l0`$, substituting $`q_𝒩=2l+1{\displaystyle \frac{2q}{k}}`$ into the above inequality (3.33), we obtain
$$s^2+\left(q\frac{1}{2}\right)^2+2lk0,$$
(3.34)
which leads to $`n=0`$, $`q={\displaystyle \frac{1}{2}}`$, $`s=0`$. In the case of $`l<0`$, $`q{\displaystyle \frac{k}{2}}(2l+1)`$ must hold. We thus obtain
$`{\displaystyle \frac{1}{2}}`$ $``$ $`{\displaystyle \frac{1}{k}}\left({\displaystyle \frac{1}{4}}+s^2\right)+{\displaystyle \frac{q^2}{k}}+{\displaystyle \frac{|q_𝒩|}{2}}{\displaystyle \frac{s^2}{k}}+{\displaystyle \frac{1}{2}},`$ (3.35)
which leads to $`s=0`$, again. Therefore the tachyonic states whose mass-squared are lower than the BF bound are successfully eliminated by the GSO projection. We can repeat the same analysis when $`q_𝒩<0`$.
### 3.2 Physical Spectrum under the Existence of Liouville Potential
In the previous argument only the principal series was allowed. In the physical sense it was a quite natural result, because we regarded the system as a free system and thus all the momenta should be real.
Now, let us try to turn on the Liouville potential term (or the screening charge term in the world-sheet action (2.1)). In this case we can expect some physical states with an imaginary momentum along the $`\rho `$-direction describing the bound states (“bound string states” in the terminology of ).
Unfortunately, a rigorous treatment of the quantum Liouville theory as an interacting theory is quite non-trivial. Instead we shall here take the operator contents as free fields and treat the Liouville potential as a small perturbation.
Recalling the $`\sigma `$-model action (2.1), this perturbation term may be identified with the operator $`S\overline{S}`$, where
$$S=\beta e^{\sqrt{\frac{2}{k}}\phi }\sqrt{\frac{k+2}{2}}iY^1e^{\sqrt{\frac{2}{k}}\rho }$$
(3.36)
is no other than the familiar screening charge which commutes with all modes of $`\widehat{SL}(2;\text{R})`$ currents. As for the spectrum generating operators, we may as well expect that at least $`_{\pm 1}(_{\pm p}^{(p)}),_0`$ remain the good DDF operators, since they commute with the screening charge (3.36).
On the other hand, because such an interaction breaks the translational invariance along the $`\rho `$-direction, the $`\rho `$-momenta $`i\left(j{\displaystyle \frac{1}{2}}\right)`$ loses its meaning. However the second Casimir $`j(j1)`$ remains well-defined as a conserved quantity characterizing the physical states even under the interacting theory. This is nothing but the standard argument of “screening out” of the extra zero-mode momentum in the free field representation of $`CFT_2`$ . We may expect the bound string states possessing the imaginary $`\rho `$-momenta as long as their second Casimirs take real values. In this way we can no longer regard $`𝒜_0^{(p)}`$ as a good DDF operator. Moreover, we must also rule out the non-zero modes $`𝒜_n^{(p)}`$ $`(n0)`$, because we have the following commutation relations $`𝒜_0^{(p)}[(_1)^n,𝒜_n^{(p)}]+\text{const. }`$ for $`n>0`$, and $`𝒜_0^{(p)}[(_1)^n,𝒜_n^{(p)}]+\text{const. }`$ for $`n<0`$.
It is a subtle problem whether the other modes of Virasoro operators $`_n^{(p)}`$ ($`n0,\pm p`$) remain the members of the spectrum generating algebra, since they also do not commute with the screening charge (3.36). However, it may be plausible to admit these operators from the point of view of the $`AdS_3/CFT_2`$ correspondence or the arguments of Brown-Henneaux . The fact that only the $`SL(2;\text{R})`$ generators $`_0,_{\pm 1}`$ commute with the screening charge and the other modes do not is supposed to reflect the following fact: In the argument of the true isometry generates only the $`SL(2;\text{R})`$ and the other modes merely correspond to the asymptotic isometries, which can be regarded as symmetries only near the boundary. We shall now propose that the DDF operators suitable for the interacting theory including (3.36) are $`\{_n^{(p)}\}`$ rather than those for the free system $`\{\stackrel{~}{}_n^{(p)},𝒜_n^{(p)}\}`$. This claim is consistent with the analyses based on the light-cone gauge for the long string configuration . Hence our assertion is likely to be consistent at least with the spectrum of the long string located near the boundary. (One should keep it in mind that our assumption of small Liouville perturbation is valid only for such a configuration of world-sheet.)
Based on this assumption we now present the complete physical spectrum as the interacting theory. For the principal series $`j={\displaystyle \frac{1}{2}}+is`$, the results are similar to those of free fields. The case of $`p=0`$ is the same as before, and when $`p0`$, only we have to do is to replace the DDF operators $`\{𝒜_n^{(p)},\stackrel{~}{}_m^{(p)}\}`$ in the expression of (3.22) by $`\{_n^{(p)}\}`$. They likewise belong to a continuous spectrum above the threshold energy $`{\displaystyle \frac{kp}{4}}`$.
A crucial difference is the existence of the physical states with $`j\text{R}`$ as we already suggested. To discuss the unitarity of this sector we again assume $`p>0`$, and the cases of $`p<0`$ can be analyzed in the same way. The unitarity condition means that the $`_0^{(p)}`$-eigenvalue of the Fock vacuum should be non-negative. Thanks to the on-shell condition
$$\frac{1}{k}j(j1)+mp\frac{k+2}{4}p^2+h_𝒩=1,$$
(3.37)
we can immediately obtain the following inequality for the unitarity
$$\frac{1}{k}\left(j\frac{1}{2}\right)^2h_𝒩+\frac{(k1)^2}{4k}.$$
(3.38)
We must also restrict the range of $`j`$ as $`j>1/2`$ because of the normalizability of wave function (see ). Especially, in the case of $`h_𝒩=0`$ we obtain the unitarity condition
$$\frac{1}{2}<j\frac{k}{2}.$$
(3.39)
These physical states do not propagate along the radial direction $`\rho `$, and are supposed to correspond to the bound string states in the argument of . In fact, we can evaluate the space-time energy for this sector
$$\frac{k+2}{4}\left(p\frac{1}{p}\right)\stackrel{<}{}_0\stackrel{<}{}\frac{h_𝒩1}{p}+\frac{k+2}{4}p,$$
(3.40)
which is consistent with the result given in .
The physical states are summarized as follows;
$`_{n_1}^{(p)}_{n_2}^{(p)}\mathrm{}|j,m,p|h_𝒩c_1|0_{\text{gh}}`$ (3.41)
$``$ $`\overline{_{\overline{n}_1}^{(p)}}\overline{_{\overline{n}_2}^{(p)}}\mathrm{}|j,\overline{m},p|\overline{h}_𝒩\overline{c}_1|0_{\text{gh}}`$
$`n_1,n_2,\mathrm{}1,\overline{n}_1,\overline{n}_2,\mathrm{}1,`$
where the on-shell conditions are
$`{\displaystyle \frac{1}{k}}j(j1)+mp{\displaystyle \frac{k+2}{4}}p^2+h_𝒩`$ $`=`$ $`1`$
$`{\displaystyle \frac{1}{k}}j(j1)+\overline{m}p{\displaystyle \frac{k+2}{4}}p^2+\overline{h}_𝒩`$ $`=`$ $`1.`$ (3.42)
and the “level matching condition” is given as
$$\underset{i}{}n_imp=\underset{i}{}\overline{n}_i\overline{m}p(\text{mod}p).$$
(3.43)
We here remark that the positive (negative) energy ($`_00`$) physical states should have $`p>0`$ ($`p<0`$). This fact will be important in the discussions in the next section about the interpretation of the long string theory.
To close this section we compare the above spectrum with the result of . For this purpose we must clarify which representations of $`\widehat{SL}(2;\text{R})`$ we should choose.
Let us assume $`p<0`$. Going back to the Wakimoto free fields $`\phi `$, $`\beta `$, $`\gamma `$ we introduce the “Wakimoto module” $`𝒲_{j,m,p}`$ which is defined as the Fock space generated by $`\alpha _{n1}^\phi `$, $`\beta _{pn}`$, $`\gamma _{pn}`$ $`(n0)`$ out of the vacuum $`|j,m,p`$ for $`mj`$, and by $`\alpha _{n1}^\phi `$, $`\beta _{pn1}`$, $`\gamma _{pn}`$ $`(n0)`$ for $`m=j`$ (corresponding to the flowed discrete series $`\widehat{𝒟}_j^{+(p)}`$ in ). $`𝒲_{j,m,p}`$ is obviously a subspace of $`_{j,m,p}`$ (and they are not isomorphic). Moreover, at least under the restriction $`{\displaystyle \frac{1}{2}}<j{\displaystyle \frac{k}{2}}`$ (3.39), we can show that $`𝒲_{j,m,p}`$ can be identified with some (reducible, in general) $`\widehat{SL}(2;\text{R})`$-module, since we have no singular vectors in the corresponding Verma module (except for the Fock vacua themselves). It is easy to see that
$$\underset{i}{}_{n_i}^{(p)}|j,m,p\underset{r\text{Z}}{}𝒲_{j,m+\frac{r}{p},p}.$$
(3.44)
Therefore we can successfully realize the actions of DDF operators $`_n^{(p)}`$ within the (reducible) representations corresponding to $`{\displaystyle \underset{r\text{Z}}{}}𝒲_{j,m+\frac{r}{p},p}`$.
For $`p>0`$, the essentially same argument works by introducing the “inverse Wakimoto representation”;
$$\{\begin{array}{c}j^3=\stackrel{~}{\beta }\stackrel{~}{\gamma }\sqrt{\frac{k}{2}}\phi \hfill \\ j^+=\stackrel{~}{\beta }\hfill \\ j^{}=\stackrel{~}{\beta }\stackrel{~}{\gamma }^2\sqrt{2k}\stackrel{~}{\gamma }\stackrel{~}{\phi }(k+2)\stackrel{~}{\gamma },\hfill \end{array}$$
(3.45)
or more explicitly,
$$\{\begin{array}{c}\stackrel{~}{\phi }=\rho \sqrt{\frac{2k}{(k+2)}}Y^+\hfill \\ \stackrel{~}{\beta }=\left(\sqrt{\frac{k+2}{2}}iY^1+\sqrt{\frac{k}{2}}\rho \right)e^{\frac{2}{\sqrt{k+2}}iY^+}\hfill \\ \stackrel{~}{\gamma }=e^{\frac{2}{\sqrt{k+2}}iY^+}.\hfill \end{array}$$
(3.46)
Consequently our choice of the representations of $`\widehat{SL}(2;\text{R})`$ is much larger than that of for the cases of $`j\text{R}`$, although the enlargement of the Hilbert space by the spectral flow in is incorporated into our setup as the discrete light-cone momentum $`p`$. By construction our physical Hilbert space also contains no ghosts. In the analysis in $`m`$ must take discrete values related with a fixed $`j`$ (belonging to the discrete series of $`\widehat{SL}(2;\text{R})`$ transformed by the spectral flow); $`jm\text{Z}`$. On the other hand, in our analysis $`m`$ is arbitrary and independent of $`j`$ as long as it satisfies the on-shell condition. This is natural from our starting point: the $`\sigma `$-model (2.1) rather than the abstract representation theory of affine Lie algebra, and thus $`j`$ and $`m`$ (and of course $`p`$, too) should correspond to independent momenta along the different directions.
Furthermore, $`\widehat{𝒟}_j^{+(p)}`$ and $`\widehat{𝒟}_{\frac{k+2}{2}j}^{(p1)}`$ are identified in , since they are equivalent as an irreducible representation of $`\widehat{SL}(2;\text{R})`$. Nevertheless they should be distinguished from our viewpoint, because they possess the different light-cone momenta $`p`$. Especially, the standard discrete series $`\widehat{𝒟}_j^+`$ (lowest weight representations), $`\widehat{𝒟}_j^{}`$ (highest weight representations) are realized in the sectors $`p=1`$, $`p=1`$ respectively, since the sectors $`p=0`$, $`j\text{R}`$ are excluded in our analysis.
We also mention that the above unitarity condition (3.39) is analogous to the result given in , but it is a slightly stronger condition. The unitarity bound proposed in reads $`{\displaystyle \frac{1}{2}}<j<{\displaystyle \frac{k+2}{2}}`$ and the one given in reads $`{\displaystyle \frac{1}{2}}<j<{\displaystyle \frac{k+1}{2}}`$ in our convention<sup>3</sup><sup>3</sup>3The unitarity bound for is stronger than that of due to the identification $`\widehat{𝒟}_j^{+(p)}=\widehat{𝒟}_{\frac{k+2}{2}j}^{(p1)}`$ mentioned above. Moreover, the same range of $`j`$ was proposed in a different context by requiring good behaviors of the two point functions of the non-normalizable primary operators. Such two point functions nicely behave, too, under our constraints (3.39), since it is more stringent than that of .. This disagreement originates from the different choices of the representations mentioned above. In fact, if one choose to restrict the value of $`m`$ as $`jm\text{Z}`$ when solving the on-shell condition, one can show the no-ghost theorem under the assumption $`{\displaystyle \frac{1}{2}}<j<{\displaystyle \frac{k+2}{2}}`$ same as rather than (3.39)<sup>4</sup><sup>4</sup>4To be precise, under this restriction $`jm\text{Z}`$ we must take $`_n`$ as the DDF operators rather than $`_n^{(p)}=p_{n/p}+\mathrm{}`$, (recall (3.44)) and hence the unitarity here means that $`_n`$-descendants should not include any negative norm states.. It is not yet clear whether our choice of the momenta $`m`$, independent of $`j`$, is completely valid even in the rigid treatment as an interacting theory, or nothing but an artifact originating from the free field approximation. However, we again emphasize that our setup of physical Hilbert space admits the whole actions of DDF operators $`\{_n^{(p)}\}`$ (and necessarily also with the space-time Virasoro algebra $`\{_n\}`$). This fact is found to be consistent with the several results about the long string sectors given by the light-cone gauge approach , as we will comment in the next section. In fact, one can readily find that our choice of $`m`$ so as to be independent of $`j`$ is crucial in order to ensure the equivalence with the spectrum in the light-cone gauge after solving the on-shell condition. In this sense we believe that our physical spectrum is valid at least for the long string configuration near the boundary, which is well described by the light-cone gauge approach. In order to justify this spectrum beyond the near boundary region we will have to carry out a further analysis with the Liouville interaction term treated more precisely.
The extension of the above arguments to superstring examples is not so difficult and we do not present it here. We instead focus on the spectrum of on-shell chiral primaries of superstring on $`AdS_3\times S^1\times 𝒩/U(1)`$ in the next section.
## 4 Chiral Primaries and Spectral Flow
In this section we further study the spectrum in the superstring cases. We especially investigate an important class of observables: chiral primaries. In other words, we shall concentrate on the topological sector of superstring vacua on $`AdS_3\times S^1\times 𝒩/U(1)`$ . They are significant from the perspective of $`AdS_3/CFT_2`$-duality. Although we have not yet achieved the complete understanding of this duality, the study of their spectrum will certainly clarify some aspects of it.
Through this section we only deal with the left moving parts of objects, and it is easy to complete our discussions by taking also the right movers.
### 4.1 Background with Space-time $`N=4`$ SUSY
We first discuss the most familiar superstring vacua with space-time $`N=4`$ SUSY;
$$AdS_3\times S^3\times T^4AdS_3\times S^1\times SU(2)/U(1)\times T^4,$$
(4.1)
where $`SU(2)/U(1)`$ means that the Kazama-Suzuki model for this coset with $`c=3{\displaystyle \frac{6}{N}}`$ ($`N2`$ is equal to the level of (bosonic) $`SU(2)`$-WZW model describing the $`S^3`$ sector), which is identified with the $`N=2`$ minimal model of $`A_{N1}`$ type and we denote it by $`M_N`$ from now on. The criticality condition gives $`k=N`$, and this background is regarded as the near horizon limit of $`NS5/NS1`$ system, as is well-known.
Let $`|\mathrm{\Phi }_l`$ ($`l=0,1,\mathrm{},N2`$) be the chiral primary states in the $`M_N`$ sector with $`h_𝒩={\displaystyle \frac{q_𝒩}{2}}={\displaystyle \frac{l}{2N}}`$. We must look for the chiral primary states in the total system by tensoring $`|\mathrm{\Phi }_l`$ with suitable vertex operators in the $`AdS_3\times S^1`$ sector. In the notation of previous section, namely,
$$|j,m,p,q\underset{z0}{lim}e^{\left(\frac{\sqrt{k}}{2}p\frac{2m}{\sqrt{k}}\right)iX^++\frac{\sqrt{k}}{2}piX^{}\sqrt{\frac{2}{k}}j\rho \sqrt{\frac{2}{k}}qiX^2}|0,$$
(4.2)
the possible candidates for the desired vertices are written as follows;
$$|j,j,p,j\frac{k}{2}p,\mathrm{\Psi }_{1/2}^+|j,(j1),p,(j1)\frac{k}{2}p.$$
(4.3)
They are primary states with respect to $`T(z)`$, $`G(z)`$ and also satisfy
$$G_{1/2}^+|j,j,p,j\frac{k}{2}p=G_{1/2}^+\mathrm{\Psi }_{1/2}^+|j,(j1),p,(j1)\frac{k}{2}p=0.$$
(4.4)
First we consider $`|j,j,p,j\frac{k}{2}p|\mathrm{\Phi }_lc_1e^\varphi |0_{\text{gh}}`$. The on-shell condition leads to
$$j=\frac{Nl}{2}.$$
(4.5)
(The GSO condition is automatically satisfied, since we have the relation $`h={\displaystyle \frac{Q}{2}}`$.) Similarly, for the second candidate $`\mathrm{\Psi }_{1/2}^+|j,(j1),p,(j1)\frac{k}{2}p|\mathrm{\Phi }_lc_1e^\varphi |0_{\text{gh}}`$ we can solve the on-shell condition and obtain
$$j=1+\frac{l}{2}.$$
(4.6)
From now on we denote the first type of chiral primary (4.5) as $`|l,p;1`$ and the second type (4.6) as $`|l,p;2`$. Namely, we set
$`|l,p;1`$ $`:=`$ $`|{\displaystyle \frac{Nl}{2}},{\displaystyle \frac{Nl}{2}},p,{\displaystyle \frac{N(p1)+l}{2}}|\mathrm{\Phi }_lce^\varphi |0_{\text{gh}}`$ (4.7)
$`|l,p;2`$ $`:=`$ $`\mathrm{\Psi }_{1/2}^+|{\displaystyle \frac{l}{2}}+1,{\displaystyle \frac{l}{2}},p,{\displaystyle \frac{Np+l}{2}}|\mathrm{\Phi }_lce^\varphi |0_{\text{gh}}.`$ (4.8)
They are both normalizable and satisfy the unitarity condition (3.38).
Remarkably one can find that (4.7), (4.8) are also chiral primaries with respect to the space-time superconformal algebra $`\{_n,𝒢_r^\pm ,𝒥_n\}`$, as suggested in . They satisfy
$`_0|l,p;1`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒥_0|l,p;1={\displaystyle \frac{1}{2}}\{l+N(p1)\}|l,p;1`$ (4.9)
$`_0|l,p;2`$ $`=`$ $`{\displaystyle \frac{1}{2}}𝒥_0|l,p;2={\displaystyle \frac{1}{2}}(l+Np)|l,p;2.`$ (4.10)
Since the light-cone momentum $`p`$ can now take an arbitrary integer, we have infinite number of on-shell chiral primaries. All of them have the same $`U(1)_R`$ charge in the sense of world-sheet because of the the on-shell condition. But they have the different $`U(1)_R`$ charges with respect to the space-time conformal algebra.
Let us study the action of the spectral flow on these states. A natural extension of the spectral flow (2.23) to the superstring case is given by
$`\mathrm{U}_pX^0(z)\mathrm{U}_p^1`$ $`=`$ $`X^0(z)p\sqrt{{\displaystyle \frac{k}{2}}}i\mathrm{ln}z`$
$`\mathrm{U}_pX^2(z)\mathrm{U}_p^1`$ $`=`$ $`X^2(z)+p\sqrt{{\displaystyle \frac{k}{2}}}i\mathrm{ln}z,`$ (4.11)
and the other fields should remain unchanged by the spectral flow. The total $`\widehat{SL}(2;\text{R})`$-currents are transformed by them as follows;
$`\mathrm{U}_pJ_n^3\mathrm{U}_p^1`$ $`=`$ $`J_n^3+{\displaystyle \frac{kp}{2}}\delta _{n,0}`$
$`\mathrm{U}_pJ_n^\pm \mathrm{U}_p^1`$ $`=`$ $`J_{np}^\pm .`$ (4.12)
We can show that
$$\mathrm{U}_r|l,p;i=|l,p+r;i.$$
(4.13)
Namely, the spectral flow maps an on-shell chiral primary to another on-shell chiral primary. This is a general feature. In fact, it is a straightforward calculation to check that
$`\mathrm{U}_pG^+(z)\mathrm{U}_p^1`$ $`=`$ $`G^+(z)`$ (4.14)
$`\mathrm{U}_pQ_{BRST}\mathrm{U}_p^1`$ $`=`$ $`Q_{BRST}\underset{z0}{lim}\left[pc\sqrt{{\displaystyle \frac{k}{2}}}i(X^0+X^2)+p\eta e^\varphi \sqrt{{\displaystyle \frac{k}{2}}}(\mathrm{\Psi }^0+\mathrm{\Psi }^2)\right].`$ (4.15)
Because the correction term in (4.15) vanishes when acting on an arbitrary on-shell chiral primaries, the spectral flow $`\mathrm{U}_p`$ preserves the on-shell condition in the space of chiral primaries. One should remark that $`\mathrm{U}_p`$ does not transform all the physical states among themselves. Indeed the physical states which are not the chiral primaries are transformed into off-shell states by $`\mathrm{U}_p`$.
Turning our attention to the space-time conformal algebra, we obtain
$`\{𝒢_{\frac{1}{2}}^{},𝒢_{\frac{1}{2}}^+\}`$ $`=`$ $`_0{\displaystyle \frac{1}{2}}𝒥_0`$ (4.16)
$`=`$ $`{\displaystyle \sqrt{\frac{k}{2}}i(X^0+X^2)},`$
and this identity is unchanged by the spectral flow. This means that the spectral flows are closed in the space of the space-time chiral primaries, which is consistent with the above observation.
Next we present some remarks from the view points of $`AdS_3/CFT_2`$-duality and the long string theory given in . Although our understanding of them by string theory is not yet complete, we believe the following remarks are useful to clarify some important aspects of them.
Tracing back to the procedure of our field redefinitions, one can find that $`|l,p=1;1`$ can be identified with the space-time chiral primary states given in (see also )<sup>5</sup><sup>5</sup>5They correspond to $`|\omega _{l/2}^0`$ in the notation of , which contain the trivial cohomology in the $`T^4`$ sector. We can also consider more general space-time chiral primary states that contain higher cohomologies of $`T^4`$, as given in . But they are not chiral primaries in the sense of world-sheet.. It has the following structure
$$|l,1;1=\underset{z0}{lim}𝒪_l(z)|0,1;1,$$
(4.17)
where
$$𝒪_l(z):=V_{\frac{l}{2},\frac{l}{2},0,\frac{l}{2}}\mathrm{\Phi }_l(z)e^{\frac{l}{\sqrt{2N}}(iX^0+iX^1+iX^2+\rho )}\mathrm{\Phi }_l(z)$$
(4.18)
naturally corresponds to the chiral operator in the light-cone gauge formalism of the long string theory given in . $`|0,1;1e^{\sqrt{\frac{N}{2}}(iX^1+\rho )}ce^\varphi |0`$ has the maximal $`j`$-value $`j=N/2k/2`$ and is the same as the “space-time vacuum” (or “long string vacuum”) presented in . It satisfies
$`_n|0,1;1`$ $`=`$ $`0,(n1),`$ (4.19)
$`𝒢_r^\pm |0,1;1`$ $`=`$ $`0,(r1/2),`$ (4.20)
$`𝒥_n|0,1;1`$ $`=`$ $`0,(n1),`$ (4.21)
$`{\displaystyle \frac{\sqrt{k}}{2}}{\displaystyle iX^+|0,1;1}`$ $`=`$ $`|0,1;1,`$ (4.22)
where the last line simply means that $`p{\displaystyle \gamma ^1\gamma }=1`$, and these identities hold up to BRST-exact terms.
As discussed in (see also ), the chiral operator $`𝒪_l`$ corresponds to a non-normalizable state. Its wave function is divergent near the boundary, where the Coulomb branch CFT is weakly coupled. On the other hand, thanks to the existence of $`|0,1;1`$, the chiral primary state $`|l,1;1`$ itself becomes normalizable state vanishing exponentially at large $`\rho `$, as expected from the observation about the Higgs branch tube in . In this sense the interpretation of $`|0,1;1`$ as the long string vacuum may be natural.
It may be also useful to define explicitly the space-time chiral primary operator $`\widehat{𝒪}_l(x){\displaystyle \underset{n}{}}{\displaystyle \frac{\widehat{𝒪}_{l,n}}{x^{n+\frac{l}{2}}}}`$ by introducing the vertex operators
$$\widehat{𝒪}_{l,n}:=V_{\frac{l}{2},n,0,\frac{l}{2}}\mathrm{\Phi }_l(\mathrm{\Psi }^0+\mathrm{\Psi }^1)e^\varphi ,$$
(4.23)
where $`n`$ runs over all (half-)integers if $`{\displaystyle \frac{l}{2}}`$ is an (half-)integer. $`\widehat{𝒪}_l(x)`$ actually behaves as a chiral primary operator with respect to the space-time SCA. For example, we obtain
$$[_m,\widehat{𝒪}_{l,n}]=\left\{(\frac{l}{2}1)mn\right\}\widehat{𝒪}_{l,m+n},$$
(4.24)
which means that $`\widehat{𝒪}_l(x)`$ is a primary operator with conformal weight $`h={\displaystyle \frac{l}{2}}`$. We can further show that
$$\widehat{𝒪}_{l,n}|0,1;1=0,(n>\frac{l}{2})$$
(4.25)
and also obtain the “operator-state correspondence”
$$\widehat{𝒪}_{l,\frac{l}{2}}|0,1;1=|l,1;1$$
(4.26)
(up to the picture changing and an overall constant).
For $`p>1`$ we can consider more general chiral primaries with the higher space-time $`U(1)_R`$-charges
$$|l,p;1=\mathrm{U}_{p1}|l,1;1=𝒪_l(0)|0,p;1=\widehat{𝒪}_{l,\frac{l}{2}}|0,p;1.$$
(4.27)
As we observed above (4.9), the spectrum of $`𝒥_0`$ charge is $`l+N(p1)`$, $`l=0,1,\mathrm{},N2`$, $`p1`$.
We have not yet known the suitable interpretation of the “graviton-like” chiral primary states $`|l,p;2`$ in the context of $`AdS_3/CFT_2`$ correspondence. We only point out that they do not seem to have the forms such as (4.26), and so it might be plausible to suppose that they do not have any counterparts in the boundary theory, as long as our identification of $`|0,1;1`$ with the space-time vacuum is justified. In any case we will need a further analysis to give a more definite statement about this problem.
The following aspect may be worthwhile to point out. Here the chiral primaries with the higher space-time $`U(1)_R`$-charges appear in the sector with higher light-cone momenta $`p`$ (or by taking account of the degrees of freedom of spectral flows). On the other hand, in the analysis of , they correspond to the $`\text{Z}_p`$-twisted sector of the symmetric orbifold theory, which describes the sector of long string with the “length” $`p`$ as in the Matrix string theory . It suggests a remarkable correspondence between the spectral flow in the covariant gauge formalism and the twisted sector of the symmetric orbifold in the light-cone gauge formalism .
To address the precise correspondence between them we should work on the second quantized framework. It is quite reasonable from the viewpoints of $`AdS_3/CFT_2`$ correspondence, since the boundary CFT should also contain multi-particle excitations. We shall now focus on the physical states with positive energies, which should have the light-cone momenta $`p0`$ as we found in section 3. The physical Hilbert space of the first quantized string states, which was studied in our previous analyses, is then decomposed to $`={\displaystyle \underset{p0}{}}_p`$, where $`_p`$ denotes the sector with the light-cone momentum $`p0`$. The Hilbert space in the (free) second quantized theory can be roughly written as
$$\widehat{}=\underset{n=0}{\overset{\mathrm{}}{}}()^n.$$
(4.28)
(To be precise, we must make some (anti-)symmetrization to assure the correct statistics in this and the expressions given below.) The second quantized space $`\widehat{}`$ has a natural decomposition with respect to the total light-cone momentum
$$\widehat{}=\underset{p0}{}\widehat{}_p,$$
(4.29)
Obviously $`\widehat{}_p`$ is decomposed to the subspaces of the forms
$$_{p_1}_{p_2}\mathrm{}_{p_n}_0_0\mathrm{},(p_1p_2\mathrm{}p_n1,\underset{i=1}{\overset{n}{}}p_i=p).$$
(4.30)
The $`p=0`$ Hilbert space $`_0`$ only contains tachyons, which are eliminated by the GSO projection, as was already shown<sup>6</sup><sup>6</sup>6The fact that the physical Hilbert space of “short string sector” $`_0`$ is vacant is not a contradiction. It rather means that only the non-normalizable physical operators can appear and the operator-state correspondence fails in the short string sector as suggested in .. Therefore we can neglect the $`_0`$ factors, and can explicitly write down
$$\widehat{}_p=\underset{n=1}{\overset{p}{}}\underset{_{i=1}^np_i=p}{}(_{i=1}^n_{p_i}).$$
(4.31)
Now let us consider the system of $`Q_5`$ NS5 and $`Q_1`$ NS1. The NS5 charge $`Q_5(N)`$ appears in the world-sheet action of $`AdS_3`$-string theory, but $`Q_1`$ does not. It only appears in the string coupling, which is stable under the near horizon limit, and hence we cannot find this effect in the first quantized theory. However, in the second quantized theory, it is quite natural to identify the NS1 charge $`Q_1`$ with the total light-cone momentum $`p`$ in the expression of $`\widehat{}_p`$. Hence we propose that the physical Hilbert space of this NS5-NS1 system should be defined as $`\widehat{}_{Q_1}`$. Notice that it has the structure characterized by the various partitions $`\{p_i(1);{\displaystyle \underset{i}{}}p_i=Q_1\}`$ which is consistent with the expected correspondence to the symmetric orbifold theory. Clearly this system can be decomposed to the subsystems of various long strings with the “lengths” (or “windings”) $`p_i`$ ($`1p_iQ_1`$ $`{\displaystyle \underset{i}{}}p_i=Q_1`$), as observed in .
It is interesting to present the spectrum of the “single-particle chiral primaries” in this framework. Let $`pQ_1`$ be a positive integer. We obtain the required states as
$$|l,p;1\underset{Q_1p\text{-times}}{\underset{}{|0,1;1\mathrm{}|0,1;1}},$$
(4.32)
which satisfies $`_0={\displaystyle \frac{1}{2}}𝒥_0={\displaystyle \frac{l+Q_5(p1)}{2}}`$ as expected. Since $`l`$, $`p`$ run over the ranges $`0lQ_52`$, $`0pQ_1`$, this spectrum is completely in agreement with the result of , in which the (multiple) long string CFT was analyzed using the symmetric orbifold theory. Quite remarkably, this has the upper bound $`{\displaystyle \frac{Q_1Q_5}{2}}`$ which is expected from the $`AdS_3/CFT_2`$-duality , as already commented in .
Notice that there are the missing states corresponding to the $`𝒥_0`$-charge $`Q_5p1`$ $`(1pQ_1)`$. They are (formally) written as
$$|Q_51,p;1|0,1;1\mathrm{}=𝒪_{Q_51}(0)|0,p;1|0,1;1\mathrm{},$$
(4.33)
and “$`𝒪_{Q_51}(z)`$” is no other than the missing chiral operator discussed in , which should correspond to the cohomology with a delta function support at the small instanton singularity. In this sense these missing states (4.33) are supposed to be the natural generalizations to the cases of $`1<pQ_1`$ of the one discussed in , in which only the $`p=1`$ sector was treated.
The above observation implies that the first quantized Hilbert space $`_p`$ ($`pQ_1`$) precisely corresponds to the $`\text{Z}_p`$-twisted sector in the $`S_{Q_1}`$-orbifold theory as we already suggested. The relation
$$_n=\frac{1}{p}_{pn}^{(p)}+\frac{Q_5}{4}\left(p\frac{1}{p}\right)\delta _{n,0}$$
(4.34)
indeed confirms this identification. On the Hilbert space $`_p`$, the space-time Virasoro algebra $`_n`$ has the central charge $`c=6pQ_5`$, and the DDF operators $`_n^{(p)}`$ generate the Virasoro algebra with $`c=6Q_5`$. This relation (4.34) is the same as the well-known formula to define the conformal algebra describing the $`\text{Z}_p`$-twisted sector of the symmetric orbifold. It is easy to define the tensor product representation of space-time conformal algebra with $`c=6Q_1Q_5`$ on the second quantized Hilbert space $`\widehat{}_{Q_1}{\displaystyle \underset{_ip_i=Q_1}{}}(_{p_i})`$ including the conformal invariant vacuum
$$\underset{Q_1\text{times}}{\underset{}{|0,1;1|0,1;1\mathrm{}|0,1;1}}_1^{Q_1}.$$
(4.35)
Such a correspondence, which was essentially suggested in , is quite expected from our standpoint as the discrete light-cone theory fitted to the spirit of Matrix string . Recall that our setup of physical Hilbert space in section 3 allows the action of $`_n^{(p)}=p_{n/p}+\mathrm{}`$, and moreover we must impose the “level matching condition” $`_0^{(p)}\overline{_0^{(p)}}p\text{Z}`$ onto the Hilbert space $`_p`$. These facts are crucial to establish the above correspondence to the symmetric orbifold.
One should keep in mind the following fact: one can also construct the representation with $`c=6Q_1Q_5`$ on the first quantized Hilbert space $`_{Q_1}`$ that is the subspace of $`\widehat{}_{Q_1}`$ describing the single long string with the maximal length $`Q_1`$. However, $`_{Q_1}`$ cannot include the conformal invariant vacuum. Recall that $`_0|0,p;10`$, unless $`p=1`$. More precisely speaking, we can show that, in our setup of the first quantized Hilbert space the BRST-invariant state with the properties (4.19), (4.20), (4.21) and non-zero $`p`$ is possible only if $`p=1`$, and the solution is unique (up to BRST exact terms and an overall normalization), $`|0,1;1`$, as suggested in . This fact leads us to the only one possibility of the conformal invariant vacuum (4.35). The large Hagedorn density suited for $`c=6Q_1Q_5`$, which may reproduce the correct entropy formula of black-hole, should be attached to $`\widehat{}_{Q_1}`$, not to $`_{Q_1}`$, since $`_{Q_1}`$ does not include the vacuum state such that $`_0=0`$ (see the discussions given in ).
### 4.2 Background with Space-time $`N=2`$ SUSY
In principle it is not difficult to generalize the above analysis on chiral primaries to more general superstring vacua with space-time $`N=2`$ SUSY .
We first give a rather generic argument. Consider superstring theory on $`AdS_3\times S^1\times 𝒩/U(1)`$, where $`𝒩/U(1)`$ is an arbitrary $`N=2`$ SCFT of center $`96/k`$. As in the $`N=4`$ case, we can construct two series of chiral primary states from chiral primaries $`V_j`$ of conformal weight $`j/2k`$ in the $`𝒩/U(1)`$-sector:
$`|j,p;1`$ $`:=`$ $`|{\displaystyle \frac{kj}{2}},{\displaystyle \frac{kj}{2}},p,{\displaystyle \frac{k(p1)+j}{2}}|V_jce^\varphi |0_{\mathrm{gh}}`$ (4.36)
$`|j,p;2`$ $`:=`$ $`\mathrm{\Psi }_{1/2}^+|{\displaystyle \frac{j+2}{2}},{\displaystyle \frac{j}{2}},p,{\displaystyle \frac{kp+j}{2}}|V_jce^\varphi |0_{\mathrm{gh}}`$ (4.37)
They have the light-cone momentum $`p`$ and the following conformal weight:
$`_0|j,p;1={\displaystyle \frac{1}{2}}𝒥_0|j,p;1`$ $`=`$ $`{\displaystyle \frac{j+k(p1)}{2}}|j,p;1`$
$`_0|j,p;2={\displaystyle \frac{1}{2}}𝒥_0|j,p;2`$ $`=`$ $`{\displaystyle \frac{j+kp}{2}}|j,p;2`$ (4.38)
Note that, if we take as $`𝒩/U(1)`$ an arbitrary $`N=2`$ SCFT of center $`96/k`$, the conformal weight $`h=j/2k`$ of $`V_j`$ runs within the range $`0h32/k`$. However, it is only if $`0h1/21/k`$ that the chiral primary states are in the spectrum allowed from unitarity and normalizability.
Let us consider a specific example. Take as $`𝒩/U(1)`$ the $`N=2`$ minimal model which we denote by $`M_N`$ as before. It was proposed in that the superstring theory on this background is marginally equivalent to the non-critical superstring theory which is holographically dual to the decoupled theory based on the $`A_{N1}`$-singular $`CY_4`$. In this case the criticality condition leads to $`k={\displaystyle \frac{N}{N+1}}`$.
Let $`|\mathrm{\Phi }_l`$ ($`l=0,1,\mathrm{},N2`$) be again the chiral primary states of weight $`l/2N`$ in the $`M_N`$ sector. We obtain the following chiral primaries;
$`|l,p;1`$ $`:=`$ $`|{\displaystyle \frac{Nl}{2(N+1)}},{\displaystyle \frac{Nl}{2(N+1)}},p,{\displaystyle \frac{N(p1)+l}{2(N+1)}}|\mathrm{\Phi }_lce^\varphi |0_{\mathrm{gh}}`$ (4.39)
$`|l,p;2`$ $`:=`$ $`\mathrm{\Psi }_{1/2}^+|{\displaystyle \frac{l}{2(N+1)}}+1,{\displaystyle \frac{l}{2(N+1)}},p,{\displaystyle \frac{Np+l}{2(N+1)}}|\mathrm{\Phi }_lce^\varphi |0_{\mathrm{gh}}.`$ (4.40)
In this way we have again infinite number of on-shell chiral primaries possessing the following space-time $`U(1)_R`$ charges;
$`𝒥_0|l,p;1`$ $`=`$ $`\left({\displaystyle \frac{l}{N+1}}+{\displaystyle \frac{N}{N+1}}(p1)\right)|l,p;1`$ (4.41)
$`𝒥_0|l,p;2`$ $`=`$ $`\left({\displaystyle \frac{l}{N+1}}+{\displaystyle \frac{N}{N+1}}p\right)|l,p;2.`$ (4.42)
Unfortunately, $`|l,p;1`$ is non-normalizable and $`|l,p;2`$ does not satisfy the unitarity constraints (3.38). Hence we cannot consider the chiral primary states within the physical Hilbert space. We must only treat these chiral primaries as operators and cannot expect the operator-state correspondence. Nevertheless, they may be regarded as an important class of operators in the context of $`AdS_3/CFT_2`$-duality, or more general holographic dualities . In particular the non-normalizable chiral primaries $`|l,p;1`$ (“tachyon-like operators”) may be important, because they possess the momentum structures which can be regarded as natural generalizations of those of the scaling operators in the space-time conformal theory proposed in . Since the light-cone momentum $`p`$ runs over an infinite range, we can obtain the infinite tower of space-time chiral operators for each of the chiral operators in the “matter sector” $`\mathrm{\Phi }_l`$. This aspect may be interesting, since they look like analogues of “gravitational descendants” in the theory of two dimensional gravity. We must make further studies to gain more precise insights about these objects. In addition, the roles of the graviton-like operators $`|l,p;2`$ are again unclear. More detailed argument for them will be surely important, although it is beyond the scope of this paper.
## 5 Conclusions and Discussions
In this paper we have studied the spectrum of the physical states in string theory on $`AdS_3`$ based on a free field realization. We have found that the system is quite simply described as a linear dilaton theory with a light-like compactification, which we called as “discrete light-cone Liouville theory”. Our key idea is to utilize the DDF operators according to the traditional approach to string theory on flat backgrounds. We have two independent sets of DDF operators; $`𝒜_n^{(p)}`$, $`\stackrel{~}{}_n^{(p)}`$. This situation is similar to the non-critical string, although we started with the critical string on $`AdS_3\times 𝒩`$ background. In fact, we can easily find that a suitable linear combination of $`𝒜_n^{(p)}`$ and $`\stackrel{~}{}_n^{(p)}`$ gives the “longitudinal DDF operator” utilized in .
Regarding the system as a free theory with no Liouville interaction term (screening charge term), the physical spectrum contains only the principal series $`j={\displaystyle \frac{1}{2}}+is`$ as asserted in , and the physical states are generated by the DDF operators $`\{𝒜_n^{(p)},\stackrel{~}{}_n^{(p)}\}`$ mentioned above.
However, once we turn on the Liouville potential, the story becomes rather non-trivial. In this interacting theory the translational invariance along the Liouville direction is broken. The physical Hilbert space is expected to be spanned only by the $`\widehat{SL}(2;\text{R})`$ currents, rather than the whole oscillators of the string coordinates $`\rho ,Y^+,Y^{}`$, because the interaction term commutes only with the $`\widehat{SL}(2;\text{R})`$ currents. In this situation the spectrum generating algebra becomes $`\{_n^{(p)}\}`$ in place of $`\{𝒜_n^{(p)},\stackrel{~}{}_n^{(p)}\}`$, and the physical states possessing the imaginary $`\rho `$-momenta ($`j\text{R}`$) are also allowed. Physically they correspond to the bound string states of that are trapped inside the $`AdS_3`$-space.
It may be worthwhile to mention that only the physical Hilbert space as the interacting Liouville theory may be consistent with the microscopic evaluation of the black hole entropy. In the free system the DDF operators should be $`\{𝒜_n^{(p)},\stackrel{~}{}_m^{(p)}\}`$ and $`\stackrel{~}{}_m^{(p)}`$ (which are identified with the Virasoro operators in $`𝒩`$-sector under the light-cone gauge) have a small central charge $`c=23{\displaystyle \frac{6}{k}}`$. (An important discussion related to such a counting of physical states was given in .) On the other hand, after turning on the Liouville potential we claimed that the full Virasoro generators $`\{_n^{(p)}\}`$ are well-defined (and it is also crucial that $`\{𝒜_n^{(p)}\}`$ should be discarded). Taking further account of the second quantized Hilbert space they seem to generate sufficiently many physical states with the Hagedorn density that can reproduce the correct entropy. We would like to study this problem in more detail elsewhere.
In the study of superstring examples, we have presented the complete set of on-shell chiral primaries. There exist infinite number of such operators and the spectral flows naturally act on them. Moreover, to describe the well-known $`Q_5(k)`$ NS5 - $`Q_1`$ NS1 system we made use of the second quantized framework. The Hilbert space of the multiple long string system was given as the form $`\widehat{}_{Q_1}={\displaystyle \underset{_ip_i=Q_1}{}}(_i_{p_i})`$, where $`_p`$ denotes the first quantized physical Hilbert space of the sector with the light-cone momentum $`p(>0)`$. This space reproduce the spectrum of chiral primaries same as that given by the symmetric orbifold theory , and among other things, we have successfully obtained the upper bound $`Q_1Q_5/2`$ consistent with the prediction of $`AdS_3/CFT_2`$ correspondence .
It may be also worth pointing out that our reformulation of superstring on $`AdS_3\times S^1\times 𝒩/U(1)`$ has the same field contents as those of the non-critical string that is holographically dual to a singular Calabi-Yau compactification (especially, the cases of $`CY_4`$) proposed in . The only difference between these models is the existence/absence of the light-like compactification. In it was discussed that these two backgrounds can be interpolated by some marginal deformation. It may be an interesting problem to clarify the rigid correspondence between them. In particular, our analyses on general chiral operators in section 4 will be readily generalized to the cases of such non-critical string theories.
## Acknowledgement
Y. S. would like to thank I. Bars and Y. Satoh for helpful discussions.
The work of K. H. is supported in part by Japan Society for Promotion of Science under the Postdoctral Research Program ($`\mathrm{}12`$-$`02721`$). The work of Y. S. is supported in part by Grant-in-Aid for Encouragement of Young Scientists ($`\mathrm{}11740142`$) and also by Grant-in-Aid for Scientific Research on Priority Area ($`\mathrm{}707`$) “Supersymmetry and Unified Theory of Elementary Particles”, both from Japan Ministry of Education, Science, Sports and Culture.
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# On the ‘Anomalous’ Resurgence of Shot Noise in Long Conductors
## I Introduction
The fine-scale investigation of carrier noise in mesoscopic conductors is now a well-established field within the transport physics of the solid state. It has reached new heights recently, experimentally and theoretically (Blanter and Büttiker 2000). Noise is a unique source of information on the dynamics of microscopic fluctuations. This is notably so at small length scales, already approaching the quantum domain in actual devices.
A particular aspect of mesoscopic noise, in metallic diffusive conductors especially, is the so-called ‘crossover’ from thermal to shot noise. There are many measurements of it, and almost as many compelling (if frequently quite disparate) theoretical explanations. The term crossover refers to the smooth evolution, with increasing voltage, of the current-noise spectral density $`S(V;\omega )`$ (normally it is sufficient to study its low-frequency limit $`\omega \tau ^1`$, where $`\tau `$ is a characteristic collision time). One sees the onset of a non-dissipative excess component in $`S(V;0)`$, over and above the dissipative Johnson-Nyquist noise which exhausts the low-field limit. Thus, typically,
$$\frac{S(V;0)}{S_0}=1+\gamma \left[\frac{eV}{2k_\mathrm{B}T}\mathrm{coth}\left(\frac{eV}{2k_\mathrm{B}T}\right)1\right]$$
(1)
where $`S_0=4Gk_\mathrm{B}T`$ is the Johnson-Nyquist value and $`\gamma `$ the suppression factor; $`G`$ is the sample conductance and $`k_\mathrm{B}T`$ the thermal energy. The factor $`\gamma `$ is a signature of the often subtle correlation effects in the microscopic fluctuations, which are responsible for the form of $`S(V;0)`$ as measured (Blanter and Büttiker 2000). At voltages $`Vk_\mathrm{B}T/e`$, Equation (1) gives $`S(V;0)=2\gamma eI`$, where $`I=GV`$ is the current. This exhibits suppression of the Schottky formula $`S=2eI`$ for classical Poissonian shot noise.
While Eq. (1) gives an impressive empirical fit to many (though not all) experiments, serious questions arise as to whether any of the prevailing models (Kogan 1996; Blanter and Büttiker 2000) possesses the internal consistency expected of a standard microscopic description. We cite Das and Green (2000) and Green and Das (2000a) for a description of what may go awry with such theories, all of which rely heavily on hydrodynamic drift-diffusion analogies for mesoscopic transport (Datta 1995; Kogan 1996). Theoretically, the concept of the ‘crossover’ may be as open to new critiques (Gillespie 2000) as it is to new empirical tests by appropriately designed experiments (Green and Das 1998a, 2000a-c).
Now, a fresh window on the ‘crossover’ appears to have been opened by the recent work of Gomila and Reggiani (2000). Foremost is their emphasis on the explicit role of carrier-number fluctuations as generators of the observable shot noise. We note that this perspective was addressed theoretically (with detailed computations) in Green and Das (1998a, 2000a).
There is a basic difference between fluctuations of carrier number, manifesting at the conductor-lead interfaces and engendering shot noise, and fluctuations of the free energy, manifesting throughout the conductor’s volume and generating thermal noise. This understanding is in sharp contrast to the usual phenomenological viewpoint (Kogan 1996) in which no difference is permitted, even in principle, between thermally related and carrier-number related processes. In this respect it is useful to bring to mind the textbook distinction between a variation with respect to chemical potential and a variation with respect to particle number. The fact that they are intimately linked by microscopics does not override the fact that they are thermodynamic conjugates, with wholly distinguishable thermodynamic consequences.
Gomila and Reggiani (2000) certainly raise weighty, if not wholly unanticipated, points regarding macroscopic solid-state shot noise (Green and Das 1998a; Naveh 1998). Furthermore, these are not easily addressed in strictly low-field, linear, drift-diffusive descriptions (Datta 1995; Kogan 1996; Blanter and Büttiker 2000). Thus it is well to revisit our own existing non-perturbative Boltzmann theory, allied to a time-of-flight intepretation of shot noise (Green and Das 1998a). In Section II we briefly recall our formalism, while in Sec. III we give simple Drude-like (but complete and exact) kinetic solutions for shot noise as a function of current, dimensionality, length, and finite radius of the conductor. These show that the ‘crossover’ anomaly exists with no reference at all to long-range Coulomb correlations. In Sec. IV we examine the results of Gomila and Reggiani (2000), based on a diffusive Langevin-Poisson scheme, and compare them with those of Green and Das (1998a) as recalled in the present paper. Section V contains our final remarks.
## II Kinetics of shot noise: theory
### A Non-equilibrium fluctuations
In the context of Gomila and Reggiani (2000) we specialise to a homogeneous metallic wire, subject to a uniform driving field. The mean carrier density is $`n`$ and the mean total carrier number is $`N=\mathrm{\Omega }n`$ in the sample volume $`\mathrm{\Omega }`$. These are independent of the external field. The kinetic Boltzmann equation for the spatially uniform distribution function $`f_𝐤(t)`$ is, in the Drude-Lorentz collision approximation,
$$\left(\frac{}{t}+\frac{eE}{\mathrm{}}\frac{}{k_x}\right)f_𝐤(t)=\frac{1}{\tau }\left(f_𝐤(t)\frac{f(t)}{f^{\mathrm{eq}}}f_𝐤^{\mathrm{eq}}\right).$$
(2)
Our electrons are positive for convenience. The field $`E`$ is in the source-drain $`x`$-direction, and $`\tau `$ is the collision time. We denote traces over wave-vector (with a factor of two for spin) as $`f2\mathrm{\Omega }^1_𝐤f_𝐤`$. The physical constraint on the traces in the right side of Eq. (2) is, naturally, $`f(t)=f^{\mathrm{eq}}=n`$. Finally,
$`f_𝐤^{\mathrm{eq}}=[1+\mathrm{exp}(\epsilon _𝐤\mu )/k_\mathrm{B}T)]^1`$
is the usual Fermi-Dirac equilibrium distribution, parametrised by the thermal energy and by $`\mu `$, the chemical potential. Note that there is no coupling to the Poisson equation, since the system is uniform.
To determine all the relevant microscopic correlation functions in this non-equilibrium system, one must generate the distribution of its electron-hole pair fluctuations. That requires variational analysis of Eq. (2). There are two steps in this, related but treated separately. First we compute the steady-state fluctuation distribution
$$\mathrm{\Delta }f_𝐤k_\mathrm{B}T\frac{\delta f_𝐤}{\delta \mu }=\frac{1}{\mathrm{\Omega }}\underset{𝐤^{}}{}\frac{\delta f_𝐤}{\delta f_𝐤^{}^{\mathrm{eq}}}\mathrm{\Delta }f_𝐤^{}^{\mathrm{eq}},$$
(3)
where the equilibrium fluctuation is the mean-square fluctuation of the occupation number, precisely as defined in statistical mechanics:
$$\mathrm{\Delta }f_𝐤^{\mathrm{eq}}=k_\mathrm{B}T\frac{\delta f_𝐤^{\mathrm{eq}}}{\delta \mu }=f_𝐤^{\mathrm{eq}}(1f_𝐤^{\mathrm{eq}}).$$
(4)
Equation (3) is directly calculable by variation of the one-body Boltzmann transport equation (BTE), and is the exact solution to the linearised BTE (Green and Das 2000b,c).
For every collision model, there exists a unique one-to-one transformation that maps $`\mathrm{\Delta }f^{\mathrm{eq}}`$ to the functional $`\mathrm{\Delta }f`$ in the steady state. Together with the fact that $`\mathrm{\Delta }f`$ exactly satisfies the variational BTE for the electron-hole pair (density-density) correlation function (Kadanoff and Baym 1962), such a mapping establishes Eq. (3) as the unique mathematical form of the mean-square particle fluctuations out of equilibrium. This explicit and crucial connection, between the equilibrium and non-equilibrium fluctuations, is not model-dependent. It is completely generic to the kinetic description of noise and completely absent from every drift-diffusive description (Kogan 1996).
The next step is to obtain the dynamic response. In our physical picture, a spontaneous thermal energy exchange with the heat reservoir sets up an initial electron-hole pair excitation with average strength $`\mathrm{\Delta }f_𝐤^{}`$, at some given initial location $`𝐫^{}`$ within volume $`\mathrm{\Omega }`$. The localised spontaneous excitation is not stable. It relaxes back to the steady state according to the time-dependent Green function, derived when Eq. (2) is perturbed in time, wave-vector space, and (weakly) in real space. This dictates the change in background distribution for state $`𝐤`$ at position $`𝐫`$, given the initial random excitation at $`𝐫^{}`$. Full technical details are in Green and Das (1998a, 2000b,c).
All of the relaxation dynamics are contained in the transient component of the calculable Green function, remaining after the stable long-time adiabatic part is removed from the full response. Let us denote the transient component, in the frequency domain, by $`𝒞_{\mathrm{𝐤𝐤}^{}}(𝐫𝐫^{};\omega )`$. The dynamic (and non-local) two-body response, namely the electron coupled with its hole, is denoted by
$$\mathrm{\Delta }\mathrm{f}_{\mathrm{𝐤𝐤}^{}}^{(2)}(𝐫𝐫^{};\omega )𝒞_{\mathrm{𝐤𝐤}^{}}(𝐫𝐫^{};\omega )\mathrm{\Delta }f_𝐤^{}.$$
(5)
The full velocity-velocity correlation function is, following the approach of Gantsevich et al. (1979)
$$\mathrm{𝐯𝐯}^{}\mathrm{\Delta }\mathrm{f}^{(2)}(𝐫𝐫^{};\omega )_c^{}\stackrel{\mathrm{def}}{=}\frac{2}{\mathrm{\Omega }^2}\underset{𝐤}{}\underset{𝐤^{}}{}𝐯_𝐤\mathrm{Re}\{𝒞_{\mathrm{𝐤𝐤}^{}}(𝐫𝐫^{};\omega )\}𝐯_𝐤^{}\mathrm{\Delta }\mathrm{f}_𝐤^{}.$$
(6)
This completely determines the response of the carrier flux at $`𝐫`$, induced by a spontaneous thermal fluctuation in the carrier flux at $`𝐫^{}`$. This is not a reciprocal process; although uniformity means that the magnitude of the correlator depends on the relative co-ordinate $`𝐫𝐫^{}`$, it matters which of the two positions is upstream. The externally driven system fluctuates asymmetrically in space, just as it fluctuates irreversibly in time.
We now have the basic tool to construct both the thermal noise and the shot noise in the conductor. We discuss the low-frequency case. The thermal spectral density becomes, in a familiar way, the volume integral of all the current-current correlations:
$$S_{\mathrm{therm}}(E)4_\mathrm{\Omega }𝑑𝐫_\mathrm{\Omega }𝑑𝐫^{}(ev_x/L)(ev_x^{}/L)\mathrm{\Delta }\mathrm{f}^{(2)}(𝐫𝐫^{};0)_c^{};$$
(7)
the conductor’s length is $`L`$. The expression is often portrayed in the literature as a real-space symmetrised form which, however, adds little or nothing to the intrinsic physics. We do not elaborate on $`S_{\mathrm{therm}}(E)`$ except to recall two major properties. The first is the Johnson-Nyquist equilibrium limit
$$S_{\mathrm{therm}}(E0)=4Gk_\mathrm{B}T$$
(8)
where, in our Drude model, the conductance becomes $`G=Ne^2\tau /m^{}L^2`$, for effective mass $`m^{}`$.
The second property is that Eq. (7), regardless of driving field, will never scale other than as $`\mathrm{\Delta }f\mathrm{\Delta }f^{\mathrm{eq}}T`$ in a degenerate metallic conductor. This is the strict and inevitable consequence of kinetics, of Fermi-liquid physics, and most of all of asymptotic equilibrium and neutrality in the metallic leads. It has been discussed exhaustively (Das and Green 2000; Green and Das 1998a, 2000a-c).
### B Shot noise
To set the work of Gomila and Reggiani in context, we must look first at the ‘smooth crossover formula’ and its inbuilt theoretical deficiency. True shot noise never scales with temperature; for example, it remains well-defined even in the zero-temperature limit. In purporting to make thermal noise integral with true shot noise, the ‘smooth-crossover formula’ of drift-diffusive theory, Eq. (1), attempts the kinetically impossible. For, every drift-diffusive phenomenology proclaims that Eqs. (1) and (7) are identical (Kogan 1996; Blanter and Büttiker 2000). Yet the rigorous kinetic-theoretical constraint on non-equilibrium thermal noise, namely its abiding proportionality to $`T`$, means that Eq. (7) cannot possibly describe shot noise in the presence of strong degeneracy.
It follows that Eq. (1) is hard to justify. At any rate, the equation’s theoretical claim (thermal noise equals shot noise) is unsustainable by a first-principles analysis. We mean an analysis that is conventionally executed, in keeping with the conventional understanding of statistical mechanics and microscopics (Green and Das 2000a,b). The ‘smooth crossover formula’ is inconsistent, purely and simply.
That is the background to the shot-noise considerations of Gomila and Reggiani (2000). We do not deny that shot-noise-like structure – linear in the current – can emerge from the thermal noise spectrum. Indeed it does, in the semiclassical ballistic limit (Green and Das 1998b; Gomila and Reggiani 2000). It is also the case that purely classical models lead to classical shot noise, $`2eI`$, based on Eq. (7). Nevertheless the leading concern is with metallic diffusive wires. There, $`T`$-independent shot noise cannot be contrived from a strictly thermal basis. Detailed experiments have been proposed to test our claim (Green and Das 2000a, 2000c).
With the knowledge that true shot noise is fundamentally different from thermal noise, one can build an operationally consistent theory for it. We introduce the idea of the response to variations in the total number of carriers transiting the conductor. In the mean, $`N`$ is constant in time but fluctuates by $`\delta N^+=+1`$ at any instant that a carrier first enters at the source. Similarly it changes by $`\delta N^{}=1`$ as the carrier finally exits at the drain (conceptually this amounts to the injection of a hole at the drain). It is not hard to see that this process is described by the non-local correlation
$$𝒞_{\mathrm{𝐤𝐤}^{}}(𝐫𝐫^{};0)\frac{\delta f_𝐤^{}}{\delta N}\delta N^\pm =𝒞_{\mathrm{𝐤𝐤}^{}}(𝐫𝐫^{};0)\frac{\mathrm{\Delta }f_𝐤^{}}{\mathrm{\Delta }N}\delta N^\pm ,$$
(9)
where $`\mathrm{\Delta }N=\mathrm{\Omega }\mathrm{\Delta }f`$. When this correlation records a particle entry, then $`𝐫^{}`$ lies in the cross-sectional region at the source and $`𝐫`$ at the drain. When it records a particle exit, then $`𝐫^{}`$ belongs to the drain area and $`𝐫`$ to the source. Figure 1 illustrates the principle. The essence of shot noise is that it is a time-of-flight process, involving many sporadic transits of carriers across a predefined geometry. On this view shot noise has nothing to do with correlations distributed throughout the volume of a conductor. This is totally unlike the guiding assumption for Eq. (1) (Kogan 1996).
The definition of the measurable shot noise, strictly across the source-drain gap, follows naturally. It simply sums the stochastically independent terms produced, on average, by each of the $`N`$ active contributors. Each contribution is of equal weight since the carriers are all equivalent (assuming temporal randomness of entry/exit). Thus
$$S_{\mathrm{shot}}(E)\stackrel{\mathrm{def}}{=}\mathrm{S}_{\mathrm{s};\mathrm{d}}(\mathrm{E})\delta \mathrm{N}^++\mathrm{S}_{\mathrm{d};\mathrm{s}}(\mathrm{E})\delta \mathrm{N}^{}=\mathrm{S}_{\mathrm{s};\mathrm{d}}(\mathrm{E})\mathrm{S}_{\mathrm{d};\mathrm{s}}(\mathrm{E}),$$
(11)
in which the two directional correlations are
$$S_{\mathrm{s};\mathrm{d}}(E)=2N𝑑𝐫\delta (xL)𝑑𝐫^{}\delta (x^{})\frac{1}{\mathrm{\Delta }N}(ev_x)(ev_x^{})\mathrm{\Delta }\mathrm{f}^{(2)}(𝐫𝐫^{};0)_c^{}$$
(12)
for injection at the source $`(x=0)`$ and, for removal at the drain $`(x=L)`$,
$$S_{\mathrm{d};\mathrm{s}}(E)=2N𝑑𝐫\delta (xL)𝑑𝐫^{}\delta (x^{})\frac{1}{\mathrm{\Delta }N}(ev_x)(ev_x^{})\mathrm{\Delta }\mathrm{f}^{(2)}(𝐫^{}𝐫;0)_c^{}.$$
(13)
Note especially that:
* the space co-ordinates in the argument of $`\mathrm{\Delta }\mathrm{f}_{\mathrm{𝐤𝐤}^{}}^{(2)}(𝐫^{}𝐫;0)`$ are reversed in $`S_{\mathrm{d};\mathrm{s}}(E)`$, and
* $`S_{\mathrm{shot}}(E=0)`$ vanishes, since the equilibrium kinetic equation is self-adjoint (time reversible) and entails the identity $`S_{\mathrm{s};\mathrm{d}}(0)=S_{\mathrm{d};\mathrm{s}}(0)`$.
One easily verifies that Eq. (II B) is explicitly independent of temperature and goes to $`2eI`$ for current $`I`$ in the semi-classical ballistic limit (Green and Das 1998a). Neverthless, the intimate microscopic link with the (conjugate) thermal effects remains. It is manifest in the role of the flux auto-correlation $`v_xv_x^{}\mathrm{\Delta }\mathrm{f}^{(2)}(𝐫𝐫^{};0)_c^{}`$.
Ultimately, a real measurement of current fluctuations detects them in the access leads for the sample. Hence we expect to detect the sum of thermal and shot-noise contributions, if these effects are statistically independent. Non-equilibrium thermal noise will itself carry a hot-electron excess; it is non-dissipative rather than Johnson-Nyquist in origin (Green and Das 2000b,c), just as shot noise is non-dissipative. In the macroscopic limit hot-electron noise goes quadratically with $`E`$, at least in simple cases. One should ask whether this term could overwhelm an emerging shot-noise signal. This is unlikely, as can be seen from a rough estimate based on the Drude model. The bulk thermal-noise excess goes as
$$S_{\mathrm{exs}}(E)=S_0\frac{\mathrm{\Delta }N}{N}\frac{m^{}\mu _\mathrm{e}^2E^2}{k_\mathrm{B}T}4Gm^{}\mu _\mathrm{e}^2E^2$$
(14)
since the ratio $`\mathrm{\Delta }N/N`$ is always less than one in a degenerate conductor. (Here $`\mu _\mathrm{e}=e\tau /m^{}`$ is the mobility.) By finding the upper bound to $`2eI=2e(GEL)S_{\mathrm{exs}}(E)`$ for typical material parameters, one concludes that pure shot noise – if there were no other mechanism to suppress it – would dominate at least up to fields $`10^6\mathrm{Vcm}^1`$ for a sample 1 mm long. At much shorter (mesoscopic) lengths this simple estimate does not hold; specific modelling is needed.
## III Kinetics of shot noise: application
We can now review our results from Green and Das (1998a), built on the form of $`S_{\mathrm{shot}}(E)`$. Despite the relatively crude form of the inelastic collision term in the Drude model, one might expect it to be more relevant at high fields than, say, linear treatments that emphasise coherent (or at least elastic) scattering. A driving potential of a just few volts is quite enough to place a conductor, some millimetres in length, beyond the validity of purely elastic scattering adrift-diffusive theory (Green and Das 2000b).
The calculation is straightforward. The functions $`f_𝐤`$ and $`\mathrm{\Delta }f_𝐤`$ are first obtained for a given field $`E=V/L`$. Then the two-point transient response $`𝒞_{\mathrm{𝐤𝐤}^{}}(𝐫𝐫^{};\omega =0)`$ is derived from the linearised BTE. Finally, all are combined to yield Eq. (II B).
In Fig. 2 we plot the sum of thermal and shot noises in a one-dimensional (1D) wire calculated within the Boltzmann-Drude model of transport, Eq. (2). This is in the strongly degenerate carrier regime, at a thermal energy chosen as $`k_\mathrm{B}T=0.1\epsilon _\mathrm{F}`$, with Fermi energy $`\epsilon _\mathrm{F}`$. The spectral density is normalised to the Johnson-Nyquist value $`S_0`$ and displayed as a function of current in units of $`I_0=2Gk_\mathrm{B}T/e`$. In this and subsequent figures, each curve corresponds to one of five values of the device length as a ratio with the mean free path $`\lambda =\tau v_\mathrm{F}`$ in terms of the Fermi velocity. The curves are always monotonic, tracking down as the ratios rise in the sequence $`L/\lambda =0.001,1,10,25,\mathrm{and}50`$. Typically, $`\lambda `$ is of the order of 50 to 100 nm.
The shortest wire exhibits full shot noise. This is the semi-classical ballistic limit: the absolute upper bound for our model. In the case of the second shortest wire $`L=\lambda `$, the curve falls a little below the ideal value $`1+2eI/S_0`$ in the topmost curve. We note a slight shoulder at $`I2G\epsilon _\mathrm{F}/e`$. The shot noise remains quasi-ballistic because Pauli blocking in the 1D free electron gas efficiently inhibits any scattering when carriers cannot gain enough energy to leave the Fermi sea. This makes the Fermi distribution fairly robust to moderate external fields.
In longer wires, the shot noise at moderate currents is attenuated more and more. Inelastic-scattering suppression is exponential in the Drude model, taking the low-current form
$$\frac{S_{\mathrm{shot}}(E)}{2eI}\left(1+\frac{L}{\lambda }+\frac{L^2}{2\lambda ^2}\right)e^{L/\lambda },$$
(15)
which dies very quickly as the length increases. This accounts for the macroscopic extinction of shot noise.
Figure 3 displays the shot noise of classical 1D carriers, with the same $`n`$ and $`T`$ as the degenerate system of Fig. 2. The mean free path is now $`\tau v_{\mathrm{th}}`$ where $`v_{\mathrm{th}}=(2k_\mathrm{B}T/m^{})^{\frac{1}{2}}`$. We have retained the same physical wire lengths $`L`$ as in Fig. 2 so that, for Fig. 3 specifically, our chosen ratios of length to mean free path are scaled up by $`v_\mathrm{F}/v_{\mathrm{th}}`$. It is significant that the classical curves fall mostly on top of the degenerate ones. An obvious exception is the second plot, where we saw that degeneracy shields the shot noise from attenuation at lower $`I`$. Otherwise the high-current behaviour of the 1D shot noise is unchanged, and thus independent of the carrier statistics. From this it is evident that degeneracy plays no role in the resurgence of high-field shot noise.
To round off our discussion we compare the 1D case with the three-dimensional (3D) case. This is of interest (i) because it is closer to actual experimental devices, and (ii) because it allows us to study the effect of finite, even narrow, wires. Similar results apply in two dimensions (Green and Das 1998a).
Fig. 4 shows spectra for the same sequence of lengths, with a wire radius $`R=100\lambda `$; very wide. There are minor differences with Fig. 2. For instance, in the second curve Pauli blocking is somewhat less effective in preventing attenuation of low-current shot noise.
Rather more interesting are Figs. 5a and 5b. The first is for a wire radius $`R=0.3\lambda `$ and the second for $`R=0.05\lambda `$. (For $`\lambda `$ 100nm, such thicknesses should be achievable by sophisticated nano-lithography.) There is a major loss of shot-noise spectral strength relative to Fig. 4. Recovery towards full noise does not occur before considerably higher current levels are reached. The freedom to explore shot noise through a new variable, the thickness, suggests a novel range of experiments to complement those proposed for shot-noise resurgence as a function of length (Gomila and Reggiani 2000). Such experiments will in any case require exploration of intrinsically high-field, high-current behaviour. This is a region dominated by inelastic collisions, and one that has been almost totally neglected in metallic mesoscopic systems.
Space prevents an extended presentation of our work. We invite readers to examine the complete exposition of our general kinetic approach to noise in Green and Das (2000b,c) as well as the specifics of shot noise in Green and Das (1998a).
Before discussing the work of Gomila and Reggiani (2000), we remark on the striking graphical similarity between their calculation of $`S_{\mathrm{shot}}(E)`$ and the present one. On the other hand, Gomila and Reggiani’s theory for the anomalous rise of macroscopic shot noise apparently depends, in large degree, on the role of long-range Coulomb correlations. In our results, Coulomb effects are completely absent in the sample. Nevertheless, the behaviour of the shot noise is practically the same in our inelastic free-electron model.
## IV The Gomila-Reggiani theory
The recent paper of Gomila and Reggiani (2000) presents a much-needed invitation to reassess the microscopic basis of shot noise (unavoidably, this brings in the questionable status of the ‘smooth crossover’). Certainly Gomila and Reggiani make their argument on the overriding idea of shot noise as a number-fluctuation phenomenon, much in the way of our own philosophy, as already discussed elsewhere (Green and Das 1998a, 2000a). However, there are some differences between the respective approaches. For example, Gomila and Reggiani propose a drift-diffusive model. This entails additional, intuitive assumptions meant to simplify the underlying transport problem. By contrast, we work directly with the semi-classical kinetic equation. The main points of difference are:
* Diffusive method. The drift-diffusive equation of motion embodies a model for current fluctuations in which the diffusion constant $`D`$ (in microscopic terms, a current-current correlator) appears as a simple scaling parameter for the evolution of the current fluctuations themselves. This means that the solution preconditions its own scaling factor $`D`$. If one stays rigidly within the linear-response limit, the Einstein relation can be invoked to constrain the results (Datta 1995; Kogan 1996). This relation is no longer valid for high-field shot noise (Green and Das 2000b). Once pushed out of the weak-field limit, drift-diffusive theories face a complex problem of self-consistency in estimating $`D`$. The problem is highly non-linear and ill-controlled. In the standard Boltzmann approach, there is no a priori distinction between ‘drift’ and ‘diffusion’. The exact non-equilibrium kinetic equation will always be linear in the basic electron-hole pair fluctuations.
* Langevin sources. Gomila and Reggiani follow the popular stratagem of generating single-particle fluctuations only, in a drift-diffusive setting. This is done by adding ad hoc stochastic source terms (Langevin’s Ansatz) to the equation of motion (Kogan 1996). The low-order correlators within the phenomenology must be set by hand to meet the presumed constraint of Einstein’s relation. The status of this low-field stochastic Ansatz is unclear in high-field situations; it is certainly no clearer for carrier populations with strong internal interactions. Such ambiguities arise simply because the imposition of extraneous, stochastic current sources has no physical or logical basis in the microscopics of an internally correlated system (van Kampen 1981). Reliable kinetic descriptions of noise can be set up with no appeal at all to Langevin sources of the kind adopted for the Gomila-Reggiani model. This is as true of non-degenerate noise (Korman and Mayergoyz 1996) as it is for the semi-classical picture of electron-hole polarisation fluctuations in metallic systems (Green and Das 2000b).
* $`T`$-dependence. Equations (12) and (16) of Gomila and Reggiani (2000) both give an overall scaling of their $`S_{\mathrm{shot}}`$ with Johnson-Nyquist noise $`S_0`$, and hence with temperature. The shot noise of Eq. (16) in particular will then be manifestly $`T`$-dependent, unlike true shot noise, unless there is a counterbalancing thermal denominator in the non-linear term giving the shot-noise contribution. Such a factor will cancel the thermal dependence introduced through $`S_0`$. While this is likely to be so in the classical high-$`T`$ limit, where high-field excess noise has little dependence on $`T`$ (Green and Das 1998a, 2000b), it is not clearly so in the degenerate regime of their model. In that limit, the explicit temperature behaviour of the relevant parameters $`L_D`$ and $`I_R`$ is not given. Unless that behaviour is known, one cannot say whether the theory of Gomila and Reggiani recovers true temperature-independent shot noise in bulk metallic wires. A direct check of Eqs. (12) and (16) yields no countervailing factor to undo the $`T`$-dependence entering through $`S_0`$. Hence $`S_{\mathrm{shot}}`$ must scale with $`T`$ in the Gomila-Reggiani model at strong carrier degeneracy.
## V Summary
We have reviewed some prior results for shot noise in degenerate conductors. We conclude that the anomalous recovery of robust shot noise at bulk scales (where it is normally extinguished by inelastic scattering) depends on pushing the system to large enough currents. In such a strongly non-equilibrium limit, the average transit time of a carrier is well given by $`eN/I`$. When this becomes less than the typical scattering time, carriers are ballistic and the high-field shot noise reaches its ideal Schottky value of $`2eI`$.
A simple but strictly kinetic model, informed by a time-of-flight understanding of shot noise, gives a consistent microscopic picture of such noise. The model emphasises plain inelastic scattering in a strongly driven conductor, rather than more subtle and higher-order field effects. On the basis of its clear results, it suggests that Coulomb-fluctuation corrections need not be fundamental to the physics of shot noise in long conductors. Furthermore, recovery of the noise at high currents is not sensitive to quantum statistics. This is because the energy scale for transport will eventually outstrip even a large Fermi energy.
We have demonstrated the resurgence of true temperature-independent shot noise at high currents, even in long thin 3D samples. For a set level of the current, shot noise should certainly be much stronger in samples that are relatively wider as well as shorter. With this prediction, namely the inhibition of high-field shot noise in a constricted geometry, we advance an altogether different and major opportunity for new experiments.
We share one common idea with the approach to shot noise proposed by Gomila and Reggiani. It is the importance of (necessarily discrete) number fluctuations as generators of shot noise. This departs from thermal noise, whose character derives from distributed and continuous random changes in the carriers’ free energy. In a real sense, we are brought right back to basic thermodynamics. That is because thermal and shot-noise fluctuations are echoes of the thermodynamic conjugacy of the equilibrium variations $`\delta \mu /k_\mathrm{B}T`$ and $`\delta \mathrm{ln}N`$.
Aside from the crucial difference over the significance of Coulomb corrections in the resurgence of bulk shot noise, there are major differences of method between Green and Das (1998a) and Gomila and Reggiani (2000). The main one is our systematic adherence to strict Boltzmannian kinetics and Fermi-liquid theory (Green and Das 2000b). Decidedly, this sets our investigations apart from each one of the drift-diffusive (or so-called Boltzmann-Langevin) works to be found in the noise literature. That includes Gomila and Reggiani’s.
Boltzmann-Langevin phenomenology relies (a) on fictitious stochastic sources – said to generate the individual, one-body current fluctuations – and (b) on essentially classical diffusion to evolve such single-particle objects. Neither (a) nor, worse, (b) makes any identifiable connection with the electron-hole pair symmetry that is essential to the actual make-up of the microscopic correlations. It is remarkable that this electron-hole asymmetry within diffusively driven transport models (Büttiker 1986) persists all the way up to the macroscopic scale of the device leads (Fenton 1994). Because this unbalanced behaviour is built right into the asymptotic boundary conditions, it is intrinsic to all drift-diffusive descriptions. They cannot be rid of it.
Each of the assumptions above is equally unfounded when it comes to charged Fermi liquids at the level of microscopic many-body physics (Green and Das 2000b). The price of their apparent phenomenological simplicity is nothing less than an unphysical electron-hole asymmetry. The consequence of that is the failure of drift-diffusive fluctuations to recover the correct electronic compressibility. Hence they also fail to meet a most basic condition: metals cannot sustain inhomogeneous electric fields beyond the Thomas-Fermi screening length. This rule is violated by every diffusively based noise model (Das and Green 2000).
Langevin stochastics and pseudo-classical diffusion fail to conform to orthodox quantum kinetics for noise in charge transport, despite frequent claims to the contrary (Kogan 1996). Such schemes are deeply foreign to the real nature of a degenerate, polarisable electron plasma. There, electron-hole pair dynamics, the conservation laws, and the sum rules are utterly central to the physics (Pines and Nozières 1966). The drift-diffusive theories’ demonstrable violation of the sum rules, and neglect of the conservation laws embodied in those rules, provides the plainest evidence of non-conformity (Green and Das 2000a,b).
In passing we have called attention, once again, to the problematic status of the ‘smooth crossover’ proclaimed by drift-diffusive phenomenology. We suggest that the time is ripe for a thorough microscopic reassessment of this largely intuitive construct, and most of all for a renewed search for decisive experimental tests of it.
New information on the ‘smooth crossover’ can be expected in two experimental contexts. The first is in the low-field noise signal from quantum-confined devices (Green and Das 2000a,c). The second is the high-field regime as indicated by us (Green and Das 1998a, 2000b,c), by Naveh to some extent (Naveh 1998) and, latterly, by Gomila and Reggiani (2000). In all situations it is important to have a unifying microscopic description, equipped to cover all of the many facets of real shot noise. This would be the only way to make the most of any fresh experimental knowledge.
## references
Blanter, Ya. M., and Büttiker, M. (2000). Phys. Rept. 336, 1.
Büttiker. M. (1986). Phys. Rev. Lett. 57, 1761.
Das, M. P., and Green, F. (2000). ‘Proceedings of the 23rd International Workshop on Condensed Matter Theories, Ithaca, Greece’, ed. G. S. Anagnostatos (Nova Science: New York) in press; Das, M. P., and Green, F. (1999). Preprint cond-mat/9910183.
Datta, S. (1995). ‘Electronic Transport in Mesoscopic Systems’ (Cambridge University Press: Cambridge).
Fenton, E. W. (1994). Superlattices and Microstructures 16, 87.
Gantsevich S., V., Gurevich, V. L., and Katilius, R. (1979). Nuovo Cimento 2, 1.
Gillespie, D. T. (2000). J. Phys.: Cond. Matter 12, 4195..
Gomila, G., and Reggiani, L. (2000). Phys. Rev. B 62, 8068. See also Preprint cond-mat/0005094.
Green, F., and Das, M. P. (1998a). Preprint cond-mat/9809339. (CSIRO-RPP3911: unpublished.)
Green, F., and Das, M. P. (1998b). ‘Recent Progress in Many-Body Theories’, ed. D. Neilson and R. F. Bishop, p. 102 (World Scientific: Singapore). See also Green, F, and Das, M. P. (1997). Preprint cond-mat/9709142.
Green, F., and Das, M. P. (2000a). ‘Proceedings of the Second International Conference on Unsolved Problems of Noise and Fluctuations (UPoN’99)’, ed. D. Abbott and L. B. Kish AIP 511, pp 422-33 (American Institute of Physics: New York). For a similar discussion see Green, F., and Das, M. P. (1999). Preprint cond-mat/9905086.
Green, F., and Das, M. P. (2000b). J. Phys.: Cond. Matter 12, 5233. See also Green, F., and Das, M. P. (2000). Preprint cond-mat/0001412.
Green, F., and Das, M. P. (2000c). J. Phys.: Cond. Matter 12, 5251. See also Green, F., and Das, M. P. (1999). Preprint cond-mat/9911251.
Kadanoff, L. P., and Baym G. (1962). ‘Quantum Statistical Mechanics’ (W A Benjamin, Reading, Massachusetts).
Kogan, Sh. M. (1996). ‘Electronic Noise and Fluctuations in Solids’ (Cambridge University Press: Cambridge).
Korman, C. E., and Mayergoyz, I. D. (1996). Phys. Rev. B 54, 17620.
Naveh, Y. (1998). Preprint cond-mat/9806348.
Pines, D., and Nozières, P. (1966). ‘The Theory of Quantum Liquids’ (Benjamin, New York).
van Kampen, N. G. (1981). ‘Stochastic Processes in Physics and Chemistry’ (North-Holland, Amsterdam), pp 246-52.
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# 1 INTRODUCTION AND RIEMANNIAN MANIFOLDS
## 1 INTRODUCTION AND RIEMANNIAN MANIFOLDS
The problem of geodesic connectedness in semi-Riemannian manifolds (i.e. the question whether each two points of the manifold can be joined by a geodesic) has been widely studied from very different viewpoints. Our purpose is to review these semi-Riemannian techniques, and possible extensions. In the Riemannian case, it is natural to state this problem on (incomplete) manifolds with (possibly non-smooth) boundary, and we will discuss different conditions on this boundary in the remainder of this Section. In this case, geometrical as well as variational methods are appliable, and accurate results can be obtained by using the associated distance and related properties of positive-definiteness. For Lorentzian manifolds, the results cannot be so general, and very different techniques have been introduced which are satisfactory for some particular classes of Lorentzian manifolds. We start by considering several geometrical notions appliable to affine manifolds and, thus, to all semi-Riemannian manifolds, Section 2. Recall that variational methods are not appliable to affine manifolds, at least in a standard way: geodesics are the critical points of the action functional, but a metric tensor must be provided for the definition of this functional. In Section 3 some general facts on geodesic connectedness of Lorentzian manifolds are pointed out, and classical results about connectedness of spaceforms, which rely in the properties of actions of isommetry groups, are summarized. In Section 4 we discuss variational methods applied to Lorentzian manifolds, which have been shown to be useful, mainly, to study stationary and splitting manifolds, with or without boundary. Finally, in Section 5 recent results, based on topological arguments and appliable to multiwarped spacetimes, are explained.
As pointed out by Gordon , the problem of the geodesic connectedness of a Riemannian manifold is important not only in its own right but also because of its relation, via the Jacobi metric, with the problem of connecting the points of the manifold by means of the trajectories determined by an autonomous potential. Moreover, it is also important because stronger results on geodesic connectedness among all semi-Riemannian manifolds are found for definite metrics. In fact, the following concept is especially useful: a Riemannian metric will be said to be convex if each two of its points can be joined by means of a distance-minimizing geodesic (not necessarily unique). By the Hopf-Rinow theorem all complete Riemannian manifolds are convex, and we discuss now when an incomplete one is either convex or geodesically connected.
We start with the simplest case. Let $`(,,)`$ be a $`n`$-dimensional Riemannian manifold (all the manifolds will be assumed to be connected and smooth, even though at most differentiability $`C^4`$ will be needed), and $`𝒟`$ an open (connected) domain with differentiable (smooth) boundary $`𝒟`$; put $`\stackrel{~}{𝒟}=𝒟𝒟`$. We will not assume a priori that $``$ is complete, because this assumption is not especially relevant for the indefinite case. The following definitions become natural:
(I) $`𝒟`$ is infinitesimally convex at $`p𝒟`$ (IC<sub>p</sub>) if the second fundamental form $`\sigma _p`$, with respect to the interior normal, is positive semidefinite. When $`\sigma _p`$ is positive definite we say that $`𝒟`$ is strictly locally convex at $`p`$ (SIC<sub>p</sub>). Equivalently, take any differentiable function $`\varphi :U\stackrel{~}{𝒟}\text{}`$, where $`U`$ is a neighborhood of $`p`$ such that
$$\varphi ^1(0)=U𝒟,\varphi >0\text{on }U𝒟\text{ }\text{and}d\varphi _q0\text{for any }qU𝒟.$$
(1.1)
Then, $`𝒟`$ is IC<sub>p</sub> (resp. SIC<sub>p</sub>) if and only if for one (and hence for all) function $`\varphi `$ satisfying (1.1):
$$H_\varphi (p)[v,v]0(\text{resp.}<0)vT_p𝒟.$$
(1.2)
(II) $`𝒟`$ is locally convex at $`p𝒟`$ (LC<sub>p</sub>) if there exists a neighborhood $`U`$ of $`p`$ such that
$$\mathrm{exp}_p\left(T_p𝒟\right)\left(U𝒟\right)=\mathrm{}.$$
(1.3)
When $`\mathrm{exp}_p\left(T_p𝒟\right)\left(U\stackrel{~}{𝒟}\right)=\{p\}`$, then $`𝒟`$ is strictly locally convex at $`p`$ (SLC<sub>p</sub>).
It is easy to check
$$IC_pLC_p,SIC_pSLC_p.$$
(1.4)
Clearly, the converse to the first implication is not true ($`=\text{}^2,𝒟=\{(x,y):y>x^3\},p=(0,0)`$). Nevertheless, if there exists a neighborhood $`U`$ of $`p`$ such that $`𝒟`$ is IC<sub>q</sub> for all $`qU𝒟`$ then $`𝒟`$ is $`LC_q`$ for all $`qU𝒟`$. Do Carmo and Warner realized that this problem is not as trivial as it sounds, and solved it (as a step for other computations) for the constant curvature case. Bishop solved it in general, even though differentiability $`C^4`$ is explicitly used in his technique. Nor does the converse to the last implication (1.4) hold ($`=\text{}^2,𝒟=\{(x,y):y>x^4\},p=(0,0)`$).
The previous definitions are appliable to each point $`p`$ in the boundary $`𝒟`$. The following definitions are appliable to all $`𝒟`$:
(1) Infinitesimally convex (IC): $`𝒟`$ is IC<sub>p</sub> for all $`p𝒟`$. This is equivalent to being:
(1A) locally convex (LC): LC<sub>p</sub>, $`p𝒟`$ (because of Bishop’s result).
(1B) variationally convex (VC): for one (and hence for all) function $`\varphi :\stackrel{~}{𝒟}\text{}`$ satisfying (1.1) with $`U=`$, inequality (1.2) holds, $`p𝒟`$. For the equivalence between this definition and IC, it is enough to prove the following straightforward property \[3, Cap. 3\]: for any differentiable manifold with boundary $`\stackrel{~}{𝒟}`$, a function $`\varphi `$ satisfying (1.1) with $`U=`$ exists. (VC has been widely used by using variational methods, because function $`\varphi `$ allows direct penalization of the action functional (1.7).)
The “strict” concepts SIC, SLC and SVC can be defined analogously, and, clearly: $`SICSVCSLC`$ (but not the converse).
(2) Geometrically convex (GC): for any $`p,q𝒟`$, the range of any geodesic $`\gamma :[0,1]\stackrel{~}{𝒟}`$ such that $`\gamma (0)=p,\gamma (1)=q`$ satisfies
$$\gamma \left([0,1]\right)𝒟.$$
(1.5)
If this also holds when $`p,q𝒟`$ the boundary is strictly geometrically convex (SGC). GC is a straightforward generalization to Riemannian manifolds of the usual notion of convexity given in Euclidean spaces. It is straightforward to check the implications $`LCGCIC`$; thus, by Bishop’s result, GC is equivalent to IC. Moreover, $`SGCSLC`$. It is worth pointing out that, in the complete case, Germinario \[27, Theorem 2.1\], obtained by variational methods (with a weaker assumption on differentiability than Bishop), a direct proof of $`VCGC`$.
(3) (Geodesically) pseudoconvex (PC) (in the sense of , see also ): for each compact set $`K𝒟`$ there is a compact set $`H𝒟`$ such that each geodesic segment $`\gamma :[a,b]𝒟`$ with $`\gamma (a),\gamma (b)K`$ satisfies $`\gamma ([a,b])H`$. This definition is intrinsic to the (open domain of the) manifold. So, it is said that $`𝒟`$ (rather than $`𝒟`$) is PC (thus, it is natural to assume $`𝒟=`$).
Easily, $`PCGC`$. Nevertheless, PC is not implied by GC; in fact, a complete Riemannian manifold may be non-PC (for example: a complete surface with infinitely many holes). Summing up: at each point $`p`$ implications (1.4) holds and, globally
$$PCGCLCVCIC.$$
It is easy to check that, when $`\stackrel{~}{𝒟}`$ is complete, there is no loss of generality assuming that so is $``$: otherwise, a new Riemannian metric on all $``$ can be defined such that it is complete and coincides with the original metric on $`\stackrel{~}{𝒟}`$ \[3, Cap. 3\]. All the above equivalent conditions on convexity for $`𝒟`$ provide different ways to prove that:
when $``$ is complete, $`𝒟`$ is convex if and only if $`𝒟`$ is convex.
Recall that the completeness assumption is essential. In fact, if $``$ is incomplete and geodesically connected, a domain $`𝒟`$ with convex boundary may be non-geodesically connected (take as $``$ a cylinder minus a small segment, and as $`𝒟`$ a small ball such that the segment lies inside it).
Now we are ready to examine the general case where $`𝒟`$ is not differentiable or $`\stackrel{~}{𝒟}`$ is not complete, which has been systematically studied in . By the quoted result, it is clear that if there exists a sequence $`\left(\stackrel{~}{𝒟}_m\right),m𝐍`$ of complete submanifolds with convex (differentiable) boundary such that
$$\stackrel{~}{𝒟}_m\stackrel{~}{𝒟}_{m+1}\text{and}𝒟=\underset{m𝐍}{}\stackrel{~}{𝒟}_m,$$
(1.6)
then $`𝒟`$ is geodesically connected (as a consequence, the results by Gordon in are reproven and generalized in a simple way). Now, it is natural to wonder if, in this case, $`𝒟`$ must be convex. The answer is quite simple: if $`\stackrel{~}{𝒟}`$ is complete then $`𝒟`$ is convex; otherwise, there are counterexamples. This result is proven by using standard geometrical methods .
Now, two questions arises naturally: (A) can the convexity of each $`𝒟_m`$ be weakened? and (B) is there any intrinsic condition for the convexity of $`𝒟`$ (independent of its boundary in $``$)? To answer these questions the (intrinsic) Cauchy boundary of $`𝒟`$ and a variational approach must be taken into account. Let $`\stackrel{~}{𝒟}^c`$ be the canonical completation of $`𝒟`$ by using Cauchy sequences, and $`_c𝒟`$ the corresponding boundary points, $`\stackrel{~}{𝒟}^c=𝒟_c𝒟`$. Recall that $`\stackrel{~}{𝒟}^c`$ is always complete as a metric space, but the boundary points in $`_c𝒟`$ are not necessarily differentiable and, if they are, the metric may be non–extendible or degenerate there. Note that any point of $`𝒟`$ determines naturally one or more points in $`_c𝒟`$, and $`\stackrel{~}{𝒟}`$ is complete if and only if all the points in $`_c𝒟`$ are determined in this way by points of $`𝒟`$. From the variational point of view, geodesics joining two fixed points $`p,q𝒟`$ are seen as critical points of the action functional
$$f(x)=\frac{1}{2}_0^1\dot{x}(s),\dot{x}(s)𝑑s$$
(1.7)
defined on the space of differentiable curves joining $`p`$ and $`q`$ or, technically better, on its $`H^1`$ Sobolev completion, i.e., the Hilbert manifold $`\mathrm{\Omega }^1(𝒟,p,q)=\{xH^{1,2}([0,1],𝒟)x(0)=p,x(1)=q\}`$.
A first result shows that a domain $`𝒟`$ as in (1.6) is convex, even if the boundary of each $`𝒟_m`$ is not convex (and, thus, each $`𝒟_m`$ may be non–geodesically connected), when a suitable local estimate of the loss of convexity of $`𝒟_m`$ and boundness of the sequence is ensured.
Theorem 1.1. Let $`(,,)`$ be a complete Riemannian manifold and $`𝒟`$ an open domain of $``$. Assume that there exists a positive differentiable function $`\varphi `$ on $`𝒟`$ such that
(i) $`lim_{x𝒟}\varphi (x)=0`$;
(ii) each $`y𝒟`$ admits a neighbourhood $`U`$ and constants $`a,b>0`$ such that
$`a\varphi (x)bx𝒟U;`$
(iii) the first and second derivatives of the normalized flow of $`\varphi `$ are locally bounded close to $`𝒟`$ (that is: each $`y𝒟`$ admits a neighbourhood $`U`$ such that the induced local flow of $`\varphi /\varphi ^2`$ on $`𝒟U`$ have first and second derivatives with bounded norms), and
(iv) there exists a decreasing and infinitesimal sequence $`\{a_m\}`$ such that each $`y𝒟`$ admits a neighbourhood $`U`$ and a constant $`M\text{}`$ satisfying:
$$H_\varphi (x)[v,v]Mv,v\varphi (x)x\varphi ^1(a_m)U,vT_x\varphi ^1(a_m),m𝐍.$$
(1.8)
Then $`𝒟`$ is convex.
It is straightforward to check that, if $`𝒟`$ is differentiable and convex, all the conditions (i)–(iv) are automatically satisfied.
To prove this result, the action functional $`f`$ is penalized with a term depending on a positive parameter $`ϵ`$: $`f_ϵ(x)=f(x)+ϵ_0^1\varphi ^2(x(s))𝑑s`$. Each penalized functional is bounded from below and satisfies the condition of Palais-Smale; so, from standard variational arguments (see, for example, ) $`f_ϵ`$ admits a critical (minimum) point. The crucial point is to prove that a critical point of $`f_ϵ`$ in a sublevel of $`f_ϵ`$ is uniformly far (with respect to $`ϵ`$) from $`𝒟`$. In the proof the critical points of the penalized functionals are projected (using the normalized flow of $`\varphi `$) on the hypersurface $`\varphi ^1(a_m)`$ for $`m`$ large enough. This makes possible to get critical points of $`f`$ (i.e. geodesics) not touching $`𝒟`$ by means of a limit process.
Technical condition (iii) and even the completeness of the ambient manifold $``$ can be weakened if (iv) is imposed on all points and directions enough close to the boundary. So, a straightforward consequence of the technique in for Theorem 1.1 is the following result (compare with ):
Theorem 1.2. Let $`(,,)`$ be a Riemannian manifold, $`𝒟`$ an open domain, and $`\stackrel{~}{𝒟}^c=𝒟_c𝒟`$ its canonical Cauchy completion. Assume that there exists a positive differentiable function $`\varphi `$ on $`𝒟`$ such that
(i) $`lim_{x_c𝒟}\varphi (x)=0`$;
(ii) each $`y_c𝒟`$ admits a neighbourhood $`U\stackrel{~}{𝒟}^c`$ and constants $`a,b>0`$ such that
$`a\varphi (x)bx𝒟U;`$
(iii) each $`y_c𝒟`$ admits a neighbourhood $`U\stackrel{~}{𝒟}^c`$ and a constant $`M\text{}`$ such that inequality (1.8) holds for all $`x𝒟U`$ and for all $`vT_x`$.
Then $`𝒟`$ is convex.
It is worth pointing out that in both previous theorems a multiplicity result can be also obtained, when the topology of the fiber is not homotopically trivial: if $`𝒟`$ is not contractible (in itself), then for any $`p,q𝒟`$ there exists a sequence $`\{x_m\}`$ of geodesics in $`𝒟`$ joining them such that $`lim_m\mathrm{}f(x_m)=\mathrm{}.`$ For the proof one uses that the Ljusternik–Schnirelman category of $`𝒟`$ is infinite, which implies the existence of infinitely many connecting critical points of $`f_ϵ`$ with diverging lengths. We refer to for a deeper discussion of these results, and for examples, (see also ).
## 2 AFFINE CONNECTIONS
Now, consider a manifold $``$ endowed with an affine connection $``$; thus, all the properties of its geodesics will hold for all semi-Riemannian manifolds. As we are interested in geodesics, there is no loss of generality assuming that $``$ is symmetric; otherwise, it is well-known that there exists another affine connection with the same geodesics and torsion-free.
All previous notions about the convexity of the (differentiable) boundary of a domain $`𝒟`$ are directly extendible to the affine case except IC<sub>p</sub>, IC, because the second fundamental form $`\sigma _p`$ is not canonically defined. Nevertheless, the characterization of IC<sub>p</sub> in terms of the function $`\varphi `$ satisfying (1.1) does makes sense. In fact, it is not difficult to check that $`(\text{1.2})`$ still holds for a function $`\varphi `$ satisfying (1.1) if and only if it holds for each such $`\varphi `$. Thus, this will be the natural definition of IC<sub>p</sub> for affine manifolds. In principle, Bishop’s result holds just in the Riemannian case and, so, the implications which hold are (1.4) and: PC or LC $``$ GC $``$ VC $``$ IC, and SLC $``$ SGC $``$ SVC $``$ SIC.
However, these concepts (except for PC) are not so useful in the affine case, because even when $`\stackrel{~}{𝒟}`$ is complete and $`𝒟`$ is convex in the strongest sense, $`𝒟`$ may be non-geodesically connected. In fact, the first problem to be solved is what conditions on a manifold without boundary must be imposed to obtain geodesic connectedness. The following example by Bates shows that a compact and complete affine manifold may be non-geodesically connected.
Example 2.1. Consider the moving frame $`(X_1=\mathrm{cos}x_x+\mathrm{sin}x_y,X_2=\mathrm{sin}x_x+\mathrm{cos}x_y)`$ on $`\text{}^2`$, and the affine connection $``$ such that $`X_1,X_2`$ are parallel. The geodesics of $``$ are the integral curves of the (complete) vector fields $`X=aX_1+bX_2,a,b\text{}`$. So, any geodesic $`\gamma (s)=(x(s),y(s))`$ is complete and $`x(s)`$ lies in an interval of length $`2\pi `$. Thus, a required example is the quotient torus $`T^2=\text{}^2/4\pi \text{}^2`$, with the induced connection.
Given $`p`$ let Con$`(p)`$ be the subset containing the points of $``$ which can be connected with $`p`$ by means of a geodesic. The previous example also shows that, even under strong hypotheses, Con$`(p)`$ may be non-closed. To study this more in depth, consider the following concepts , .
Let $`G()`$ be the space of the geodesics of $`(,)`$, that is, the projective bundle $`P`$ (obtained by identifying two vectors of the reduced bundle $`T^{}=T\backslash \{\text{zero section}\}`$ if they are proportional) where two clases of non-zero vectors $`[v],[w]P`$ are identified if there exists a geodesic $`\gamma :]a,b[`$ such that $`\gamma ^{}(t_0)[v],\gamma ^{}(t_1)[w]`$ for some $`t_0,t_1]a,b[`$. Thus, each geodesic $`\gamma `$ determines a unique class $`[\gamma ]G()`$. In the remainder of this Section, all the geodesics will be assumed inextendible. The natural quotient topology of $`G()`$ can be characterized as follows. Consider first a sequence of geodesics $`\gamma _n:]a_n,b_n[`$ and a geodesic $`\gamma :]a,b[`$. If there exists $`t_0]a,b[`$ contained in all but a finite number of $`]a_n,b_n[`$, and $`\{\gamma _n^{}(t_0)\}\gamma ^{}(t_0)`$, then lim sup$`\{a_n\}a<b`$lim inf$`\{b_n\}`$ and $`\{\gamma _n^{}\}`$ converges uniformly on compact subsets of $`]a,b[`$ to $`\gamma ^{}`$ (for any distance on $`T`$ compatible with its topology) \[41, Prop. 2.1\]. Now, a sequence of geodesics $`\{\beta _n\}`$ is said to converge tangentially to a geodesic $`\beta `$ if there exists $`\gamma _n[\beta _n],\gamma [\beta ]`$ and a $`t_0`$ such that $`\{\gamma _n^{}(t_0)\}\gamma ^{}(t_0)`$; this convergence holds if and only if $`\{[\beta _n]\}[\beta ]`$ in $`G()`$, .
It is worth pointing out:
(I) Tangential convergence is completely independent of completeness, even if $``$ is compact. In fact, even if $`\gamma _n`$ converges tangentially to a unique $`\gamma `$, all $`\gamma _n`$ may be incomplete and $`\gamma `$ complete, or viceversa; counterexamples can be found in Lorentzian tori .
(II) The sequence $`\gamma _n`$ may converge tangentially to more than one limit and, in this case, $`G()`$ is not Hausdorff. Remarkably, this happens when a point $`p`$ of any affine manifold is removed, if not all the geodesics starting at $`p`$ are closed. But, unfortunately, this is not by any means the only case. In fact, consider the standard flat torus $`\text{}^2/\text{}^2`$; the sequence of geodesics constantly equal to $`\alpha `$, where $`\alpha `$ is induced on the torus by a geodesic in $`\text{}^2`$ with irrational slope, has infinitely many tangential limits.
Recall that $`(,)`$ is called disprisoning if given any geodesic $`\gamma :]a,b[`$ and any compact subset $`K`$ of $``$ there are sequences $`\{t_n\}a^+`$, $`\{s_n\}b^{}`$ such that $`\gamma (t_n)`$ and $`\gamma (s_n)`$ do not lie in $`K`$. Disprisonment, pseudoconvexity and the topology of $`G()`$ have proven to be fruitful in order to study some geometrical properties, including a Cartan-Hadamard type theorem and applications to Relativity, see , , . For geodesic connectedness, the following result by Beem and Parker holds , .
Theorem 2.2. Let $`(,)`$ be an affine manifold; the following implications are fulfilled:
$`\text{Disprisonment and pseudoconvexity}G()\text{is Hausdorff}\text{Con}(p)\text{is closed, }p.`$
Moreover, if this last property holds and there are no conjugate points, then $`(,)`$ is geodesically connected.
The last assertion is obvious because the absence of conjugate points implies that Con$`(p)`$ is open for all $`p`$; thus, the two implications in Theorem 2.2 yield sufficient conditions for geodesic connectedness.
## 3 THE INDEFINITE SEMI-RIEMANNIAN CASE. SPACEFORMS
We refer to the standard books , for definitions and general background about semi-Riemannian, and especially Lorentzian, manifolds. For indefinite metrics, the absence of an associated canonical distance and, so, of any analog to the Hopf-Rinow theorem, makes the problem of geodesical connectedness very subtle. Perhaps the only non-trivial result with a clear resemblance to the Riemannian case is that of Avez and Seifert : in a globally hyperbolic Lorentzian manifold, any pair of causally related points (i.e. which can be joined with a causal curve) can be joined by means of a causal geodesic. For this result is essential that, in the Lorentzian case, causal geodesics maximize locally the “time-separation” (or “Lorentzian distance”) between causally related points. Global hyperbolicity introduces then a sort of compactness in the space of causal curves $`𝒞_{p,q}^{cau}`$ joining any two fixed points $`p,q`$, in such a way that the lengths of curves in $`𝒞_{p,q}^{cau}`$ are bounded, and its supremum at every connected part of $`𝒞_{p,q}^{cau}`$ is reached by a curve (necessarily pregeodesic) therein.
Nevertheless, neither compactness nor completeness implies geodesic connectedness. It is interesting to study the geodesic connectedness of Lorentzian surfaces (Smith, ), in comparison with the affine case. Recall first that a plane $`𝒮`$ endowed with a Lorentzian metric (or conformal class of metrics) is called normal if there exists a diffeomorphism of $`𝒮`$ onto $`\text{}^2`$ which takes every null-geodesic into an axis-parallel line; this can be characterized in terms of the absence of barriers (see also ). Moreover, a null-complete plane is normal if its Gaussian curvature does not change sign off a compact subset, and the integral of its curvature is finite.
Theorem 3.1. (1) A normal Lorentzian plane is geodesically connected . (2) The universal covering of a complete Lorentzian torus with a Killing vector field $`K(0)`$ is normal .
Thus, a complete Lorentzian torus with such a $`K`$ is geodesically connected (this can be also checked more directly, ); the completeness in assertion (2) can be replaced by any of the following conditions (a posteriori equivalent): (i) $`K`$ has a definite causal sense on all the torus (timelike or null or spacelike), or (ii) the metric is (globally) conformally flat. Now, consider the following Lorentzian torus \[52, Sect. 5\] (see also \[10, Sect. 3\], \[48, Sect. 5\]).
Example 3.2. Take on $`\text{}^2`$ the moving frame $`X_1,X_2`$ as in Example 2.1, and consider the Lorentzian metric $`g`$ such that $`X_1`$ and $`X_2`$ are null and $`g(X_1,X_2)=1`$. The velocities of any timelike or spacelike curve $`(x(s),y(s))`$ must remain in any of the four cones continuously determined by $`X_1,X_2`$; so, the length of the projection of $`x(s)`$ is again bounded. Thus, a geodesically disconnected Lorentzian torus is induced.
It is worth pointing out that in this example $`K=_y`$ is a Killing vector field (in Example 2.1 $`K`$ is an affine vector field) inducible on the torus. So, this Lorentzian torus is geodesically incomplete (the causal character of $`K`$ changes); we do not know any example of complete and geodesically disconnected torus.
The fact that a complete semi-Riemannian manifold may be geodesically disconnected can be stressed studying spaceforms. We say that a semi-Riemannian $`n`$-dimensional manifold $``$ of index $`\nu `$ is a spaceform if it is complete with constant curvature $`C`$. In this case, $``$ is covered by the corresponding model (1-connected) space $`M(n,\nu ,C)`$, that is, $`=M(n,\nu ,C)/\mathrm{\Gamma }`$, where $`\mathrm{\Gamma }`$ is the fundamental group of $``$. The case $`C=0`$ is trivial (the model space $`M(n,\nu ,0)\text{}_\nu ^n`$ is geodesically connected) and, up to a homothety, we can assume $`C=1`$ (the homothetic factor may be positive as well as negative; recall that the Levi-Civita connection remains unchanged). If $`n3`$, the model space is then the pseudosphere $`\text{𝕊}_\nu ^n`$ (spacelike vectors of norm 1 in $`\text{}_\nu ^{n+1}`$); the Lorentzian pseudosphere $`\text{𝕊}_1^n`$ is also called de Sitter spacetime, and it is globally hyperbolic. Recall that an affine manifold is called starshaped from a point $`p`$ if $`\mathrm{exp}_p`$ is onto. It is not difficult to check \[39, Propos. 5.38\] that no indefinite pseudosphere $`\text{𝕊}_\nu ^n`$, $`0<\nu <n`$ is geodesically connected. The following result by Calabi and Markus , in particular, solves completely the geodesic connectedness of Lorentzian spaceforms of positive curvature with $`n3`$.
Theorem 3.3. For $`n2`$:
(1) Two points $`p,q\text{𝕊}_1^n`$ are connectable by a geodesic if and only if $`p,q_1>1`$, where $`,_1`$ is the usual Lorentzian inner product of $`\text{𝕃}^{n+1}\text{}_1^{n+1}`$.
(2) Every spaceform $`=S_1^n/\mathrm{\Gamma },S_1^n`$ is starshaped from some point $`p`$.
(3) A spaceform $`=S_1^n/\mathrm{\Gamma }`$ is geodesically connected if and only if it is not time-orientable.
The proof of (1) follows by a direct computation of the geodesics. For the remainder, it is essential that, whenever $`2\nu n`$, the group of isommetries $`\mathrm{\Gamma }`$ is finite. Then, up to conjugacy, $`\mathrm{\Gamma }O(1)\times O(n)O_1(n+1)`$, and the proof follows by studying the barycenter of the orbits, which must lie in the timelike axis of $`\text{}_1^{n+1}`$.
For arbitrary index $`\nu `$ (including the case $`\nu =n1`$, which is equivalent to the Lorentzian case of constant negative curvature) the only general results we know are extensions of Theorem 3.3, with $`n3`$, and assuming as an additional hypothesis (when $`2\nu >n`$) that the fundamental group $`\mathrm{\Gamma }`$ is finite. For these extensions, the non time-orientability must be replaced by the inexistence of a proper time-axis \[54, Theor. 11.2.3\]. We recall that a time-axis $`T`$ is a one-dimensional $`\mathrm{\Gamma }`$-invariant negative-definite linear subspace of $`\text{}_\nu ^{n+1}`$. $`T`$ is proper if $`\mathrm{\Gamma }`$ acts trivially on $`T`$; that is, if $`T`$ is not proper then $`\mathrm{\Gamma }`$ also acts as a multiplication by -1.
## 4 VARIATIONAL METHODS.
The systematic application of variational methods on infinite-dimensional manifolds to Lorentzian Geometry started with a seminal paper by Benci and Fortunato (see also ), who studied geodesic connectedness of standard stationary spacetimes. Since then, the geodesic connectedness of stationary as well as splitting spacetimes, with or without boundary, has been widely studied by variational methods. We will review briefly these results, explaining mainly the stationary case, and giving some references and comments for the splitting case (see also the next Section). We refer to for general background on variational methods and applications to this and other problems in Lorentzian Geometry, and to for properties of Killing fields which will be borne in mind.
For stationary manifolds we will follow the approach in , because it makes more intrinsic assumptions and recovers the previous results. Recall that a Lorentzian manifold $`(,,)`$ is called stationary if it admits a globally defined timelike Killing vector field $`K`$. Fixing two points $`p,q`$, geodesics joining them are, still in the semi-Riemannian case, the critical points of the action functional $`f`$ on $`\mathrm{\Omega }^1(,p,q)`$ given in (1.7). But now this functional is bounded neither from above nor from below, and it may not satisfy Palais-Smale condition. Nevertheless, this problem can be solved by taking into account that any critical point of $`f`$ (or geodesic) $`z`$ satisfies $`K,z^{}C_z`$, where $`C_z`$ is a constant. This suggests that variations in the $`K`$–direction are irrelevant, and that, under the orthogonal splitting of the tangent bundle $`TM=\text{Span}(K)\text{Span}(K)^{}`$, only projections on the spatial part $`\text{Span}(K)^{}`$ are important (similar ideas work for geodesic completeness , ).
More precisely, let $`𝒞_{p,q}=\{zC^1([0,1],)z(0)=p,z(1)=q,K,z^{}C_z\}`$ and $`𝒩_{p,q}`$ its Sobolev $`H^{1,2}`$-completation; for technical computations, the Riemannian metric $`g_R`$ obtained by reversing $`,`$ on Span$`(K)`$ can be used. The action functional $`f`$ can be written as the sum of two functionals $`f_1,f_2`$,
$$f_1(z)=\frac{1}{2}_0^1\left(\dot{z}(s),\dot{z}(s)\frac{\dot{x}(s),K^2}{K,K}\right)𝑑s,f_2(z)=\frac{1}{2}_0^1\frac{\dot{z}(s),K^2}{K,K}𝑑s.$$
It is not difficult to check that $`f_1^{}`$ vanishes on all tangent vectors $`\zeta `$ to $`\mathrm{\Omega }^1(,p,q)`$ which are pointwise parallel to $`K`$, that is, $`f_1`$ is the “spatial” part with respect to $`K`$ of the functional $`f`$. Moreover, $`𝒩_{p,q}`$ can be characterized as the set of $`z\mathrm{\Omega }^1(,p,q)`$ such that $`f^{}(z)[\zeta ]=f_2^{}(z)[\zeta ]=0`$, for any such $`\zeta `$ parallel to $`K`$; by the implicit function theorem, $`𝒩_{p,q}`$ is a submanifold of $`\mathrm{\Omega }^1(,p,q)`$. Finally, one can check that critical points of $`f`$ on $`\mathrm{\Omega }^1(,p,q)`$ coincide with the critical points of the restriction $`J`$ of $`f`$ to $`𝒩_{p,q}`$ (and thus to $`𝒞_{p,q}`$).
Summing up, if some conditions are imposed on $`𝒞_{p,q}`$ so that the restriction $`J`$ of $`f`$ to $`𝒩_{p,q}`$ must reach a critical point, then a geodesic connecting $`p`$ and $`q`$ is obtained. And the most natural variational such conditions are:
> $`𝒞_{p,q}`$ is (i) non-empty and (ii) $`c`$precompact, for some $`c>\mathrm{Inf}_{𝒞_{p,q}}f(z)`$
($`𝒞_{p,q}`$ is said $`c`$precompact if every sequence $`\{z_n\}𝒞_{p,q}`$ with $`f(z_n)c`$ has a uniformly convergent subsequence in $``$). Thus, if (i) and (ii) hold for all $`p,q`$, then $`(,,)`$ is geodesically connected (one can check that $`J`$ is bounded from below, its sublevels $`J^c^{}`$ are complete metric subspaces of $`𝒩_{p,q}`$, for all $`c^{}c`$, and Palais-Smale condition is fulfilled).
Now, we can wonder when (i) and (ii) hold. If the restriction of $`f`$ to $`𝒞_{p,q}`$ is pseudocoercive (that is, $`c`$precompact for all $`c\mathrm{Inf}_{𝒞_{p,q}}f(z)`$) for any $`p,q`$ then $`(,,)`$ is globally hyperbolic, but the converse is not true. Condition (i) holds if either $`p`$ and $`q`$ are causally related or $`K`$ is complete (this happens, for example, when the auxiliary Riemannian metric $`g_R`$ is complete, and $`g_R(K,K)(=K,K)`$ is bounded); clearly, the converse does not hold. Assume that $`(,,)`$ is a standard stationary spacetime, that is, $``$ is a product manifold $`=\text{}\times _0`$ ($`_0`$ any manifold) and $`,`$ can be written, with natural identifications, as:
$$,=\beta dt^2+2\omega dt+g_0,$$
(4.1)
where $`dt^2`$ is the usual metric on , and $`g_0,\beta ,\omega `$ are, resp., a Riemannian metric, a positive function and a 1-form, all on $`_0`$ (locally, stationary spacetimes look like standard stationary ones).
Theorem 4.1. A standard stationary spacetime is geodesically connected, if: (a) $`g_0`$ is complete, (b) $`0<`$Inf$`(\beta )`$ Sup$`(\beta )<\mathrm{}`$, and (c) the $`g_0`$-norm of $`\omega (x)`$ has a sublinear growth in $`_0`$.
(For (c), we mean that the norm of $`\omega (x)`$ has an upper bound $`Ad_0(x,p_0)^\alpha +B`$, for some $`A,B\text{},\alpha [0,1[,p_0_0`$, where $`d_0`$ is the $`g_0`$-distance on $`_0`$). In fact, under these three conditions, assumptions (i) and (ii) are always satisfied; we refer to \[32, Prop. A.3\] for a more intrinsic way to express these conditions, in terms of stationary manifolds admitting a differentiable time function (which are standard stationary a posteriori). In the standard static case ($`\omega 0`$) condition $`0<`$Inf$`(\beta )`$ can be dropped (see , ); however, we remark that the imposed inequalities always imply global hyperbolicity \[50, Cor. 3.4, 3.5\].
It is also worth pointing out:
(I) This technique provides also consequences for the existence of infinitely many connecting geodesics or timelike geodesics when $``$ is not contractible. In fact,
two points $`p,q`$ of a stationary manifold can be joined by a sequence of spacelike geodesics with divergent lengths if $`K`$ is complete, $`𝒞_{p,q}`$ is pseudocoercive and $``$ is non-contractible;
(the essential step for the proof is that $`\mathrm{\Omega }_{p,q}(M)`$ is homotopically equivalent to $`𝒩_{p,q}`$ and, thus, the Ljusternik-Schnirelman category of $`𝒩_{p,q}`$ is infinite).
For timelike geodesics, recall that, under our type of assumptions, Avez-Seifert’s technique is appliable and chronologically related points can be joined by timelike geodesics. Let $`p,q`$, and $`\gamma _q(t)`$ be an integral curve of $`K`$ starting at $`q`$; when $`K`$ is complete then $`p`$ belongs to the chronological past of $`\gamma _q(t)`$ for $`t`$ big enough (a direct proof is not difficult, see also \[50, Sect. 4\]). Then, under precompactness, $`𝒩_{p,\gamma _q(t)}`$ contains at least a timelike geodesic for $`t`$ big enough and, when the topology is not homotopically trivial, the following result on multiplicity holds (see \[32, Theorem 1.4\]):
If $``$ is non-contractible, $`K`$ is complete and there exist $`c_0<0,t_0>0`$ such that $`𝒩_{p,\gamma _q(t)}`$ is $`c_0`$ precompact for all $`t>t_0`$, then the number of timelike geodesics joining $`p`$ and $`\gamma _q(t)`$ goes to $`\mathrm{}`$ when $`t\mathrm{}`$.
(II) Analogous techniques should work if a semi-Riemannian manifold $`(,g)`$ of index $`s`$ admits $`s`$ Killing vector fields $`K_1,\mathrm{},K_s`$ independent at each $`p`$ such that $`g`$ restricted to $`𝒦=`$ Span$`\{K_1,\mathrm{},K_s\}`$ is negative definite. We have now the natural splitting $`TM=𝒦𝒦^{}`$, and, for each geodesic $`z`$, the projection of $`\dot{z}`$ on $`𝒦`$ can be recovered from the constants $`C_{z,i}g(\dot{z},K_i)`$ (the analogous problem for geodesic completeness was solved in ).
Moreover, even in the Lorentzian case, one can consider the case when there exist two pointwise independent Killing vector fields $`K_1,K_2`$ such that $`\{K_1(p),K_2(p)\}`$ spans a Lorentzian plane at each $`p`$ but neither $`K_1`$ nor $`K_2`$ are timelike on (all) $``$. Remarkably, this happens in Gödel type spacetimes; for the modifications of the technique in this case, see , .
(III) Let us discuss the case with boundary briefly (see also, , , ). Consider first a standard stationary spacetime $`=\text{}\times _0`$ satisfying the assumptions of Theorem 4.1 (thus, geodesically connected), and let $`𝒟_0_0`$ be an open domain with differentiable boundary $`𝒟_0`$. We have seen that, in general, a domain $`𝒟`$ of a complete or geodesically connected semi-Riemannian manifold does not inherit good properties for geodesic connectedness. Nevertheless, if $`𝒟=\text{}\times 𝒟_0`$ is variationally convex (VC), then the corresponding function $`\mathrm{\Phi }`$ can be chosen independent of $`t`$, and the functionals $`f`$, $`J`$, can be penalized in a similar way to the Riemannian case. So, penalized functionals $`J_ϵ`$ do satisfy Palais-Smale condition, and one obtains geodesic connectedness. Moreover, one obtains again that $`𝒟`$ is VC if and only if it is GC . Of course, if we are interested just in $`𝒟`$, it is not exactly relevant for $``$ to fulfill assumptions of Theorem 4.1: one needs just the convexity of $`𝒟`$ and the possibility to extend the metric on $`𝒟`$ to all $``$, in such a way that the assumptions of Theorem 4.1 are fulfilled; this possibility can be expressed as more intrinsic conditions on $`𝒟`$. Furthermore, the assumption that the stationary manifold is standard can also be dropped. In fact, in order to penalize the functionals $`f,f_1,f_2`$ as in the Riemannian case, one needs only that the boundary $`𝒟`$ can be determined by a function $`\varphi `$ as in the definition of VC, which is also invariant by the flow of $`K`$.
When the boundary is not smooth, the problem has been studied in the standard static case ($`\omega 0`$). A result in the spirit of Theorem 1.2 can be obtained, proving the geodesic connectedness of spacetimes such as outer Reissner-Nordström’s and Schwarzschild’s (see also \[3, Cap. 6\]). On the other hand, for the standard static case, a different variational principle in solves completely the problem of connecting a point and a integral curve of $`_t`$, for arbitrary $`_0`$ (or $`𝒟_0`$), .
We will mean by a splitting spacetime a product manifold $`=\text{}\times _0`$ endowed with a Lorentzian metric as (4.1) but allowing $`g_0,\beta ,`$ and $`\omega `$ to depend differentiably on the time variable $`t`$. In this case, bounds on the derivatives with respect to $`t`$ of these elements must be also imposed in order to obtain geodesic connectedness. As a typical result we have :
Theorem 4.2. A splitting spacetime $`(,,=\beta (t,x)dt^2+2\omega (t,x)dt+g_t(x))`$ is geodesically connected, if:
(1) $`g_0`$ is complete and there exists $`\lambda >0`$ such that $`g_t>\lambda g_0`$ for all $`t`$.
(2) $`0<`$Inf$`(\beta )`$ and $`\beta (x,0),\omega (x,0)`$ are bounded.
(3) $`g_t/\beta (t,x)`$ and $`\omega /\beta (t,x)`$ are bounded by a function on $``$ type: $`b_0(x)+b_1(x)|t|^\mu ,`$ ($`\mu [0,1[`$, and $`\mu [0,2[`$, resp.)
(4) Consider the natural derivatives $`_t\alpha ,_t\beta ,_t\delta `$ of $`\alpha ,\beta ,\delta `$ with respect to $`t`$, resp. Then the $`g_t`$–norms of $`_t\alpha /\alpha ,_t\beta /\beta ,_t\delta `$ are bounded at each hypersurface with constant $`t`$, and its supremum when $`t\pm \mathrm{}`$ goes to 0.
Moreover, domains of type $`𝒟=]a,b[\times _0`$ (strips) are shown to inherit geodesic connectedness, provided that $`𝒟`$ is VC (see for orthogonal splittings $`\omega 0`$, for non-orthogonal splittings, and for strips; see also ).
Nevertheless, now the $`t`$-dependence does not allow a reduction to an equivalent Riemannian problem, as in the stationary case. Instead, Rabinowitz’s saddle point theorem is used, but two technical complications must be circumvented: (A) $`f`$ does not satisfy a Palais-Smale condition, which is solved by approximating by a family of functionals $`f_\eta ,\eta 0,f_0=f`$, and making some a priori estimates of the critical points of $`f_\eta ,\eta >0`$ in order to ensure a good behavior under the limit $`\eta 0`$, and (B) in Rabinowitz’s theorem, the independent directions where the functional goes to $`\mathrm{}`$ are finite; so, a Galerkin finite-dimensional approximation is carried out. For the existence of infinitely many connecting geodesics when $``$ is not contractible, the relative category, a topological invariant somewhat subtler then the Ljusternik-Schnirelman category, is used.
Finally, it is worth pointing out that the problem of geodesic connectedness is naturally generalized to others like: (1) the connectedness of two submanifolds by normal geodesics, studied for splitting manifolds in , or (2) the connection by trajectories of some more general Lagrangian systems, studied for stationary manifolds and potential vector fields independent of time in .
## 5 MULTIWARPED SPACETIMES. A TOPOLOGICAL METHOD
A multiwarped spacetime is a product manifold $`I\times F_1\times \mathrm{}\times F_m,`$ $`I=]a,b[\text{}`$ endowed with a metric $`g=dt^2+_{i=1}^mf_i^2(t)g_i,tI`$, where $`f_1,\mathrm{}f_m`$ are positive functions on $`I`$, and each $`g_i`$ is a Riemannian metric on the manifold $`F_i`$. These spacetimes include classical examples of spacetimes: when $`m=1`$ they are the Generalized Robertson-Walker (GRW) spacetimes, standard models of inflationary spacetimes ; when $`m=2`$, the intermediate zone of Reissner-Norsdström spacetime and the interior of Schwarzschild spacetime appear as particular cases ; moreover, multiwarped spacetimes may also represent relativistic spacetimes together with internal spaces attached at each point (see ).
The geodesic connectedness of this type of spacetimes with $`m=2`$ have been studied by using variational methods in manifolds without and with boundary , . Nevertheless, more accurate results are obtained by using a topological method introduced in ; in fact, the following result is proven:
Theorem 5.1. (1) In a multiwarped spacetime with convex fibers $`(F_1,g_1),\mathrm{},(F_m,g_m)`$ each two causally related points can be joined by a causal geodesic.
(2) The multiwarped spacetime is geodesically connected if the fibers are and:
$$\begin{array}{cc}_c^bf_i^2(f_1^2+\mathrm{}+f_m^2)^{1/2}=\mathrm{}\hfill & _a^cf_i^2(f_1^2+\mathrm{}+f_m^2)^{1/2}=\mathrm{}\hfill \end{array}$$
(5.1)
for all $`i`$ and for some $`c(a,b)`$. Moreover, if one of the fibers $`F_j`$ is not contractible then:
(a) each two points can be joined by infinitely many geodesics, and
(b) for any $`zI\times F_1\times \mathrm{}\times F_m,`$ and $`xF_1\times \mathrm{}\times F_m,`$ the number of timelike geodesics joining $`z`$ and $`(t,x)`$ goes to $`\mathrm{}`$ when $`t`$ goes to an endpoint of the interval $`I`$.
It is worth pointing out:
(I) Equality (5.1) is equivalent to the following condition: any point $`z_0`$ of the spacetime can be joined with any line $`L[x]`$ by means of both, a future directed and a past directed causal curve. Nevertheless, Theorem 5.1 does not cover all the possibilities of the technique, and more general versions of this theorem can be given. These general versions give very accurate results; in fact, a necessary and sufficient condition for geodesic connectedness when $`m=1`$ can be given with a reasonably long distinction of cases (multiplicity, existence of timelike geodesics, conjugate points, etc. are also completely characterized in this reference; see also ). In particular, geodesic connectedness of Reissner-Nordström Intermediate spacetime is reproven (previous proofs and results in this direction were obtained in , ). The accuracy of the technique is shown by proving the geodesic connectedness of Schwarzschild inner spacetime; in fact, a good behaviour of the warping functions yields geodesic connectedness, but the warping functions of Schwarzschild inner spacetime do not have such good behaviour. Nevertheless, in this case, this problem can be skipped because one of the fibers of Schwarzschild’s is a sphere (and so, each pair of its points can be joined by geodesics of arbitrarily large length); if this fiber is replaced by a plane, the resulting spacetime is not geodesically connected.
(II) If the fibers are assumed to have boundary, the problem is reduced to considering when this boundary implies geodesic connectedness or convexity (Section 1). The general technique also works when a strip $`]a^{},b^{}[I`$ is considered; so, the problem with boundary is also solved, for boundaries which preserve the multiwarped structure.
For the proof of Theorem 5.1, the Avez-Seifert type result (1) relies on a partial integration of the geodesics. For (2), the idea is the following, under an assumption somewhat stronger than (5.1) (under (5.1) some technicalities must be also taken into account). As there are points in each $`L[x]`$ both, future and past related with $`z`$, these points can be joined with causal geodesics. Thus, we have just to connect $`z=(z_0,z_1,\mathrm{}z_m)`$ and $`(t,x),x=(x_1,\mathrm{}x_m)`$ for $`t`$ in a compact interval. For this: (A) fixing the geodesics in the fibers joining each $`z_i`$ and $`x_i`$, a continuous map $`\overline{\mu }(c,K)`$, $`\overline{\mu }:]0,1[^{m1}\times [K^{},K^+]\text{}^{m1}`$, is constructed in such a way that each zero of $`\overline{\mu }`$ represents the initial condition of a connecting geodesic, (B) for $`K=K^{}`$ and $`K=K^+`$, the result above on causal geodesics ensures the existence of at least one zero of (a continuous extension of) $`\overline{\mu }`$ for some $`c[0,1]^{m1}`$; then, the problem is solved if a continuous set of zeros $`𝒵`$ containing these two zeros is found, (C) for each $`K]K^{},K^+[`$ the behavior of $`\overline{\mu }(c,K)`$ at the boundary of $`[0,1]^{m1}`$, allows to find $`𝒵`$, by means of arguments on continuity of solutions of equations depending on a parameter, based in Brower’s topological degree. When $`F_j`$ is not contractible, the result follows by applying the technique for each geodesic joining $`z_j`$ and $`x_j`$.
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# Untitled Document
RI-02-00, EFI-2000-7, CERN-TH/2000-095
hep-th/0005052
D-Branes in the Background of NS Fivebranes
Shmuel Elitzur<sup>1</sup>, Amit Giveon<sup>1,2</sup>, David Kutasov<sup>3</sup>, Eliezer Rabinovici<sup>1</sup> and Gor Sarkissian<sup>1</sup>
<sup>1</sup>Racah Institute of Physics, The Hebrew University
Jerusalem 91904, Israel
<sup>2</sup>Theory Division, CERN
CH-1211, Geneva 23, Switzerland
<sup>3</sup>Department of Physics, University of Chicago
5640 S. Ellis Av., Chicago, IL 60637, USA
We study the dynamics of $`D`$-branes in the near-horizon geometry of $`NS`$ fivebranes. This leads to a holographically dual description of the physics of $`D`$-branes ending on and/or intersecting $`NS5`$-branes. We use it to verify some properties of such $`D`$-branes which were deduced indirectly in the past, and discuss some instabilities of non-supersymmetric brane configurations. Our construction also describes vacua of Little String Theory which are dual to open plus closed string theory in asymptotically linear dilaton spacetimes.
5/00
1. Introduction
In the last few years it was found that embedding various supersymmetric gauge theories in string theory, as the low energy worldvolume dynamics on branes, provides an efficient tool for studying many aspects of the vacuum structure and properties of BPS states in these theories. One class of constructions (reviewed in ) involves systems of $`D`$-branes ending on and/or intersecting $`NS5`$-branes. For applications to gauge theory one is typically interested in taking the weak coupling limit $`g_s0`$ (as well as the low energy limit). In this limit one might expect the system to be amenable to a perturbative worldsheet treatment, but the presence of the $`NS5`$-branes and branes ending on branes complicate the analysis.
In the absence of a derivation of the properties of $`D`$-branes interacting with $`NS5`$-branes from first principles, some of their low energy properties were postulated in the past based on symmetry considerations and consistency conditions. One of the main purposes of this paper is to derive some of these properties by a direct worldsheet study of $`D`$-branes in the vicinity of $`NS5`$-branes.
In the analysis, we will use the improved understanding of the dynamics of $`NS5`$-branes achieved in the last few years. It is now believed that in the weak coupling limit (but not necessarily at low energies) fivebranes decouple from gravity and other bulk string modes and give rise to a rich non-gravitational theory, Little String Theory (LST) . In it was proposed to study LST using holography. The theory on a stack of $`NS5`$-branes was argued to be holographically dual to string theory in the near-horizon geometry of the fivebranes. Many properties of LST can be understood by performing computations in string theory in this geometry. In particular, we will find below that it is an efficient way to study properties of $`D`$-branes ending on or intersecting $`NS5`$-branes.
From the point of view of LST, $`D`$-branes in the vicinity of $`NS5`$-branes give rise to new particle and extended object states in the theory, as well as new vacua, typically with reduced supersymmetry (when the branes are space-filling). Another motivation of this work is to understand the non-perturbative spectrum and dynamics of extended objects in LST, and more generally study the interplay between the non-trivial dynamics on fivebranes and the worldvolume physics of $`D`$-branes in their vicinity.
The plan of the paper is the following. In section 2 we briefly review some facts regarding the relevant brane configurations and, in particular, describe the conjectured properties of these configurations that we will try to verify. Section 3 is a review of the near-horizon geometry of $`NS5`$-branes (the CHS geometry ) and its holographic relation to LST. We also describe a modification of this geometry corresponding to fivebranes positioned at equal distances around a circle, which plays a role in the analysis. In section 4 we review some facts about $`D`$-branes in flat space and on a three-sphere (or $`SU(2)`$ WZW model). In section 5 we study $`D`$-branes in the CHS geometry as well as its regularized version. We verify some of the properties described in section 2, and describe additional features which follow from our analysis. Some of the technical details are presented in the appendices.
2. Some properties of brane configurations
One class of brane constructions, that gave rise to many insights into gauge dynamics, involves $`D`$-branes suspended between $`NS5`$-branes. We start with a brief description of a particular example of such a construction, which realizes four dimensional $`N=1`$ supersymmetric gauge theory with gauge group $`G=U(N_c)`$ and $`N_f`$ chiral superfields in the fundamental representation (more precisely, $`N_F`$ fundamentals $`Q^i`$, $`i=1,\mathrm{},N_F`$, and $`N_F`$ anti-fundamentals $`\stackrel{~}{Q}_i`$). This will help introduce the issues that will be discussed later<sup>1</sup> For a more detailed discussion, references to the original literature and other constructions, see ..
We will consider brane configurations in type IIA string theory consisting of $`NS5`$-branes, $`D4`$ and $`D6`$-branes, oriented as follows<sup>2</sup> In later sections we will discuss more general configurations, obtained by rotating some of the branes.:
$$\begin{array}{cc}\hfill NS5:& (x^0,x^1,x^2,x^3,x^4,x^5)\hfill \\ \hfill NS5^{}:& (x^0,x^1,x^2,x^3,x^8,x^9)\hfill \\ \hfill D4:& (x^0,x^1,x^2,x^3,x^6)\hfill \\ \hfill D6:& (x^0,x^1,x^2,x^3,x^7,x^8,x^9)\hfill \\ \hfill D6^{}:& (x^0,x^1,x^2,x^3,x^4,x^5,x^7)\hfill \end{array}$$
One can check that any combination of two or more of these branes preserves four or eight supercharges. The Lorentz symmetry is broken in the presence of the branes to
$$SO(9,1)SO(3,1)_{0123}\times SO(2)_{45}\times SO(2)_{89}$$
(assuming that all branes that are pointlike in $`(x^4,x^5)`$ are placed at the same point in the $`(4,5)`$ plane, and similarly for $`(x^8,x^9)`$). To study situations with non-trivial $`3+1`$ dimensional physics (in $`(x^0,x^1,x^2,x^3)`$) some of the branes must be made finite.
Consider a configuration of $`N_c`$ $`D4`$-branes stretched between an $`NS5`$-brane and an $`NS5^{}`$-brane separated by a distance $`L`$ along the $`x^6`$ direction (fig. 1). At distances greater than $`L`$, the five dimensional theory on the $`D4`$-branes reduces to a four dimensional theory with $`N=1`$ supersymmetry. The boundary conditions provided by the fivebranes imply that out of all the massless degrees of freedom on the fourbranes (described by $`44`$ strings), only the $`4d`$ $`N=1`$ vector multiplets for $`G=U(N_c)`$ survive. To decouple the gauge dynamics from the complications of string theory one takes the limit $`g_s0`$, $`L/l_s0`$, with the four dimensional gauge coupling $`g^2=g_sl_s/L`$ held fixed .
Fig. 1
One can add matter in the fundamental representation of $`G`$ in one of the two ways illustrated in fig. 2 (which, as explained in , are related). One is to place $`N_f`$ $`D6`$-branes between the fivebranes (fig. 2a). The $`46`$ strings give $`N_f`$ fundamentals of $`U(N_c)`$, $`Q^i`$, $`\stackrel{~}{Q}_i`$, whose mass is proportional to the separation between the sixbranes and the fourbranes in $`(x^4,x^5)`$. Alternatively, one can add to the configuration $`N_f`$ $`D4`$-branes stretching from the $`NS5`$-brane to infinity (fig. 2b). In this case, the $`N_f`$ fundamentals of $`U(N_c)`$ should arise from $`44`$ strings stretched between the two kinds of fourbranes. Their mass is the separation between the fourbranes in $`(x^4,x^5)`$.
Fig. 2
In the configuration of fig. 2a, the low energy physics should be independent of the locations of the $`D6`$-branes along the interval between the fivebranes. In it was pointed out that a particularly natural location for the sixbranes is at the same value of $`(x^4,x^5,x^6)`$ as the $`NS5^{}`$-brane. As is clear from (2.1), in this case the $`NS5^{}`$-brane is embedded in the $`D6`$-branes; it divides them into two disconnected parts (fig. 3).
Fig. 3
Consequently, the configuration has two separate $`U(N_f)`$ symmetries, acting on the two semi-infinite sixbranes, and one may attempt to interpret them as the $`U(N_f)\times U(N_f)`$ global symmetry of $`N=1`$ SQCD, under which $`Q`$ and $`\stackrel{~}{Q}`$ transform as $`(N_F,1)`$ and $`(1,\overline{N}_F)`$, respectively. This would imply that $`46`$ strings connecting the $`D4`$-branes to the upper part of the $`D6`$-branes give rise at low energies to the $`N_f`$ chiral multiplets in the fundamental of $`U(N_c)`$, $`Q^i`$, while those that connect the fourbranes to the lower part of the sixbranes give the $`N_f`$ chiral superfields in the anti-fundamental of $`U(N_c)`$, $`\stackrel{~}{Q}_i`$, as indicated in fig. 3.
Despite the fact that the two groups of sixbranes are independent, none of them can be removed from the configuration. From the string theory point of view the reason is that this would lead to violation of the RR charge that couples to $`D6`$-branes. In the low energy gauge theory the inconsistency is seen as a non-vanishing chiral anomaly. However, the same basic mechanism for getting chiral matter can be used to produce chiral spectra in lower dimensions or in the presence of orientifolds \[9,,10,,11\].
The low-lying excitations of the brane configurations discussed above can be divided into two classes: those that are bound to one of the fivebranes, and those that are not. In this paper we will analyze the properties of the first class of excitations. It includes the following:
(a) $`46`$ strings connecting $`N_c`$ $`D4`$-branes to $`N_F`$ $`D6^{}`$-branes, all of which end on an $`NS5`$-brane<sup>3</sup> The configuration of fig. 4a can be obtained from that of fig. 3 by exchanging $`(x^4,x^5)(x^8,x^9)`$ and removing some branes. (fig. 4a). The prediction is that they give rise to a chiral spectrum: a chiral superfield $`Q`$ in the $`(N_c,N_F)`$ of $`U(N_c)\times U(N_F)`$.
(b) $`44`$ strings connecting fourbranes ending on an $`NS5`$-brane from opposite sides (fig. 4b). They should give rise to chiral superfields, $`Q`$ in the $`(N_c,N_F)`$ and $`\stackrel{~}{Q}`$ in the $`(\overline{N}_c,\overline{N}_F)`$, or hypermultiplets $`(Q,\stackrel{~}{Q})`$.
(c) $`46`$ strings connecting $`D4`$-branes ending on $`NS5`$-branes to $`D6`$-branes intersecting the fivebranes (fig. 4c). They should also give rise to hypermultiplets $`(Q,\stackrel{~}{Q})`$.
We will verify the predictions (a), (b), (c) below. Excitations which belong to the second class and which we will not analyze include the $`44`$ strings that give rise to the $`N=1`$ vector superfield in fig. 1, and the $`46`$ strings that give rise to the fundamental chiral superfields $`Q^i`$, $`\stackrel{~}{Q}_i`$ in fig. 2a. The former are complicated to study, since their wavefunction is spread throughout the interval between the fivebranes and is influenced by them only via the boundary conditions they provide. The latter can be studied using standard $`D`$-brane techniques, at least when the intersection of the $`D4`$ and $`D6`$-branes is far enough from the edges of the fourbrane (see section 4). When that intersection lies on the $`NS5`$-brane (fig. 4c), one can study it using the techniques of this paper, and we will discuss this case below.
Fig. 4
3. The near-horizon geometry of $`NS5`$-branes and holography
The background fields around a stack of $`k`$ parallel $`NS5`$-branes are :
$$\begin{array}{cc}& e^{2(\mathrm{\Phi }\mathrm{\Phi }_0)}=1+\underset{j=1}{\overset{k}{}}\frac{l_s^2}{|\stackrel{}{x}\stackrel{}{x}_j|^2}\hfill \\ & G_{IJ}=e^{2(\mathrm{\Phi }\mathrm{\Phi }_0)}\delta _{IJ}\hfill \\ & G_{\mu \nu }=\eta _{\mu \nu }\hfill \\ & H_{IJK}=ϵ_{IJKL}^L\mathrm{\Phi }\hfill \end{array}$$
$`I,J,K,L=6,7,8,9`$ label the directions transverse to the fivebranes (see (2.1)). $`\mu ,\nu =0,1,\mathrm{},5`$ are the directions along the brane. $`\{\stackrel{}{x}_j\}`$ are the locations of the fivebranes in $`\stackrel{}{x}=(x^6,\mathrm{},x^9)`$. $`H`$ is the field strength of the NS-NS $`B`$-field; $`G`$, $`\mathrm{\Phi }`$ are the metric and dilaton, respectively.
The background (3.1) interpolates between flat ten dimensional spacetime far from the fivebranes, and a near-horizon region, in which the $`1`$ on the right hand side of the first line of (3.1) can be neglected (fig. 5). This near-horizon region is an asymptotically linear dilaton solution. E.g. if the fivebranes are coincident, $`\stackrel{}{x}_j=0`$, the near-horizon solution is
$$\begin{array}{cc}& e^{2(\mathrm{\Phi }\mathrm{\Phi }_0)}=\frac{kl_s^2}{|\stackrel{}{x}|^2}\hfill \\ & G_{IJ}=e^{2(\mathrm{\Phi }\mathrm{\Phi }_0)}\delta _{IJ}\hfill \\ & G_{\mu \nu }=\eta _{\mu \nu }\hfill \\ & H_{IJK}=ϵ_{IJKM}^M\mathrm{\Phi }\hfill \end{array}$$
Fig. 5
String propagation in the near-horizon geometry (3.1) can be described by an exact worldsheet Conformal Field Theory (CFT) . The target space is
$$\mathrm{IR}^{5,1}\times \mathrm{IR}_\varphi \times SU(2)$$
$`\mathrm{IR}_\varphi `$ corresponds to the radial direction $`r=|\stackrel{}{x}|`$:
$$\begin{array}{cc}& \varphi =\frac{1}{Q}\mathrm{log}\frac{|\stackrel{}{x}|^2}{kl_s^2}\hfill \\ & \mathrm{\Phi }=\frac{Q}{2}\varphi \hfill \end{array}$$
where we set $`\mathrm{\Phi }_0=0`$ by rescaling $`\stackrel{}{x}`$. $`Q`$ is related to the number of fivebranes via
$$Q=\sqrt{\frac{2}{k}}$$
The CFT describing the three-sphere at constant $`|\stackrel{}{x}|`$ in (3.1) is an $`SU(2)`$ WZW model at level $`k`$. The $`SU(2)`$ group element $`g`$ is related to the coordinates on the three-sphere via (e.g.):
$$g(\stackrel{}{x})=\frac{1}{|\stackrel{}{x}|}\left[x^61+i(x^8\sigma _1+x^9\sigma _2+x^7\sigma _3)\right]$$
The $`SO(4)SU(2)_L\times SU(2)_R`$ global symmetry corresponding to rotations in the $`\mathrm{IR}^4`$ labeled by $`(x^6,x^7,x^8,x^9)`$ acts on $`g`$ as $`gh_Lgh_R`$, where $`h_{L(R)}SU(2)_{L(R)}`$. Denoting the generators of $`SU(2)_L`$ ($`SU(2)_R`$) by $`J^a`$ ($`\overline{J}^a`$), one finds that $`J^3\overline{J}^3`$ generates rotations in the $`(x^6,x^7)`$ plane, while $`J^3+\overline{J}^3`$ is the generator of rotations in $`(x^8,x^9)`$.
Since we are dealing with the superstring, we are interested in the $`N=1`$ superconformal $`\sigma `$-model on (3.1). Thus, in addition to the bosonic coordinates $`(x^\mu ,\varphi ,g)`$, there are worldsheet fermions ($`\psi ^\mu ,\chi ^r,\chi ^a`$), $`a=1,2,3`$, which are free (after a certain chiral rotation). The worldsheet $`N=1`$ superconformal generators are:
$$\begin{array}{cc}\hfill T(z)=& \frac{1}{2}(x^\mu )^2\frac{1}{2}\psi ^\mu \psi _\mu \frac{1}{2}(\varphi )^2\frac{Q}{2}^2\varphi \frac{1}{2}\chi ^r\chi ^r\frac{1}{k}J^aJ^a\frac{1}{2}\chi ^a\chi ^a\hfill \\ \hfill G(z)=& i\psi _\mu x^\mu +i\chi _r\varphi +iQ(\chi _aJ^ai\chi _1\chi _2\chi _3+\chi _r)\hfill \end{array}$$
Here $`J^a`$ are the bosonic $`SU(2)`$ currents of level $`k_Bk2`$. The total $`SU(2)`$ current algebra of level $`k`$ is generated by the currents $`J_{\mathrm{total}}^a=J^a+J_F^a`$, where $`J_F^a=(i/2)ϵ^{abc}\chi _b\chi _c`$ is the contribution of the fermions. Note, in particular, that this construction only makes sense for $`k_B0`$, i.e. for two or more fivebranes. One also imposes a chiral GSO projection $`()^{F_L}=()^{F_R}=1`$. The GSO projected theory on the background (3.1), (3.1) preserves sixteen supercharges – the $`NS5`$-brane is a half BPS object.
An interesting feature of the near-horizon geometry (3.1) is that in the vicinity of the fivebranes, $`|\stackrel{}{x}|0`$, an infinite “throat” appears (for two or more fivebranes), corresponding to $`\mathrm{IR}_\varphi `$ in (3.1). In it was proposed to interpret it in terms of holography. String theory in the background (3.1) was conjectured to be equivalent to the theory on the fivebranes (LST).
The map between the “bulk” ($`10d`$) and “boundary” ($`6d`$) theories is the following. On-shell observables in the bulk theory, such as vertex operators in the background (3.1) which correspond to non-normalizable wavefunctions supported at the “boundary,” $`\varphi \mathrm{}`$ on $`\mathrm{IR}_\varphi `$, are mapped to off-shell observables in the $`6d`$ fivebrane theory. Non-normalizable vertex operators on (3.1) depend among other things on the six dimensional momentum $`k^\mu `$, which is interpreted as off-shell momentum in the $`6d`$ LST.
Normalizable eigenstates of the Hamiltonian on (3.1) correspond to on-shell states in LST. One way to compute the spectrum of these states is to study correlation functions of non-normalizable vertex operators and look for singularities as a function of $`k_\mu `$. Poles at $`k_\mu ^2=M_i^2`$ signal the presence of normalizable states with that mass in the spectrum of the theory.
The worldsheet theory (3.1)-(3.1) is singular. The string coupling $`g_s\mathrm{exp}(Q\varphi /2)`$ diverges as $`\varphi \mathrm{}`$ (i.e. as one approaches the fivebranes (3.1)). Therefore, the weakly coupled ten dimensional description is only useful for studying those aspects of LST that can be analyzed at large positive $`\varphi `$. This includes identifying a large set of observables (such as the aforementioned non-normalizable vertex operators), and their transformation properties under the symmetries of the theory (e.g. $`6d`$ super-Poincare symmetry). Normalizable states are difficult to analyze, since their wave-functions tend to be supported in the strongly coupled region $`\varphi \mathrm{}`$. Equivalently, correlation functions of non-normalizable operators in linear dilaton vacua are typically not computable without specifying the cutoff at large negative $`\varphi `$.
If one is interested in studying the theory on $`k>1`$ coincident fivebranes, one must face this strong coupling problem. However, to make contact with section 2 we are in fact mainly interested in the case where the fivebranes are separated. In that case one might hope that the strong coupling region will be absent, e.g. because the CHS solution with the throat only makes sense when there are at least two coincident fivebranes. Indeed, one can show that this is the case.
If, for example, we distribute the $`k`$ fivebranes at equal distances around a circle in the $`(x^6,x^7)`$ plane, which breaks the global symmetry
$$SO(5,1)\times SO(4)SO(5,1)\times SO(2)\times Z_k$$
the background (3.1) changes as follows. Decompose $`SU(2)U(1)\times SU(2)/U(1)`$. The $`U(1)`$ is the CSA generator $`J_{\mathrm{total}}^3`$, which can be bosonized as
$$J_{\mathrm{total}}^3=2i\sqrt{\frac{k}{2}}Y$$
where $`Y`$ is a scalar field normalized such that $`Y(z)Y(w)=\frac{1}{4}\mathrm{log}(zw)`$, a normalization that will be convenient later; it differs by a factor of four from that of \[12,,13\]. In the original CHS solution, $`(\varphi ,Y)`$ describe an infinitely long cylinder with a string coupling that varies along the cylinder, diverging at one end and going to zero at the other. Separating the fivebranes replaces it \[14,,15,,12,,13\] with a semi-infinite cigar $`SL(2)/U(1)`$ in which the strong coupling region is absent, or equivalently with $`N=2`$ Liouville, in which the strong coupling region is suppressed by a superpotential which goes like $`\mathrm{exp}[(\widehat{\varphi }+i\widehat{Y})/Q]`$, where $`(\widehat{\varphi },\widehat{Y})`$ are the superfields whose bosonic components are $`(\varphi ,Y)`$. For reasons explained in it is more convenient to use the coset (cigar) description in this case. The full background replacing (3.1) is
$$\mathrm{IR}^{5,1}\times \frac{SL(2,\mathrm{IR})_k}{U(1)}\times \frac{SU(2)_k}{U(1)}$$
The global $`SO(2)`$ symmetry in (3.1) corresponds to translations in $`Y`$ (i.e. rotations of the cigar around its axis). The $`Z_k`$ charge in (3.1) is the winding number around the cigar, a $`U(1)`$ which is broken down to $`Z_k`$, since winding number can slip off the tip of the cigar. It should be noted that the product in (3.1) is not direct, since the GSO projection relates the different factors. In the GSO projected theory, the winding number around the cigar can be fractional ($`Z/k`$). The fractional part of the winding number is the $`Z_k`$ charge mentioned above.
Since the background (3.1) can be made arbitrarily weakly coupled<sup>4</sup> by decreasing the value of the string coupling at the tip of the cigar. This value is related to the radius of the circle on which the fivebranes lie ., one can use it to compute correlation functions in LST in a weak coupling regime in its moduli space of vacua . This is done by constructing BRST invariant observables on (3.1) and computing their worldsheet correlation functions, using the results of . For a detailed analysis we refer the reader to ; here we mention a few facts that will play a role below.
Consider the (NS,NS) sector of the theory<sup>5</sup> The other sectors can be reached by applying the spacetime supercharges .. The observables are primaries of the $`N=1`$ superconformal algebra (3.1) with scaling dimension $`(h,\overline{h})=(\frac{1}{2},\frac{1}{2})`$. The $`\mathrm{IR}^{5,1}`$ and $`SU(2)/U(1)`$ are well known SCFT’s. The former is a free field theory; the latter, an $`N=2`$ minimal model with $`c=3(6/k)`$. The $`SL(2)/U(1)`$ SCFT is $`N=2`$ superconformal as well; it has central charge $`c=3+(6/k)`$. The $`N=2`$ primaries $`V_{j;m,\overline{m}}`$ have scaling dimensions
$$(h,\overline{h})=\frac{1}{k}(m^2j(j+1),\overline{m}^2j(j+1))$$
and
$$(m,\overline{m})=\frac{1}{2}(p+wk,pwk)$$
where $`p,wZ`$ are momentum and winding around the cigar, respectively. As mentioned before, in the GSO projected theory (3.1) one finds
$$w\frac{1}{k}Z$$
while the momentum $`p`$ is still integer.
Unitarity and non-normalizability limit the range of $`j`$ to
$$j\mathrm{IR},\frac{1}{2}<j<\frac{k1}{2}$$
There are also delta-function normalizable operators with $`j=\frac{1}{2}+is`$, $`s\mathrm{IR}`$. These do not give rise to off-shell observables in the theory, but rather should be thought of as producing a continuum of states above a gap in LST<sup>6</sup> The relation between states and operators in the cigar CFT is subtle and very similar to that in Liouville theory, described in .. All other observables (in addition to $`V_{j;m,\overline{m}}`$) can be obtained by acting with $`N=2`$ superconformal generators on these primaries.
By analyzing correlation functions of $`V_{j;m,\overline{m}}`$, one finds that poles in correlators correspond to discrete representations<sup>7</sup> In some cases one also finds poles corresponding to $`m=j`$ . of $`SL(2)`$
$$|m|=j+n,|\overline{m}|=j+\overline{n};n,\overline{n}=1,2,3,\mathrm{}$$
This leads to a discrete spectrum of states in LST, which exhibits Hagedorn growth at high energy.
4. Some properties of $`D`$-branes
In this section we review some properties of $`D`$-branes in flat space and on $`S^3`$, in preparation for our discussion of $`D`$-branes in the CHS geometry (3.1), and its regularized version (3.1).
4.1. 4-6 strings in flat space
Later, we will analyze $`46`$ strings connecting $`D4`$-branes and $`D6^{}`$-branes, both of which end on a stack of $`NS5`$-branes (fig. 4a). We start by reviewing the simpler case of intersecting infinite $`D4`$ and $`D6^{}`$-branes in flat space \[20,,21\].
Consider an open string, one of whose ends is on a $`D4`$-brane. The other end may be either on the same brane or on another brane. We want to study the emission of a $`46^{}`$ string from the $`D4`$ boundary of this string. After the emission, this boundary of the emitting string will lie on a $`D6^{}`$-brane (fig. 6).
Fig. 6
The worldsheet of an open string is a strip, or equivalently the upper half plane. On the upper half plane, worldsheet time evolves radially. Equal time surfaces are half circles around the origin, with early times corresponding to small circles. The boundary of the upper half plane (the real line $`\mathrm{Im}z=0`$) corresponds to the ends of the string and is divided into two parts: $`z>0`$ or $`\sigma =0`$ in fig. 6, which lies on the $`D4`$-brane, and $`z<0`$ or $`\sigma =\pi `$.
The emission of an open $`46^{}`$ string from the boundary $`\sigma =0`$ is described by an insertion of a vertex operator $`V`$ at some point $`z>0`$ on the boundary. The boundary now splits into three components: (i) $`z^{}<0`$, the spectator boundary at $`\sigma =\pi `$, (ii) $`0<z^{}<z`$, which lies on the fourbrane, and (iii) $`z^{}>z`$, which is on the sixbrane (see fig. 7). The boundary conditions for, say, the coordinate $`x^4`$ are Dirichlet, $`x^4(z^{})+\overline{}x^4(z^{})=0`$, in region (ii), and Neumann, $`x^4(z^{})\overline{}x^4(z^{})=0`$, in region (iii).
Fig. 7
The vertex operator $`V(z)`$ changes the boundary conditions of $`x^4`$ from Neumann to Dirichlet. If one inserts the holomorphic current $`x^4`$ into the worldsheet, and moves it along a small upper semicircle around $`z`$ (see fig. 7), from $`z^{}<z`$ to $`z^{}>z`$ (the operator $`\overline{}x^4`$ is necessarily transported along the mirror image lower semicircle between these two points), the relative sign of the two operators has to flip. This means that the operator product $`V(z)x^4(z^{})`$ has to have a square root branch cut in $`zz^{}`$. Operators with such a cut are twist operators, familiar from orbifold CFT (see e.g. ). The lowest dimension operator of this type, $`\sigma `$, has dimension $`1/16`$.
The same arguments apply to all the coordinates for which the $`46^{}`$ string has Dirichlet-Neumann boundary conditions; hence, the operator $`V`$ contains a twist field for $`(x^4,x^5,x^6,x^7)`$, $`\sigma _{4567}`$, with total dimension $`1/4`$. We will be mainly interested in emission of spacetime bosons, in which case $`V`$ belongs to the $`NS`$ sector of the worldsheet CFT. Therefore, it has to be local with respect to $`G=i\psi _ax^a`$, the worldsheet superconformal generator. Since $`G`$ has a square root branch cut with respect to $`x^{4,5,6,7}`$, it must also have a cut with respect to $`\psi _{4,5,6,7}`$. This implies that $`V`$ has to contain the spin field $`S_{4567}`$ for these worldsheet fermions. This operator also has dimension $`1/4`$.
The directions $`x^{0,1,2,3}`$ are common to the $`D4`$ and $`D6^{}`$-branes, and are treated as in standard open string theory with Neumann boundary conditions. The remaining coordinates $`(x^8,x^9)`$ have Dirichlet-Dirichlet boundary conditions. If the $`D4`$-brane is, say, at $`x^8=0`$, while the $`D6`$-brane is at $`x^8=a`$ (and both at $`x^9=0`$), then similarly to the discussion above, a bulk operator of the form $`\mathrm{exp}(ikx^8)`$ when moved around the above mentioned semi-circle from a point on the $`D4`$ to a point on the $`D6^{}`$, has to pick a factor of $`\mathrm{exp}(ika/2)`$, together with another such contribution from the mirror path of the lower semicircle. This means that the operator product $`V(z)\mathrm{exp}(ikx^8(z^{}))`$ has a branch cut of the form $`(zz^{})^{ka/2\pi }`$ and the product with the corresponding right-moving operator has a cut with the opposite sign. The appropriate boundary operator is $`\mathrm{exp}(i(a/\pi )(x_L^8(z)x_R^8(z)))`$. This operator generates “winding number” $`a/\pi `$ along $`x^8`$ – in agreement with the geometrical picture of a $`46^{}`$ string stretched a distance $`a`$ along $`x^8`$. Its dimension is $`a^2/2\pi ^2`$.
Collecting all the factors, we get the following vertex operator describing the (bosonic) ground state of a $`46^{}`$ string stretched between a $`D4`$-brane at $`(x^8,x^9)=(0,0)`$ and a $`D6^{}`$-brane at $`(x^8,x^9)=(a,b)`$,
$$V=e^\phi \sigma _{4567}S_{4567}e^{\frac{i}{\pi }(a(x_L^8x_R^8)+b(x_L^9x_R^9))}e^{ik_\mu x^\mu }$$
where $`k_\mu `$ ($`\mu =0,1,2,3`$) is the $`4d`$ spacetime momentum. $`\phi `$ is the bosonized superconformal ghost. The vertex operator (4.1) is written in the $`1`$ picture; one can check that the coefficient of $`e^\phi `$ in (4.1) is an $`N=1`$ superconformal primary, which is a necessary condition for its BRST invariance. The requirement that it has worldsheet dimension $`1/2`$, which is also necessary for BRST invariance, implies that the mass squared of the ground state of the $`46^{}`$ string is $`k_\mu ^2=\frac{1}{\pi ^2}(a^2+b^2)`$, as one would expect (we work in a convention $`\alpha ^{}=1/2`$, in which the scalars $`x`$ are canonically normalized on the boundary and the tension of the fundamental string is $`T=1/\pi `$).
In particular, when the $`D6^{}`$-brane intersects the $`D4`$-brane, i.e. when $`a=b=0`$, this mass vanishes. The vertex operator (4.1) describes a particle which transforms as a scalar under $`3+1`$ dimensional Lorentz rotations. The spin field $`S_{4567}`$ has $`4`$ components, half of which are projected out by the GSO projection, so (4.1) actually describes two real scalar particles. For the applications described in section 2 it is sometimes useful to consider not one but a stack of $`N`$ $`D4`$-branes. In that case, the two scalars (4.1) transform in the fundamental representation ($`𝐍`$) of the $`U(N)`$ gauge symmetry on the fourbranes.
In addition to the $`46^{}`$ strings described above there are also $`6^{}4`$ strings which have similar properties but transform in the $`\overline{𝐍}`$ of $`U(N)`$. Altogether we have two complex scalars, $`Q`$ in the $`𝐍`$ and $`\stackrel{~}{Q}`$ in the $`\overline{𝐍}`$ of $`U(N)`$. The system of $`D4`$ and $`D6^{}`$-branes preserves eight supercharges ($`N=2`$ SUSY in the four dimensions $`(0123)`$) and acting with the supercharges on (4.1) completes a hypermultiplet transforming in the fundamental representation of $`U(N)`$.
4.2. $`D`$-branes on the $`SU(2)`$ group manifold
We next turn to some facts regarding $`D`$-branes on a group manifold $`G`$, focusing on the case $`G=SU(2)`$ ($`D`$-branes on group manifolds have been studied, for instance, in \[23,,24,,25,,26,,27,,28,,29,,30,,31,,32,,33,,34,,35,,36,,37,,38\]). In the absence of $`D`$-branes, the WZW model has an affine $`G_L\times G_R`$ symmetry. If $`g(z,\overline{z})`$ is a map from the worldsheet to the group $`G`$, the symmetry acts on it as:
$$gh_L(z)gh_R(\overline{z})$$
If the worldsheet has a boundary, there is a relation between left-moving and right-moving modes, and the $`G_L\times G_R`$ symmetry is broken. One can still preserve some diagonal symmetry $`G`$, say the symmetry
$$ghgh^1$$
corresponding to $`h_L=h_R^1=h`$ in (4.1). The presence of this symmetry constrains the boundary conditions that can be placed on $`g`$. Allowing $`g(\mathrm{boundary})=f`$ for some $`fG`$ we must also allow $`g(\mathrm{boundary})=hfh^1`$ for every $`hG`$. This means that $`g`$ on the boundary takes value in the conjugacy class containing $`f`$ . For $`G=SU(2)`$, conjugacy classes are parametrized by a single angle $`\theta `$, $`0\theta \pi `$, corresponding to the choice $`f=\mathrm{exp}(i\theta \sigma _3)`$. Thinking of $`SU(2)`$ as the group of three dimensional rotations, the conjugacy class $`C_\theta `$ is the set of all rotations of angle $`2\theta `$ about any axis. A boundary condition which preserves (4.1) is then
$$g(\mathrm{boundary})C_\theta $$
Since $`h_L`$ is generated by the currents $`J^a`$ while $`h_R`$ is generated by $`\overline{J}^a`$, the invariance of the boundary condition (4.1) under (4.1) implies that the currents satisfy :
$$J^a=\overline{J}^a;a=1,2,3$$
on the boundary. Not any value of $`\theta `$ in (4.1) gives rise to a consistent model \[24,,27,,29\]. Recall that for a general group $`G`$, the level $`k`$ WZW action has the form
$$S=_\mathrm{\Sigma }d^2zL^{sm}+_BL^{WZ}$$
where $`L^{sm}=\frac{k}{4\pi }Tr(g^1\overline{}g)`$ is the sigma model part of the action and $`L^{WZ}=\frac{k}{4\pi }\omega ^{(3)}`$ is the Wess-Zumino term. Here $`\mathrm{\Sigma }`$ is the worldsheet Riemann surface, $`\omega ^{(3)}=\frac{1}{3}Tr((g^1dg)^3)`$ is a closed three-form, and $`B`$ is a three dimensional manifold whose boundary is the worldsheet. When the worldsheet $`\mathrm{\Sigma }`$ has itself boundaries, it cannot be the boundary of a three dimensional manifold, since a boundary cannot have boundary. To define the WZ term in (4.1) for this case, one should fill the holes in the worldsheet by adding discs, and extend the mapping from the worldsheet into the group manifold to these discs. One further demands that the whole disc $`D`$ is mapped into a region (which we will also refer to as $`D`$) inside the conjugacy class in which the corresponding boundary lies. $`B`$ will then be defined as a three-manifold bounded by the union $`\mathrm{\Sigma }D`$, which now has no boundaries.
The resulting action should preserve the symmetry (4.1) in the bulk of the worldsheet which tends to (4.1) on the boundary. Take $`C`$, the conjugacy class containing $`D`$, to be the class of a fixed group element $`f`$, i.e.
$$C=\{hfh^1|hG\}$$
Let $`\delta _Lg`$ be an infinitesimal variation of $`g(z,\overline{z})`$ in the bulk such that $`\overline{}(\delta _Lgg^1)=0`$ and $`\delta _Rg`$ a variation for which $`(g^1\delta _Rg)=0`$. On the boundary of $`\mathrm{\Sigma }`$, where $`gC`$ is parametrized as $`g=hfh^1`$, $`\delta _Lgg^1=g^1\delta _Rg=\delta hh^1`$. By the above symmetry the variation $`(\delta _L+\delta _R)S`$ should vanish. In the interior of the disc we should allow an arbitrary variation $`\delta h`$ since the location of the auxiliary disc inside $`C`$ has no physical significance and cannot influence the action. Using the identity
$$\delta \omega ^{(3)}=dTr((g^1\delta g)(g^1dg)^2)$$
one gets
$$(\delta _L+\delta _R)S=\frac{k}{4\pi }\left[_DTr((g^1\delta g)(g^1dg)^2)+_\mathrm{\Sigma }A\right]$$
where $`A`$ is the one-form
$$A=Tr[(h^1\delta h)(f^1h^1dhffh^1dhf^1)]$$
The first term in (4.1) comes from the change in the region of integration of $`L^{WZ}`$ resulting from the variation of $`D`$. The second term is a boundary correction to the symmetry (4.1). Clearly the right hand side of (4.1) is not zero for a non-trivial mapping of $`\mathrm{\Sigma }`$ and of $`D`$ into $`C`$.
To fix the symmetry we have to modify the action by adding to it an integral over the disc $`D`$ of some two form $`\omega ^{(2)}`$, defined on $`C`$, such that its variation will cancel (4.1). The proper action has now the form
$$S=_\mathrm{\Sigma }d^2zL^{sm}+_BL^{WZ}+\frac{k}{4\pi }_D\omega ^{(2)}$$
The variation of this action is
$$(\delta _L+\delta _R)S=\frac{k}{4\pi }\left[_DTr((g^1\delta g)(g^1dg)^2)_DA+_D\delta \omega ^{(2)}\right]$$
where we used the fact that $`\mathrm{\Sigma }`$ and $`D`$ are identical curves with opposite orientations. The vanishing of (4.1) fixes $`\delta \omega ^{(2)}`$ as
$$\delta \omega ^{(2)}=dATr((g^1\delta g)(g^1dg)^2)$$
Expressing the r.h.s. of (4.1) in terms of the group element $`h`$ parametrizing $`C`$ as in (4.1) one gets $`\delta \omega ^{(2)}`$ explicitly on $`C`$. Its solution is \[27,,29\]
$$\omega ^{(2)}=Tr[(h^1dh)f^1(h^1dh)f]$$
Notice that (as implied by (4.1)) on the conjugacy class $`C`$, $`d\omega ^{(2)}=\omega ^{(3)}`$.
The modified action (4.1) is independent, by construction, of continuous deformations of $`D`$ inside $`C`$. However, in general, the second homotopy of a conjugacy class is non-trivial. If we compare then the value of the action for $`D`$ and $`D^{}`$, two different choices of embedding the disc in $`C`$ with the same boundary, $`D^{}`$ may not be a continuous deformation of $`D`$ in $`C`$. In that case the above analysis does not imply that the two ways to evaluate the action (4.1) agree. Since there is no natural way to choose between the two embeddings, (4.1) is not yet a well defined action. In particular, for $`G=SU(2)`$ the conjugacy classes $`C`$ have the topology of $`S^2`$, the two-sphere generated by all possible axes of rotation by a fixed angle in three dimensions. One may then choose $`D`$ and $`D^{}`$ such that their union covers the whole of $`S^2`$. In that case the difference between the action $`S_D`$, the value of (4.1) with embedding $`D`$, and $`S_D^{}`$ with embedding $`D^{}`$ is
$$\mathrm{\Delta }S=\frac{k}{4\pi }\left[_B\omega ^{(3)}+_C\omega ^{(2)}\right]$$
where $`B`$ is the three-volume in $`SU(2)`$ bounded by the two-sphere $`C`$ (fig. 8).
Fig. 8
For the case of $`SU(2)`$, which has the topology of $`S^3`$, the form $`\omega ^{(3)}`$ is $`\omega ^{(3)}=4\mathrm{\Omega }^{(3)}`$ where $`\mathrm{\Omega }^{(3)}`$ is the volume form on the unit three-sphere. For $`C`$ in (4.1) corresponding to $`f=\mathrm{exp}(i\theta \sigma _3)`$, the first term in (4.1) is
$$_B\omega ^{(3)}=8\pi (\theta \frac{1}{2}\mathrm{sin}(2\theta ))$$
As to the two-form $`\omega ^{(2)}`$, from eq. (4.1) we see that $`\omega ^{(2)}[h]=\omega ^{(2)}[qh]`$ for any fixed element $`qSU(2)`$. This implies that this form is proportional to $`\mathrm{\Omega }^{(2)}`$, the volume form of the unit two-sphere. The expression (4.1) for $`h`$ near the identity gives then for the conjugacy class $`C_\theta `$,
$$\omega ^{(2)}=\mathrm{sin}(2\theta )\mathrm{\Omega }^{(2)}$$
This gives for the change in the action for two topologically different embeddings in (4.1)
$$\mathrm{\Delta }S=2k\theta $$
Although this is non-zero, the quantum theory is still well defined if $`\mathrm{\Delta }S`$ is an integral multiple of $`2\pi `$. We find then that the possible conjugacy classes on which a boundary state can live are quantized, the corresponding $`\theta `$ must satisfy
$$\theta =2\pi \frac{j}{k}$$
with $`j`$ integer or half integer satisfying $`0j\frac{k}{2}`$.
Thus, there are $`k+1`$ different boundary conditions preserving the symmetry (4.1) which one can impose on the $`SU(2)_k`$ group manifold. This matches nicely with the algebraic analysis of Cardy , who found that in a general rational CFT on a worldsheet with boundary, to each primary field of the chiral algebra there corresponds a boundary state preserving the diagonal chiral algebra, the analog of (4.1). In the case of the $`SU(2)`$ WZW model, the chiral algebra is affine $`SU(2)_k`$, and there are $`k+1`$ different primary fields corresponding to representations of (integer or half integer) spin $`0j\frac{k}{2}`$. It is natural, \[27,,29\], to identify each of the $`k+1`$ different geometric boundary conditions corresponding to $`j`$ in (4.1) with the algebraic boundary state corresponding to the same $`j`$.
Since all of these boundary states preserve (4.1), the open strings stretching between them should have well defined transformation properties under the diagonal chiral algebra. According to , the open strings stretched between a boundary state corresponding to a primary field $`j`$ and another boundary state corresponding to primary field $`j^{}`$, belong to the representations of the chiral algebra which appear in the fusion of $`j`$ and $`j^{}`$.
We have chosen the boundary conditions (4.1) such that the particular diagonal subgroup $`G`$ of $`G\times G`$ defined in eq. (4.1) will survive them. This is of course not a unique choice. One can act on these boundary conditions with any element of the $`G\times G`$ symmetry group to get equivalent boundary conditions which preserve a different diagonal subgroup. Thus we can multiply the conjugacy class $`C_\theta `$ in eq. (4.1) from the right (which is the same as multiplying from the left) by any group element $`f`$ to get modified boundary conditions
$$g(\mathrm{boundary})C_\theta f$$
These boundary conditions also preserve a diagonal subgroup, since the set $`C_\theta f`$ satisfies
$$C_\theta f=h(C_\theta f)f^1h^1f$$
for any $`hG`$. Therefore, the boundary conditions (4.1) preserve the diagonal subgroup of $`G_L\times G_R`$ defined by $`h_R=f^1h_L^1f`$. In terms of the infinitesimal generators of $`G_L\times G_R`$, i.e. the left and right handed currents, the invariance of (4.1) under this subgroup implies for the corresponding boundary state the condition
$$J^a=(Ad_{f^1}\overline{J})^a$$
which modifies (4.1) by conjugating the right handed currents by $`f^1`$.
5. $`D`$-branes in the near-horizon geometry of $`NS5`$-branes
After assembling the necessary tools, we are now ready to study the physics of the configurations of fig. 4.
5.1. $`D4`$ and $`D6^{}`$-branes ending on $`NS5`$-branes
We start with the configuration of fig. 4a. A stack of $`N_c`$ $`D4`$-branes ends from the left, i.e. from negative $`x^6`$, on $`k`$ coincident $`NS5`$-branes. $`N_F`$ $`D6^{}`$-branes end on the $`NS5`$-branes from above (positive $`x^7`$). From the point of view of the geometry (3.1) (fig. 5), the $`D`$-branes extend into the CHS throat, as indicated in fig. 9a.
The $`D4`$-branes intersect the three-sphere at the point $`x^7=x^8=x^9=0`$; the $`D6^{}`$-branes at $`x^6=x^8=x^9=0`$ (see fig. 9b). Thus, (3.1), they correspond to the boundary states $`g|_{\mathrm{boundary}}=1`$ and $`i\sigma _3`$, respectively. The $`D4`$-branes are described by the boundary state with $`\theta =0`$ and $`f=1`$ (see (4.1), (4.1)), while the $`D6^{}`$-branes correspond to a transformed state (4.1), with $`\theta =0`$ and $`f=\mathrm{exp}(i\pi \sigma _3/2)=i\sigma _3`$. The $`SU(2)`$ currents $`J^a`$ satisfy the boundary conditions (4.1) and (4.1) for strings ending on the $`D4`$ and $`D6^{}`$-branes, respectively. In order to preserve worldsheet supersymmetry one has to impose analogous boundary conditions on the fermions. For example, for a string ending on the $`D4`$-branes one has
$$\chi ^a=\overline{\chi }^a$$
while for a boundary on a $`D6^{}`$-brane the $`\chi ^a`$ satisfy an analog of (4.1). Since, as is clear from fig. 9, both the fourbranes and the sixbranes extend into the throat, the boundary conditions on $`\varphi `$ are Neumann.
Fig. 9
We would like to construct the vertex operator for emitting the lowest lying $`46^{}`$ string in the geometry of fig. 9, i.e. generalize (4.1) to the fivebrane near-horizon geometry. Some parts of the discussion leading to (4.1) are unchanged. In particular, the geometry is the same as there in $`(x^0,x^1,x^2,x^3,x^4,x^5)`$. The presence of $`\varphi `$ allows also a contribution $`\mathrm{exp}(\beta \varphi )`$ to the vertex operator. Thus, the analog of (4.1) for this case has the form
$$V=e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }e^{\beta \varphi }V_2$$
where $`V_2`$ is the contribution of the $`SU(2)`$ group manifold to the vertex operator, to which we turn next.
The $`46^{}`$ vertex operator $`V_2`$ changes the worldsheet boundary conditions from $`g=1`$ to $`g=f=\mathrm{exp}(i\alpha \sigma _3/2)`$. The $`D6^{}`$-brane corresponds to $`\alpha =\pi `$, but it is instructive to discuss the general case, in which the angle between the $`D4`$ and $`D6^{}`$-branes is $`\alpha /2`$ (see fig. 10).
Fig. 10
As discussed above, $`g=1`$ corresponds to the boundary condition (4.1), (5.1) for the $`SU(2)`$ currents and fermions, while $`g=f`$ gives rise to (4.1)
$$\begin{array}{cc}& J^3=\overline{J}^3\hfill \\ & \chi ^3=\overline{\chi }^3\hfill \\ & J^\pm =\mathrm{exp}(i\alpha )\overline{J}^\pm \hfill \\ & \chi ^\pm =\mathrm{exp}(i\alpha )\overline{\chi }^\pm \hfill \end{array}$$
Note that the symmetry generated by $`J^3+\overline{J}^3`$ is preserved by both (4.1) and (5.1). As discussed after eq. (3.1), it corresponds to the rotation symmetry $`SO(2)_{89}`$ which is unbroken by the brane configuration. As in section 4.1, we conclude that $`V_2`$ must include a twist field $`\sigma _\alpha `$ with the following locality properties w.r.t. the left-moving currents $`J^a`$ and fermions $`\chi ^a`$:
$$\begin{array}{cc}& \sigma _\alpha (z)J^3(z^{})(zz^{})^{m_1}O_1(z^{})\hfill \\ & \sigma _\alpha (z)\chi ^3(z^{})(zz^{})^{m_2}O_2(z^{})\hfill \\ & \sigma _\alpha (z)J^\pm (z^{})(zz^{})^{\frac{\alpha }{2\pi }+n_1}O_3(z^{})\hfill \\ & \sigma _\alpha (z)\chi ^\pm (z^{})(zz^{})^{\frac{\alpha }{2\pi }+n_2}O_4(z^{})\hfill \end{array}$$
with $`m_i,n_iZ`$; $`O_I`$ on the r.h.s. are operators whose precise form will not be specified here. A similar twist holds for the right-movers. Equation (5.1) implies that $`V_2`$ belongs to a representation of a twisted affine $`SU(2)`$, in which the currents $`J^\pm `$ have fractional modes. This twisted algebra reads
$$\begin{array}{cc}& [J_n^3,J_m^3]=n\frac{k_B}{2}\delta _{n,m}\hfill \\ & [J_n^3,J_{m\pm \frac{\alpha }{2\pi }}^\pm ]=\pm J_{n+m\pm \frac{\alpha }{2\pi }}^\pm \hfill \\ & [J_{n+\frac{\alpha }{2\pi }}^+,J_{m\frac{\alpha }{2\pi }}^{}]=(n+\frac{\alpha }{2\pi })k_B\delta _{n,m}+2J_{n+m}^3\hfill \end{array}$$
where $`k_B=k2`$ is the level of the bosonic $`SU(2)`$ algebra (as in section 3). The algebra (5.1) can be mapped into the standard (untwisted) affine Lie algebra by using spectral flow. If $`J^a`$ satisfy the twisted algebra (5.1), the generators $`\stackrel{~}{J}`$ defined by
$$\begin{array}{cc}& \stackrel{~}{J}_n^3=J_n^3+\frac{k_B}{2}\frac{\alpha }{2\pi }\delta _{n,o}\hfill \\ & \stackrel{~}{J}_n^\pm =J_{n\pm \frac{\alpha }{2\pi }}^\pm \hfill \end{array}$$
satisfy the ordinary untwisted algebra. Thus one can use standard facts about the representations of untwisted affine Lie algebra to study the twisted one.
The modes of the energy-momentum tensor ($`\stackrel{~}{L}_n`$) constructed from $`\stackrel{~}{J}`$ are related to those of the original energy-momentum tensor ($`L_n`$) via
$$\stackrel{~}{L}_n=L_n+\frac{\alpha }{2\pi }J_n^3+\frac{k_B}{4}\frac{\alpha ^2}{4\pi ^2}\delta _{n,0}$$
To understand the properties of the operator $`\sigma _\alpha `$, consider first the case $`\alpha =0`$. $`\sigma _0`$ describes an open string connecting two boundary states, both corresponding to $`j=0`$ in eq. (4.1). As discussed in section 4.2, this open string should transform in a representation of the diagonal $`SU(2)`$ contained in the fusion of two $`j=0`$ representations, which consists of only the spin $`0`$ representation. If $`h`$ is the energy ($`L_0`$) and $`m`$ the $`J^3`$ charge, then the lowest energy state in this sector has $`h=m=0`$ and all the excited states satisfy
$$h|m|$$
Turning on $`\alpha `$ continuously, the open string created by $`\sigma _\alpha `$ remains in the spin 0 representation, now of the affine algebra generated by the currents $`\stackrel{~}{J}`$ (5.1). Hence the corresponding states satisfy
$$\stackrel{~}{h}|\stackrel{~}{m}|$$
or, using (5.1),
$$h(\pm 1\frac{\alpha }{2\pi })\stackrel{~}{m}+\frac{k_B}{4}\frac{\alpha ^2}{4\pi ^2}$$
where $`\stackrel{~}{m}Z`$. For $`0\alpha <2\pi `$, the lowest energy state corresponds to $`\stackrel{~}{m}=0`$. The dimension and charge of the corresponding operator are
$$\begin{array}{cc}\hfill h=& \frac{k_B}{4}\frac{\alpha ^2}{4\pi ^2}\hfill \\ \hfill m=& \frac{k_B}{2}\frac{\alpha }{2\pi }\hfill \end{array}$$
To construct this operator it is convenient to decompose the $`SU(2)_{k_B}`$ bosonic CFT under $`U(1)\times SU(2)/U(1)`$, where $`U(1)`$ represents the subgroup generated by $`J^3`$. One can bosonize<sup>8</sup> As in section 4.1, here and below scalar fields are canonically normalized on the boundary; e.g. $`u(z)u(w)=\mathrm{log}|zw|`$ for $`z,wR`$. $`J^3`$ as in (3.1), $`J^3=2i\sqrt{k_B/2}u`$; $`J^\pm `$ are represented by $`\mathrm{exp}(\pm 2i\sqrt{2/k_B}u)`$ multiplied by operators from the $`SU(2)/U(1)`$ coset. The operators
$$\sigma _\alpha ^B=\mathrm{exp}\left[i\sqrt{\frac{k_B}{2}}(\frac{\alpha }{2\pi }+n)u\right]$$
with $`nZ`$, have the right locality properties (5.1) with respect to the currents. The dimension and charge (5.1) correspond to those of the operator (5.1) with $`n=0`$. Note that for $`\pi \alpha <2\pi `$, setting $`n=1`$ in (5.1) would give a lower dimension than that of $`n=0`$, however, the charge and dimension do not satisfy the inequality (5.1) in this case. Hence, these operators are not in the spectrum.
The operator $`\sigma _\alpha `$ must also flip the boundary conditions of the fermions $`\chi ^a`$ from (5.1) to (5.1). This means that it must include a spin field for $`\chi ^\pm `$. Bosonizing the two fermions,
$$iH=\chi ^+\chi ^{}$$
one finds that the operators with the right locality properties w.r.t. $`\chi ^\pm `$ are
$$\mathrm{\Sigma }_n=e^{i(n\frac{\alpha }{2\pi })H}$$
where $`nZ`$. For $`0\alpha <\pi `$ the operator of lowest dimension out of these is $`\mathrm{\Sigma }_0`$ whose scaling dimension is $`h=\alpha ^2/8\pi ^2`$ and $`U(1)`$ charge $`\alpha /2\pi `$. Thus, in this range of $`\alpha `$, the lowest lying state in the sector twisted by $`\alpha `$ is
$$\sigma _\alpha =e^{i\frac{\alpha }{2\pi }H}e^{i\frac{\alpha }{2\pi }\sqrt{\frac{k_B}{2}}u}$$
Notice that $`\sigma _\alpha `$ is a twist field for $`J_{\mathrm{total}}`$, i.e. one can rewrite it as
$$\sigma _\alpha =e^{i\sqrt{\frac{k}{2}}\frac{\alpha }{2\pi }Y}$$
where $`Y`$ is defined in (3.1) and is related to $`u`$ and $`H`$ via
$$\sqrt{\frac{k}{2}}Y=\sqrt{\frac{k_B}{2}}u+H$$
The operators (5.1) with $`n0`$ can be thought of as “excited” twist fields. For $`\pi <\alpha 2\pi `$ the operator $`\mathrm{\Sigma }_1`$ in (5.1) has lower dimension then that of $`\mathrm{\Sigma }_0`$. Then the excited twist operator
$$\sigma _\alpha ^{}=e^{i(\frac{\alpha }{2\pi }1)H}e^{i\sqrt{\frac{k_B}{2}}\frac{\alpha }{2\pi }u}$$
has the smallest dimension in the $`\alpha `$ twisted sector<sup>9</sup> Applying the GSO projection to (5.1) will pick up different spin fields $`S_{45}`$ for the different twist fields (5.1) and (5.1).. For $`\alpha =\pi `$, the supersymmetric case, the two operators in (5.1) and in (5.1) are degenerate on the worldsheet. Nevertheless, we show in Appendix B that the operator (5.1) creates from the vacuum excited $`46^{}`$ strings, even for $`\alpha =\pi `$.
To summarize, the lowest lying open string connecting the $`D4`$ and $`D6^{}`$-branes in the configuration of fig. 10 (which we will refer to as a $`46_+^{}`$ string, since it connects the fourbrane to a half-infinite $`D6^{}`$-brane at $`x^7>0`$) is described by the vertex operator
$$V_{46^{}}^+=e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }e^{\beta \varphi }e^{i\frac{\alpha }{2\pi }\sqrt{\frac{k}{2}}Y}$$
The mass shell condition requires that
$$\frac{1}{2}k_\mu ^2\frac{1}{2}\beta (\beta +Q)+\frac{k}{4}\left(\frac{\alpha }{2\pi }\right)^2=\frac{1}{4}$$
We see that the $`46_+^{}`$ string vertex (5.1) describes an on-shell excitation in five non-compact dimensions, $`(x^0,x^1,x^2,x^3,\varphi )`$. This is clearly the same phenomenon as that discussed for closed strings in section 3. The extra non-compact dimension is the infinite throat of associated with the radial direction. Non-normalizable open string vertex operators such as (5.1) with $`\beta >Q/2`$ correspond to off-shell observables in the theory on the $`D`$-branes.
To find the spectrum of low lying on-shell states, which is relevant for analyzing the physics of the brane configurations described in section 2, one has to find the normalizable states corresponding to $`46_+^{}`$ strings. One way of doing that is to compute the correlation functions of the non-normalizable operators (5.1) and extract the spectrum by analyzing their analytic structure. As explained in section 3, to do this within a weakly coupled string theory one has to regularize the strong coupling region $`\varphi \mathrm{}`$. Following , we will do this by considering the geometry (3.1) corresponding to fivebranes separated along a circle in the $`(x^6,x^7)`$ plane<sup>10</sup> One might wonder whether the separation of the fivebranes introduces a finite mass shift for $`46^{}`$ strings, when the fourbrane and sixbrane end on different fivebranes. To see that this does not happen note that, as discussed in , the cigar geometry corresponds to the limit $`g_s,r_0/l_s0`$ with $`r_0/(l_sg_s)`$ fixed and large ($`r_0`$ is the typical separation between fivebranes). The mass shift for $`46^{}`$ strings associated with this separation is $`r_0/l_s^2`$; it goes to zero in string units in the above limit..
The effect of this regularization on the vertex operator (5.1) is to change the wavefunction on the infinite cylinder labeled by $`(\varphi ,Y)`$ to a wavefunction on the semi-infinite cigar,
$$e^{\beta \varphi }e^{i\frac{\alpha }{2\pi }\sqrt{\frac{k}{2}}Y}V_{jm}$$
where
$$\begin{array}{cc}\hfill j=& \beta \sqrt{\frac{k}{2}}\hfill \\ \hfill m=& \frac{k}{2}\frac{\alpha }{2\pi }\hfill \end{array}$$
$`j`$ can be thought of as momentum along the cigar, while $`m`$ is the momentum around the cigar (both the $`D4`$-branes and $`D6^{}`$-branes are wrapping the cylinder, and in particular the circle labeled by $`Y`$).
The vertex operator (5.1) becomes in the geometry (3.1)
$$V_{46^{}}^+=e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }V_{jm}$$
with (5.1)
$$\frac{k_\mu ^2}{2}+\frac{m^2j(j+1)}{k}=\frac{1}{4}$$
To compute the spectrum of low-lying normalizable excitations of the $`46_+^{}`$ string, one computes the two point function of the operators (5.1) on the disk. This amplitude exhibits first order poles in $`j`$. Using the mass-shell condition (5.1), these can be interpreted as poles in $`k_\mu ^2`$; they correspond to on-shell particles in four dimensions, created from the vacuum by the operator (5.1).
The calculation of this two point function is very similar to its closed string analog \[17,,13\], and is described in appendix A. The result is, as in the closed string case, a series of poles corresponding to the discrete representations of $`SL(2)`$, (3.1). The lowest lying state corresponds to $`n=1`$, i.e. $`j=|m|1`$. Plugging this together with (5.1) into (5.1) we find that the mass of the lowest lying normalizable state of the $`46_+^{}`$ string is
$$M^2(\alpha )=\frac{1}{2}(\frac{\alpha }{\pi }1)$$
In particular, we find that as expected, for $`\alpha =\pi `$ the lowest lying state is massless. The vertex operator which creates this massless particle from the vacuum is
$$V_{46^{}}^+(k_\mu ^2=0)=e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }V_{\frac{k}{4}1,\frac{k}{4}}$$
Note that the degeneracy is correct too. The spin field $`S_{45}`$ in (5.1) takes apriori two possible values, but the GSO projection picks one of them. Thus, (5.1) describes a single complex scalar field $`Q`$ which transforms in the $`(N_c,N_F)`$ of the $`U(N_c)\times U(N_F)`$ symmetry on the $`D`$-branes, as explained in section 2. Applying the spacetime supercharges gives rise to a chiral superfield <sup>11</sup> In particular, one can check that (5.1), (5.1) is annihilated by half of the spacetime supercharges. $`Q`$. The complex conjugate field $`Q^{}`$ arises from $`46_+^{}`$ strings with the opposite orientation (a $`6_+^{}4`$ string). The field $`\stackrel{~}{Q}`$ ($`\stackrel{~}{Q}^{}`$) transforming in the $`(\overline{N}_c,\overline{N}_F)`$ representation that is needed for anomaly cancellation of $`N=1`$ SQCD comes from $`6_{}^{}4`$ ($`46_{}^{}`$) strings connecting the fourbranes to sixbranes attached to the $`NS5`$-branes from below, as anticipated in . These are described by vertex operators of the form (5.1), (5.1), (5.1) with $`\alpha <0`$ and with the $`S_{45}^{}`$ spin field for the $`46_{}^{}`$ operator $`\stackrel{~}{Q}^{}`$ (replacing $`S_{45}`$ for the $`46_+^{}`$ operator $`Q`$).<sup>12</sup> The reason for the $`S_{45}^{}`$ is the following. The $`D6_{}^{}`$ and $`D6_+^{}`$-branes must have the same orientation. They can be obtained by starting with parallel $`D6^{}`$ and $`\overline{D}6^{}`$-branes (at $`\alpha =0`$). The $`D6_+^{}`$-brane is obtained by rotating the $`D6^{}`$-brane by an angle $`\frac{\alpha }{2}=\frac{\pi }{2}`$, while the $`D6_{}^{}`$-brane is obtained by rotating the $`\overline{D}6^{}`$-brane by an angle $`\frac{\alpha }{2}=\frac{\pi }{2}`$. The conjugation of $`S_{45}`$ in passing from the $`D6^{}`$ to the $`\overline{D}6^{}`$-brane is due to the different GSO projection on branes and anti-branes.
Some comments are in order at this point:
(1) The above analysis is valid as long as $`j=|m|1>\frac{1}{2}`$ satisfies the unitarity bound (3.1), i.e. for $`\alpha >\alpha _c=\frac{2\pi }{k}`$. For $`\alpha <\alpha _c`$, the lowest lying normalizable states belong to the continuum ($`j=\frac{1}{2}+is`$) and the physics is different.
(2) An interesting question is whether $`m`$, (5.1), is quantized on the cigar. Before separating the $`NS5`$-branes on a circle, it is clear that $`\alpha `$ in (5.1) (and hence $`m`$) is arbitrary – the angle between the $`D4`$ and $`D6^{}`$-branes is not constrained. After the separation the situation is different. The $`D`$-branes are wrapped around the cigar, and the momentum around the cigar $`m`$ appears to be quantized, $`mZ`$. For $`D`$-branes wrapped around a circle one can shift the momentum by an arbitrary fractional amount, by turning on a Wilson line of the gauge field on the branes (around the circle). However, on the cigar there are no non-contractible cycles, and hence one expects the quantization of $`m`$ to persist.
(3) The quantization of $`m`$ can be seen in other ways as well. One is to use a T-dual description, in which the cigar is replaced by $`N=2`$ Liouville \[13,,40,,41\], and the $`D`$-branes are at points on the dual $`S^1`$. The $`N=2`$ Liouville superpotential effectively pins down the (T-dual) field $`Y`$ (3.1) to points on the circle when one goes far down the $`N=2`$ Liouville throat to the strong coupling region. Since the $`D`$-branes extend into the throat and lie at points (independent of $`\varphi `$) on the circle, their positions on the circle must coincide with those determined by the superpotential. Another way to see the quantization is to note that quantization of $`\alpha `$ has a natural geometric interpretation when the branes are ending on fivebranes which lie on a circle. Since the $`D4`$ and $`D6^{}`$-branes point towards the center of the circle on which the fivebranes lie, the angle between the $`D`$-branes, $`\frac{\alpha }{2}`$, must be an integer multiple of $`\frac{2\pi }{k}`$ – the angular separation between the $`k`$ fivebranes. This leads to the same conclusion, (5.1), $`mZ`$. Thus, the $`D`$-branes seem to know that the smooth cigar is in fact associated with $`k`$ fivebranes on a circle.
(4) It is interesting to compare the symmetry properties of (5.1) to those expected of $`Q`$, $`\stackrel{~}{Q}`$. The vertex operator $`V_{46^{}}^+`$ is charged under the unbroken rotation symmetry $`SO(2)_{45}\times SO(2)_{89}`$ of the configuration of fig. 10. The $`SO(2)_{45}`$ charge is carried by the spin field $`S_{45}`$ in (5.1). In units in which the supercharges have $`SO(2)_{45}`$ charge $`\pm \frac{1}{2}`$, $`V_{46^{}}^+`$ has charge $`\frac{1}{2}`$. The $`SO(2)_{89}`$ charge of (5.1) is carried by the last factor, $`\mathrm{exp}(i\frac{\alpha }{2\pi }\sqrt{\frac{k}{2}}Y)`$. Normalizing it such that the supercharges again have charge $`\pm \frac{1}{2}`$, one finds that the charge of $`V_{46^{}}^+`$ is $`\frac{k\alpha }{4\pi }`$. Therefore, the $`SO(2)_{89}`$ charge of (5.1) is $`k/4`$. From the discussion above one learns that $`\stackrel{~}{Q}`$ carries the same charge as $`Q`$. These assignments are similar but not identical to those postulated in brane theory in the past. In our notation, the $`SO(2)_{45}\times SO(2)_{89}`$ charge of $`Q,\stackrel{~}{Q}`$ was postulated to be $`(\frac{1}{2},0)`$ (see e.g. , discussion between eqs. (182) and (183)). In the present construction, we find charge $`(\frac{1}{2},\frac{k\alpha }{4\pi })`$ for $`Q`$ and $`\stackrel{~}{Q}`$. The discrepancy in the $`SO(2)_{89}`$ charges might be related to that found in ; a better understanding of its origin is left for future work.
For $`\alpha <\pi `$ the lowest lying state (5.1) is tachyonic; the stable vacuum is obtained by its condensation. This process has a very natural interpretation from the point of view of brane theory, which also makes it clear what is the endpoint of the condensation. For $`\alpha \pm \pi `$, the configuration of fig. 10 is not supersymmetric, hence stability needs to be checked. For $`\alpha <\pi `$ the configuration of fig. 10 can reduce its energy by having the $`D4`$-brane slide away from the $`NS5`$-brane so that it ends on the $`D6^{}`$-brane instead (fig. 11). The resulting vacuum is stable. For $`\alpha >\pi `$, the configuration of figs. 10, 11 is stable under small deformations. Indeed, the lowest lying open string state is massive in this case (5.1).
Fig. 11
From the point of view of the brane configurations describing four dimensional gauge theories like that of fig. 3, one can change $`\alpha `$ from $`\pi `$ in a number of ways. One is to change the relative position of the $`NS5`$ and $`NS5^{}`$-branes in $`x^7`$ by the amount $`\mathrm{\Delta }x^7`$. In the gauge theory on the $`D4`$-branes this corresponds to turning on a Fayet-Iliopoulos D-term \[6,,1\]. Depending on the sign of the FI term, either $`Q`$ or $`\stackrel{~}{Q}`$ should condense to minimize the D-term potential
$$V_D(Q^{}Q\stackrel{~}{Q}^{}\stackrel{~}{Q}r)^2$$
where $`r`$ is proportional to $`\mathrm{\Delta }x^7`$. The gauge theory analysis is nicely reproduced by our string theory considerations<sup>13</sup> Both in gauge theory and in string theory, in the context of the full configuration of fig. 3 the foregoing discussion is valid for $`N_fN_c`$ (see for details).. $`\mathrm{\Delta }x^7>0`$ corresponds to $`0<\alpha <\pi `$. In this case the ground state of the $`46_+`$ string is tachyonic, $`Q`$ condenses and the $`D4`$-branes detach from the $`NS5^{}`$-branes and attach to the $`D6`$-branes. $`\mathrm{\Delta }x^7<0`$ corresponds to $`\alpha >\pi `$ for the $`46_+`$ string which is hence massive, but since $`0<\alpha <\pi `$ for the $`46_{}`$ strings ($`\stackrel{~}{Q}`$), a similar process of condensation to the above occurs for them.
Another way of changing $`\alpha `$ is to tilt the $`D6`$-branes by some angle in the $`(x^6,x^7)`$ plane, which must lead to a potential similar to (5.1).
More comments on the supersymmetric case $`\alpha =\pi `$:
(1) Like in , for values of $`j`$ that satisfy the unitarity constraint (3.1), the full spectrum obtained from (3.1), (5.1) and its generalization to other observables is non-tachyonic.
(2) The excited twist field (5.1) does not create from the vacuum massless states; this is shown in appendix B.
(3) Massive states form hypermultiplets, as they should; this is also shown in appendix B.
(4) Like in the closed string case \[3,,13\], there is a continuum of $`\delta `$-function normalizable $`46^{}`$ string states corresponding to $`j1/2+i\mathrm{IR}`$. These states are separated by an energy gap of order $`1/l_s`$ from the massless discrete states discussed above.
5.2. $`D4`$-branes ending on $`NS5`$-branes from opposite sides
We next turn to the configuration of fig. 4b. A stack of $`N_L`$ $`D4`$-branes ends from the left on $`k`$ coincident $`NS5`$-branes, and a stack of $`N_R`$ $`D4`$-branes ends on the $`NS5`$-branes from the right. In this case, in the parametrization (3.1), the $`D4`$-branes on the left intersect the group manifold at $`g=1`$ while the ones on the right intersect it at $`g=1`$ (see fig. 12).
Fig. 12
We first consider the case where the two stacks of $`D4`$-branes are at the same point in the $`(x^4,x^5)`$ plane. To construct the vertex operator for emitting the lowest lying $`4_L4_R`$ string in the geometry of fig. 12, we can follow the discussion of the previous subsection. The dependence on $`(x^0,x^1,x^2,x^3,\varphi )`$ is the same as in the $`46^{}`$ case. In the $`(x^4,x^5)`$ directions we now have Dirichlet-Dirichlet boundary conditions. Therefore, the $`44`$ vertex operator is:
$$V_{44}=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }V_2$$
The contribution $`V_2`$ of the $`SU(2)`$ SCFT to (5.1) is the same as in section 5.1, with $`\alpha =\pm 2\pi `$:
$$V_2^\pm =\sigma _{2\pi }=e^{\pm i\sqrt{\frac{k}{2}}Y}$$
Hence
$$V_{44}^\pm =e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }e^{\pm i\sqrt{\frac{k}{2}}Y}$$
To describe $`V_{44}^\pm `$ in the cigar background (3.1), we use eqs. (5.1), (5.1), with $`\alpha =2\pi `$. The vertex operators (5.1) become in this geometry:
$$V_{44}^\pm =e^\phi e^{ik_\mu x^\mu }V_{j,\pm \frac{k}{2}}$$
The mass shell condition reads:
$$\frac{k_\mu ^2}{2}+\frac{k}{4}\frac{j(j+1)}{k}=\frac{1}{2}k_\mu ^2=\frac{2}{k}\left[(\frac{k}{2}1)\frac{k}{2}j(j+1)\right]$$
As in and section 5.1, the two point function of (5.1) has a series of poles corresponding to discrete representations of $`SL(2)`$, and the lowest lying state corresponds to
$$j=|m|1=\frac{k}{2}1$$
Plugging (5.1) into (5.1) we find that in this case $`V_{44}^\pm `$ create massless states from the vacuum.
The vertex operators which couple to these massless particles take the form (5.1), (5.1):
$$V_{44}^\pm (k_\mu ^2=0)=e^\phi e^{ik_\mu x^\mu }V_{\frac{k}{2}1,\pm \frac{k}{2}}$$
As expected from gauge theory, $`V_{44}^\pm `$ describe two complex scalar fields $`Q`$, $`\stackrel{~}{Q}^{}`$, in the $`(N_L,N_R)`$ of the $`U(N_L)\times U(N_R)`$ symmetry on the $`D`$-branes, respectively. The complex conjugates $`Q^{}`$, $`\stackrel{~}{Q}`$ arise from $`4_L4_R`$ strings with the opposite orientation.
In this case, the charges of $`Q,\stackrel{~}{Q}`$ under the $`SO(2)_{45}\times SO(2)_{89}`$ R-symmetry are $`(0,\frac{k}{2})`$. The $`SO(2)_{45}`$ agrees with expectations. The $`SO(2)_{89}`$ charge predicted by brane theory is $`1/2`$, which is the value obtained by extrapolating our results to $`k=1`$ (which, as discussed above, is in fact outside the range of validity of our approach). For $`k>1`$, the situation is less clear, but it is worth pointing out that in addition to the fact that the same issue already arose in the previous subsection and in , it also appeared in brane theory before. As reviewed in , the theory on $`D4`$-branes ending on $`k`$ $`NS5`$-branes typically contains chiral superfields $`\mathrm{\Phi }`$ with a polynomial superpotential $`s_0\mathrm{\Phi }^{k+1}`$. The coupling $`s_0`$ appears to be charged under the analog of $`SO(2)_{89}`$. The resolution of all these problems is again left for future work.
Comments:
(1) If the $`N_L`$ $`D4`$-branes are located at $`x^4=x^5=0`$, and the $`N_R`$ $`D4`$-branes are at $`(x^4,x^5)=(a,b)`$, one can apply the discussion of section 4.1. As in (4.1), the vertex operators $`V_{44}^\pm `$ have to be multiplied by the “winding” generating factor $`\mathrm{exp}\{\frac{i}{\pi }[a(x_L^4x_R^4)+b(x_L^5x_R^5)]\}`$. The lowest lying $`4_L4_R`$ states have mass squared $`\frac{1}{\pi ^2}(a^2+b^2)`$, in agreement with expectations.
(2) For the case $`\alpha =\pm 2\pi `$, the $`SU(2)`$ current algebra is untwisted (see (5.1)). The discussion of section 4.2 shows that the $`D4_L`$ and $`D4_R`$ branes correspond in this case to the boundary states with $`j=0`$ and $`j=\frac{k_B}{2}`$, respectively. Hence, the strings that connect the two transform in the $`j=\frac{k_B}{2}`$ representation of the bosonic $`\widehat{SU(2)}`$. Denoting the corresponding vertex operators by $`\sigma _{\frac{k_B}{2},m}`$ ($`m=\frac{k_B}{2},\frac{k_B}{2}+1,\mathrm{},\frac{k_B}{2}1,\frac{k_B}{2}`$), we find that $`V_2`$ in (5.1), (5.1) is given by
$$V_2^\pm =e^{\pm i\sqrt{\frac{k}{2}}Y}=\chi ^\pm \sigma _{2\pi }^B=\chi ^\pm \sigma _{\frac{k_B}{2},\pm \frac{k_B}{2}}$$
The last expression can be thought of as the highest and lowest weight states in a spin $`\frac{k}{2}=\frac{k_B}{2}+1`$ representation of the total $`SU(2)`$. The full set of primaries at the lowest level of $`4_L4_R`$ strings is the following:
$$\begin{array}{cc}\hfill (i)& e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\left(\chi \sigma _{\frac{k_B}{2}}\right)_{j=\frac{k}{2}}\hfill \\ \hfill (ii)& e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\left(\chi \sigma _{\frac{k_B}{2}}\right)_{j=\frac{k}{2}2}\hfill \\ \hfill (iii)& e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\left[\underset{M}{}a_M\psi ^M\sigma _{\frac{k_B}{2}}+a_8\left(\chi \sigma _{\frac{k_B}{2}}\right)_{j=\frac{k}{2}1}\right]\hfill \end{array}$$
In $`(i)`$ and $`(ii)`$, $`\chi `$ and $`\sigma _{\frac{k_B}{2}}`$ are coupled to a representation of total spin $`\frac{k}{2}`$ and $`\frac{k}{2}2`$, respectively. $`(iii)`$ contains six physical combinations (after imposing the physical state conditions on the polarization coefficients $`a_M`$, and eliminating null states); $`\{\psi ^M\}=\{\psi ^\mu ,\psi ^4,\psi ^5,\chi ^r\}`$. It can be shown that out of the operators (5.1) only those whose $`SU(2)`$ part is (5.1) create massless particles (see appendix C).
(3) One can study $`4_L4_R`$ strings when the angle between the two kinds of branes is generic, $`\alpha /2`$. For $`\alpha /2\pm \pi `$, an analysis similar to the one in section 5.1 shows that the lowest lying $`4_L4_R`$ states are tachyonic. For $`\alpha =0`$, one finds a system of parallel $`D4\overline{D}4`$. This is discussed in the next subsection.
5.3. Rotating $`D4D4`$ systems into $`D4\overline{D}4`$
The vertex operators creating the low lying states of $`4_L4_R`$ strings when the angle between the $`D4_L`$ and $`D4_R`$-branes is $`\pi <\alpha /2<\pi `$ are (see subsection 5.1):
$$V_{44}(\alpha /2)=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }e^{i\frac{\alpha }{2\pi }\sqrt{\frac{k}{2}}Y}$$
Repeating the analysis of sections 5.1, 5.2, one finds that the lowest lying states that (5.1) creates from the vacuum are tachyonic. At $`\alpha =0`$, when the two $`D4`$-branes are anti-parallel,
$$V_{44}(0)=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }$$
is the tachyon vertex operator creating the low lying state of a string connecting a $`D4`$-brane parallel to an anti-$`D4`$-brane ($`\overline{D}4`$-brane), both ending on the $`NS5`$-branes from the same side; it is a $`4_L\overline{4}_L`$ string.
As before, the tachyon signals that the $`D4`$-branes could reduce their energy by reconnecting. Imagine for simplicity that the ends of the $`D4`$-branes are pinned down far from the $`NS5`$-branes (see fig. 13). Tachyon condensation corresponds to a process where the $`D4`$-branes (solid lines) connect to each other and detach from the fivebrane. They can then reduce their energy by stretching straight between the two pinned endpoints (dashed line). As is clear from figure 13, their energy decreases in the process. As $`\alpha 0`$, the fourbranes annihilate (or more generally become lower dimensional branes). This generalizes the results of on brane annihilation to $`D`$-branes ending on fivebranes.
Fig. 13
Similarly, the vertex operator creating the low lying states of $`4_L4_L`$ strings when the angle between the $`D4_L`$-branes is $`\alpha /2`$ is
$$V_{44}(\alpha /2)=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }e^{i[(\frac{\alpha }{|\alpha |}\frac{\alpha }{2\pi })H\sqrt{\frac{k_B}{2}}\frac{\alpha }{2\pi }u]}$$
(the last exponent is the excited twist field (5.1)). The extra factor of $`\mathrm{exp}(\pm iH)`$ in (5.1) compared to (5.1) is due to the different signs of the GSO projection in the two cases resulting from the reversed orientation of the $`D4_R`$-brane. At $`\alpha =0`$, we obtain the $`4_L4_L`$ vertex operator which in the cigar geometry is (recall (5.1) and see appendix B for notation)
$$V_{44}^\pm (0)=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\chi ^\pm e^\phi e^{ik_\mu x^\mu }z_1^{k2}V_{j,\pm 1}$$
(the $`\pm `$ in (5.1) is correlated with $`\alpha 0_+`$ or $`\alpha 0_{}`$ in (5.1)). At first sight one might be puzzled why for two parallel $`D4`$-branes ending on a stack of $`NS5`$-branes one does not find additional massless particles. For example, in section 2 it was argued that fundamental strings stretching between the fourbranes give rise to massless non-Abelian gauge bosons on the fourbranes. The reason we are not supposed to see those here is that we are studying the physics of the $`D`$-branes in the near-horizon geometry of the fivebranes. Only states whose wavefunctions are bound to the fivebranes can give rise to normalizable modes in our analysis. The wavefunctions of the gauge bosons discussed in section 2 are in contrast spread out in $`x^6`$ and, in particular, are not localized at the fivebranes. From the point of view of the near-horizon geometry, they are non-normalizable.
At $`\alpha /2=\pm \pi `$ (5.1) turns into
$$V_{44}(\pm \pi )=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\sigma _{\pm 2\pi }^B$$
which is the tachyon vertex operator for a $`4_L\overline{4}_R`$ string connecting a $`D4`$-brane and a $`\overline{D}4`$-brane located on opposite sides of the stack of $`NS5`$-branes (see fig. 14).
Fig. 14
Brane constructions involving $`D4`$ and $`\overline{D}4`$-branes ending on $`NS5`$-branes were recently studied in \[43,,44\]. The present work can be used to shed more light on such configurations.
5.4. A $`D6`$-brane intersecting a stack of $`NS5`$-branes
We now turn to the configuration of fig. 4c. A stack of $`D4`$-branes ends from the left on $`k`$ coincident $`NS5`$-branes, and a $`D6`$-brane intersects the $`NS5`$-branes at the same value of $`(x^4,x^5)`$ as the fourbranes.
From the brane geometry point of view in Type IIA, the notion of putting a $`D6`$-brane on top of $`k`$ coincident $`NS5`$-branes is not well defined, due to the HW transition . For instance, for a single $`NS5`$-brane one must specify if the $`D6`$-brane is located to the left of the $`NS5`$-brane or to its right. In the first case, the $`D6`$-brane can be moved freely away from the $`NS5`$-brane to the left in the $`x^6`$ direction, but if moved away to the right an extra $`D4`$-brane is created, stretched between the $`D6`$-brane and the $`NS5`$-brane. In the second case, the $`D4`$-brane will be created when moving the $`D6`$-brane to the left. The configuration space of this system is thus separated into two disconnected components.
Similarly, for the case of $`k`$ $`NS5`$-branes one has to specify which of these fivebranes are to the left of the $`D6`$-brane and which of them are located to its right. If $`nk`$ fivebranes are to the left of the $`D6`$-brane, $`n`$ $`D4`$-branes will be created when the $`D6`$-brane moves to the left away from the stack of $`NS5`$-branes, and $`kn`$ $`D4`$-branes will be formed if it is taken away to the right. Thus the configuration space consists of $`k+1`$ sectors – there are $`k+1`$ possibilities of what is meant by placing a $`D6`$-brane on top of $`k`$ coincident fivebranes.
This ambiguity in type IIA string theory can be understood in M-theory on $`\mathrm{IR}^{10}\times S^1`$ when the radius of the $`S^1`$ is large in Planck units . The $`D6`$-brane becomes a Kaluza-Klein monopole, i.e. a bundle whose fiber is the eleventh circular coordinate $`x^{10}`$ of radius $`R_{10}`$. The bundle is non-trivial and two patches are needed to trivialize it. These patches can be chosen as the two halves of the ten dimensional space, one with $`x^60`$ to the right of the sixbrane and the other with $`x^60`$ to its left. Denote by $`x_\pm ^{10}`$ the fiber coordinate over the two different patches. If $`\vartheta `$ is the azimuthal angle in the $`(x^4,x^5)`$ plane, i.e. $`\mathrm{tan}(\vartheta )=\frac{x^5}{x^4}`$, then the transition between the fiber coordinates in the two patches on the transition $`(x^4,x^5)`$ plane at $`x^6=0`$ is: $`\frac{x_+^{10}}{R}=\frac{x_{}^{10}}{R_{10}}+\vartheta `$. From eleven dimensional standpoint, an $`NS5`$-brane is an $`M5`$-brane which is a point in $`x^{10}`$, stretching along the $`(x^4,x^5)`$ plane. Starting with $`k`$ $`NS5`$-branes in presence of a $`D6`$-brane, one must specify which of these $`M5`$-branes have a definite $`x_+^{10}`$ coordinate, and which have a well defined $`x_{}^{10}`$ coordinate.
The division of configuration space into sectors is visible also in the worldsheet CFT description of section 4. Unlike the $`D6^{}`$-brane case studied in section 5.1, the $`D6`$-brane intersects the group manifold (3.1) not at a point but along a two-sphere of constant $`x^6`$. This corresponds to the conjugacy class $`C_\theta `$ (4.1), where $`\theta `$ satisfies $`\mathrm{cos}(\frac{\theta }{2})=\frac{x^6}{|\stackrel{}{x}|}`$. In section 4 we have seen that consistent boundary conditions of this type exist only for $`k1`$ values of $`\theta `$ given by eq. (4.1). The $`D6`$-brane is thus allowed to intersect the stack of $`k`$ $`NS5`$-branes only at fixed quantized values of $`x^6`$, namely, on two-spheres with quantized sizes $`S_n^2`$, $`n=1,2,\mathrm{},k1`$ (see fig. 15).
Fig. 15
It is natural to identify these $`k1`$ different boundary states with the various possible positions of the $`D6`$-brane among the $`k`$ $`NS5`$-branes, discussed above. There we had $`k+1`$ possible states, two of which are absent in the CFT approach. A possible interpretation of this fact is that when the $`D6`$-brane is either to the right or to the left of all $`NS5`$-branes, it is not bound to the fivebranes and hence does not give rise to a boundary state in the CHS geometry.
Anyhow, we can again study the massless $`46`$ open strings when the $`D6`$-brane is characterized by $`\theta =2\pi \frac{l}{k_B}`$ in (4.1), with $`2l=0,1,\mathrm{},k_B`$. Such a string connects the $`g=1`$ boundary state with $`l=0`$, corresponding to the $`D4`$-brane, to a spin $`l`$ boundary state corresponding to the $`D6`$-brane (see fig. 16). The $`46`$ string belongs to representations contained in the fusion of spin $`0`$ and spin $`l`$, i.e. the spin $`l`$ representation. This consists of $`2l+1`$ operators $`\sigma _{lm}`$, $`m=l,l+1,\mathrm{},l1,l`$, with scaling dimension
$$h(\sigma _{lm})=\frac{l(l+1)}{k}$$
As in subsection 5.2, we get for each $`l`$ the vertex operators:
$$V_{lm}=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\left(\chi \sigma _l\right)_{l+1,m}$$
and some additional operators that do not create massless states. The two operators $`V_{lm}`$ (5.1) with $`m=\pm (l+1)`$ do give rise to poles at $`k_\mu ^2=0`$ in amplitudes and thus create massless particles of charge $`\pm (l+1)`$.
Fig. 16
Recently, it was argued that $`n`$ $`D0`$-branes in an $`SU(2)`$ WZW background, which are $`n`$ points on $`S^3`$, can turn into a single $`D2`$-brane which wraps a two-sphere with a quantized radius $`S_n^2`$ \[37,,38\]. This meshes nicely with the discussion above. In the near horizon of $`k`$ $`NS5`$-branes, a $`D4`$-brane is a point on $`S^3`$ while a $`D6`$-brane wraps an $`S_n^2S^3`$, $`n=1,2,\mathrm{},k1`$. Hence, $`n`$ $`D4`$-branes can condense into a single $`D6`$-brane in the $`\theta =\pi \frac{n1}{k2}`$ conjugacy class (4.1), namely, a $`D6`$-brane wrapping the sphere $`S_n^2`$. This is the HW transition. The $`D6`$-brane on $`S_n^2`$ is located, say, to the left of $`n`$ out of the $`k`$ $`NS5`$-branes. Moving it to the right, past these $`n`$ fivebranes, creates $`n`$ $`D4`$-branes stretched from the $`D6`$-brane to each of the $`n`$ fivebranes. In the near horizon limit these look like $`n`$ points on $`S^3`$.
Acknowledgements: We thank M. Berkooz, A. Hanany, N. Itzhaki, Y. Oz and O. Pelc for useful discussions. We also thank the ITP in Santa Barbara for hospitality during the course of this work. S.E. thanks the IAS in Princeton. This research was supported in part by NSF grant #PHY94-07194 and by the Israel Academy of Sciences and Humanities – Centers of Excellence Program. The work of A.G. and E.R. is also supported in part by BSF – American-Israel Bi-National Science Foundation. D.K. is supported in part by DOE grant #DE-FG02-90ER40560.
Appendix A. $`V_{jm}V_{j,m}`$ on the disc
We first present some useful formulae:
$$C(a,b)=_{\mathrm{}}^{\mathrm{}}|x|^{a1}|1x|^{b1}𝑑x=B(a,1ab)+B(a,b)+B(1ab,b)$$
Here $`B`$ is the Euler beta function:
$$B(a,b)=_0^1x^{a1}(1x)^{b1}𝑑x=\frac{\mathrm{\Gamma }(a)\mathrm{\Gamma }(b)}{\mathrm{\Gamma }(a+b)}$$
where
$$\mathrm{\Gamma }(a+1)=a\mathrm{\Gamma }(a),\mathrm{\Gamma }(n+ϵ)=\frac{()^n}{n!ϵ}+O(1),n=0,1,2,\mathrm{}$$
To find the analytic structure of the two point functions of boundary operators discussed in this work, like (5.1), we need to compute the two point functions $`V_{jm}V_{j,m}`$ on the disc. This is the purpose of this appendix.
The calculation is similar to its closed string analog \[17,,13\], which we follow below. The CFT on the cigar is obtained by the coset construction from that on $`AdS_3`$. The natural observables in CFT on the Euclidean version of $`AdS_3`$ on the disc are functions $`\mathrm{\Phi }_j(x;z)`$ which transform as primaries under the diagonal (left $`+`$ right) $`SL(2)`$ current algebra (see section 4.2). $`x`$ is an auxiliary real variable.
The two point function of $`\mathrm{\Phi }_j`$ is<sup>14</sup> Here and below we suppress the dependence of correlation functions on the worldsheet location $`z`$ on the real line, the boundary of the upper half plane.
$$\mathrm{\Phi }_j(x)\mathrm{\Phi }_j(x^{})=A(j,k)|xx^{}|^{2(j+1)}$$
where $`A(j,k)`$ is an analytic function of $`j`$ in the domain (3.1); its precise form will not play a role below.
To study the coset it is convenient to “Fourier transform” the fields $`\mathrm{\Phi }_j(x)`$ and define the mode operators
$$\mathrm{\Phi }_{jm}=_{\mathrm{}}^{\mathrm{}}\mathrm{\Phi }_j(x)|x|^{j+m}𝑑x$$
The two point functions of the modes $`\mathrm{\Phi }_{jm}`$ are equal to those of the $`SL(2)/U(1)`$ coset theory:
$$V_{jm}V_{j,m}=\mathrm{\Phi }_{jm}\mathrm{\Phi }_{j,m}$$
Using (A.1), (A.1) and (A.1) one finds that
$$V_{jm}V_{j,m}=A^{}(j,k)_{\mathrm{}}^{\mathrm{}}|y|^{jm}|1y|^{2(j+1)}𝑑y=C(jm+1,2j1)$$
Using eqs. (A.1), (A.1), (A.1) one can finally express the two point functions in terms of gamma functions as:
$$\begin{array}{cc}& V_{jm}V_{j,m}=A^{}(j,k)\times \hfill \\ & \left(\frac{\mathrm{\Gamma }(jm+1)\mathrm{\Gamma }(j+m+1)}{\mathrm{\Gamma }(2j+2)}+\frac{\mathrm{\Gamma }(jm+1)\mathrm{\Gamma }(2j1)}{\mathrm{\Gamma }(jm)}+\frac{\mathrm{\Gamma }(j+m+1)\mathrm{\Gamma }(2j1)}{\mathrm{\Gamma }(j+m)}\right)\hfill \end{array}$$
Using eq. (A.1) we can now find the analytic structure of $`V_{jm}V_{j,m}`$: it has single poles for $`j,m`$ satisfying (3.1).<sup>15</sup> For $`j=m`$ one finds that the three terms in (A.1) conspire to give $`0`$. This implies that the extra poles at $`j=m`$, found in some special cases in , appear on the sphere but not on the disc.
Appendix B.
In this appendix we show that the excited twist field (5.1) does not give rise to massless particles even for $`\alpha =\pi `$, and in the case $`\alpha =\pi `$ we verify that the massive states created by (5.1) and (5.1) have degeneracy two, and are hence organized into hypermultiplets.
The twist field (5.1) has a conformal weight
$$h(\sigma _\alpha ^{})=h(\sigma _\alpha )+\frac{1}{2}\left(1\frac{\alpha }{\pi }\right)=\frac{k}{4}\left(\frac{\alpha }{2\pi }\right)^2+\frac{1}{2}\left(1\frac{\alpha }{\pi }\right)$$
and a $`J_{\mathrm{total}}^3`$ charge (recall (3.1), (5.1))
$$m(\sigma _\alpha ^{})=m(\sigma _\alpha )+1=\frac{k}{2}\frac{\alpha }{2\pi }+1$$
Its decomposition on $`SU(2)/U(1)\times U(1)`$ thus reads
$$\sigma _\alpha ^{}=z_1^{k2}e^{i\sqrt{\frac{2}{k}}(\frac{k}{2}\frac{\alpha }{2\pi }1)Y}$$
where $`z_1^{k2}`$ (the notation will become clear soon) is an operator in the $`N=2`$ minimal model $`SU(2)_k/U(1)`$ with
$$h(z_1^{k2})=h(\sigma _\alpha ^{})\frac{2}{k}\frac{(\frac{k}{2}\frac{\alpha }{2\pi }1)^2}{2}=\frac{1}{2}\frac{k2}{k}$$
It can be shown that $`z_1^{k2}`$ is a chiral operator in the $`N=2`$ minimal model with a $`U(1)_R`$ charge $`R(z_1^{k2})=(k2)/k`$; it is the highest charge operator in the chiral ring of $`SU(2)_k/U(1)`$, $`\{z_1^i|i=0,1,\mathrm{},k2\}`$, $`R(z_1^i)=i/k`$.
Collecting the above, recalling (5.1) (with $`\sigma _\alpha `$ replaced by $`\sigma _\alpha ^{}`$) and (5.1), (5.1), (5.1), we find that the $`46^{}`$ vertex operator under consideration, and what it becomes in the cigar geometry (3.1), is
$$V_{46^{}}^{}=e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }e^{\beta \varphi }\sigma _\alpha ^{}e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }z_1^{k2}V_{j,\frac{k}{2}\frac{\alpha }{2\pi }+1}$$
The lowest lying state corresponds to $`n=1`$ in (3.1), namely, $`j=|m|1`$; its mass is
$$(M_\alpha ^{})^2=\frac{1}{2}+\frac{2}{k}j=\frac{1}{2}\frac{2}{k}+\left|\frac{\alpha }{2\pi }\frac{2}{k}\right|$$
Since $`j>1/2`$ (3.1), we have
$$(M_\alpha ^{})^2>\frac{1}{2}\frac{1}{k}0$$
for any $`\alpha `$ allowed in the unitarity range.
A particular case is the supersymmetric $`\alpha =\pi `$ configuration. In this case, the twist operators (5.1) and (5.1) take the form
$$\begin{array}{cc}\hfill \sigma _\pi ^{}& =e^{\frac{i}{2}H}e^{\frac{i}{2}\sqrt{\frac{k_B}{2}}u}=e^{\frac{i}{2}\sqrt{\frac{k}{2}}Y}\hfill \\ \hfill \sigma _\pi ^+& =e^{\frac{i}{2}H}e^{\frac{i}{2}\sqrt{\frac{k_B}{2}}u}=z_1^{k2}e^{i\sqrt{\frac{2}{k}}(\frac{k}{4}1)Y}\hfill \end{array}$$
The corresponding $`46^{}`$ vertex operators turn in the cigar geometry (3.1) into
$$\begin{array}{cc}\hfill V_{46^{}}^{}& =e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }e^{\beta \varphi }\sigma _\pi ^{}e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }V_{j,m_{}}\hfill \\ \hfill V_{46^{}}^+& =e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }e^{\beta \varphi }\sigma _\pi ^+e^\phi \sigma _{45}S_{45}e^{ik_\mu x^\mu }z_1^{k2}V_{j,m_+}\hfill \end{array}$$
where
$$m_\pm =\frac{k2}{4}\pm \frac{1}{2}$$
Physical states in spacetime are obtained when $`j`$ in $`V_{46^{}}^\pm `$ is (3.1)
$$j_\pm ^n=|m_\pm |n,n=1,2,\mathrm{}$$
respectively. Notice that<sup>16</sup> Below we restrict to the case $`k4`$; for $`k=2,3`$ there are no massive excitations in the range (3.1).
$$j_+^n=j_{}^{n+1}$$
The on-shell conditions $`h(V_{46^{}}^\pm )=1`$ imply
$$\frac{(M_\pm ^n)^2}{2}=\frac{k4}{16}\frac{j_\pm ^n(j_\pm ^n+1)}{k}$$
Using (B.1), (B.1), (B.1), (B.1), one finds
$$(M_{}^n)^2=\frac{1}{k}(n1)(k2n)$$
$$(M_+^n)^2=(M_{}^{n+1})^2=\frac{1}{k}n(k22n)$$
The unitarity bound (3.1) restricts $`n`$ to a certain range. In this range $`(M_+^n)^2>0`$ for any $`n`$, while $`(M_{}^n)^20`$ for any $`n`$ and equality is satisfied iff $`n=1`$. We thus see (B.1) that there is a degeneracy two for all massive states; they get organized into $`4d`$ hypermultiplets $`(Q,\stackrel{~}{Q})`$, as they should. The only massless state $`M_{}^1=0`$ (B.1) is non-degenerate; it corresponds to the chiral superfield $`Q`$ considered in section 5.1.
Appendix C.
In this appendix we show that, except for (5.1), the other operators in (5.1) do not create massless states from the vacuum.
The operators in (5.1) are linear combinations of
$$V_m^a=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\chi ^a\sigma _{\frac{k_B}{2},m},a=\pm ,3$$
and
$$V_m^M=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\psi ^M\sigma _{\frac{k_B}{2},m},M=\mu ,4,5,r$$
The operators $`\sigma _{\frac{k_B}{2},m}`$ in (C.1) and (C.1) carry $`m`$ units of charge under $`J^3`$ and have scaling dimension
$$h(\sigma _{\frac{k_B}{2},m})=\frac{1}{k}\frac{k_B}{2}(\frac{k_B}{2}+1)=\frac{k2}{4}$$
The decomposition of $`\sigma _{jm}`$ on $`SU(2)/U(1)\times U(1)`$ reads
$$\sigma _{jm}=V_{jm}^{}e^{i\sqrt{\frac{2}{k}}mY}$$
where $`V_{jm}^{}`$ is an operator in the $`SU(2)/U(1)`$ SCFT with scaling dimension
$$h(V_{jm}^{})=\frac{j(j+1)m^2}{k}$$
Now, going from the cylinder $`(\varphi ,Y)`$ to the $`SL(2)/U(1)`$ cigar we take
$$e^{\beta \varphi }e^{i\sqrt{\frac{2}{k}}mY}V_{jm}$$
(see (5.1), (5.1)). Altogether, in the background (3.1), the operators (C.1) take the form
$$V_m^M=e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\psi ^M\sigma _{\frac{k_B}{2},m}e^\phi e^{ik_\mu x^\mu }\psi ^MV_{\frac{k_B}{2},m}^{}V_{jm}$$
When $`k_\mu ^2=0`$, the on-shell condition for $`V_m^M`$ reads
$$\frac{1}{k}\frac{k_B}{2}(\frac{k_B}{2}+1)\frac{j(j+1)}{k}=0j=\frac{k_B}{2}$$
Since $`|m|\frac{k_B}{2}`$, the condition (3.1) is not satisfied. Hence, no massless particle is emitted by $`V_m^M`$. Similarly, $`V_m^3`$ in (C.1) does not create massless states from the vacuum.
The operators $`V_m^\pm `$ in (C.1) include a factor $`\chi ^\pm \sigma _{\frac{k_B}{2},m}`$ which decomposes as
$$\chi ^\pm \sigma _{\frac{k_B}{2},m}=\mathrm{\Psi }_me^{i\sqrt{\frac{2}{k}}(m\pm 1)Y}$$
where $`\mathrm{\Psi }_m`$ is in the $`SU(2)/U(1)`$ SCFT and has
$$h(\mathrm{\Psi }_m)=\frac{k}{4}\frac{(m\pm 1)^2}{k}$$
In the geometry (3.1), the operators $`V_m^\pm `$ now take the form
$$V_m^\pm =e^\phi e^{ik_\mu x^\mu }e^{\beta \varphi }\chi ^\pm \sigma _{\frac{k_B}{2},m}e^\phi e^{ik_\mu x^\mu }\mathrm{\Psi }_mV_{j,m\pm 1}$$
When $`k_\mu ^2=0`$, the on-shell condition gives again $`j=\frac{k_B}{2}`$ (C.1). Since $`m\pm 1`$ take values between $`\frac{k_B}{2}\pm 1`$ and $`\frac{k_B}{2}\pm 1`$, only the operators $`V_{\frac{k_B}{2}}^+`$ and $`V_{\frac{k_B}{2}}^{}`$ – the operators in (5.1) whose $`SU(2)`$ part is (5.1) – satisfy (3.1). Therefore, the other operators $`V_m^\pm `$ in (C.1) do not emit massless particles.
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# 1 W-pair production at LEP2
## 1 W-pair production at LEP2
The investigation of W-pair production at LEP2 plays an important role in the verification of the Electroweak Standard Model (SM). Apart from the direct observation of the triple-gauge-boson couplings in $`\mathrm{e}^+\mathrm{e}^{}\mathrm{W}^+\mathrm{W}^{}`$, the increasing accuracy in the W-pair-production cross-section and W-mass measurements has put this process into the row of SM precision tests.
To account for the high experimental accuracy $`^{\mathrm{?},\mathrm{?}}`$, on the theoretical side is a great challenge: the W bosons have to be treated as resonances in the full four-fermion processes $`\mathrm{e}^+\mathrm{e}^{}4f`$, and radiative corrections need to be included. While several lowest-order predictions are based on the full set of Feynman diagrams, only very few calculations include radiative corrections beyond the level of universal effects (see Refs. ?,? and references therein). Fortunately, to match the experimental precision for W-pair production at LEP2 a full one-loop calculation for the four-fermion processes is not needed, and it is sufficient to take into account only those radiative corrections that are enhanced by two resonant W bosons. A naive estimate of the neglected corrections yields $`(\alpha /\pi )(\mathrm{\Gamma }_\mathrm{W}/M_\mathrm{W})0.5\%`$. The theoretically clean way to carry out this approximation is the expansion about the two resonance poles, which is called double-pole approximation (DPA) $`^\mathrm{?}`$. A full description of this strategy and of different variants used in the literature $`^{\mathrm{?},\mathrm{?},\mathrm{?}}`$ (some of them involving further approximations) is beyond the scope of this article. We can only briefly sketch the approach pursued in RacoonWW $`^{\mathrm{?},\mathrm{?}}`$.
## 2 Radiative corrections with RacoonWW
In DPA, $`𝒪(\alpha )`$ corrections to $`\mathrm{e}^+\mathrm{e}^{}\mathrm{WW}4f`$ can be classified into two types: factorizable and non-factorizable corrections. We first focus on virtual corrections.
Factorizable corrections are those that correspond either to W-pair production or to W decay. Virtual factorizable corrections are represented by the schematic diagram on the l.h.s. of Fig. 1, in which the shaded blobs contain all one-loop corrections to the on-shell production and on-shell decay processes, and the open blobs include the corrections to the W propagators.
For the corresponding matrix element $``$ the application of the DPA amounts to the replacement
$$=\frac{R(k_{\mathrm{W}^+}^2,k_\mathrm{W}^{}^2)}{(k_{\mathrm{W}^+}^2M_\mathrm{W}^2)(k_\mathrm{W}^{}^2M_\mathrm{W}^2)}\frac{R(M_\mathrm{W}^2,M_\mathrm{W}^2)}{(k_{\mathrm{W}^+}^2M_\mathrm{W}^2+\mathrm{i}M_\mathrm{W}\mathrm{\Gamma }_\mathrm{W})(k_\mathrm{W}^{}^2M_\mathrm{W}^2+\mathrm{i}M_\mathrm{W}\mathrm{\Gamma }_\mathrm{W})},$$
(1)
where the originally gauge-dependent numerator $`R(k_{\mathrm{W}^+}^2,k_\mathrm{W}^{}^2)`$ is replaced by the gauge-independent residue $`R(M_\mathrm{W}^2,M_\mathrm{W}^2)`$. The one-loop corrections to this residue can be deduced from the known results for the pair production and the decay of on-shell W bosons. However, the spin correlations between the two W decays should be taken into account.
Non-factorizable corrections $`^\mathrm{?}`$ comprise all those doubly-resonant corrections that are not yet contained in the factorizable ones, and include, in particular, all diagrams involving particle exchange between the subprocesses. Such diagrams only lead to doubly-resonant contributions if the exchanged particle is a photon with energy $`E_\gamma \stackrel{<}{}\mathrm{\Gamma }_\mathrm{W}`$; all other non-factorizable diagrams are negligible in DPA. A typical diagram for a virtual non-factorizable correction is shown on the r.h.s. in Fig. 1, where the full blob represents tree-level subgraphs. We note that diagrams involving photon exchange between the W bosons contribute both to factorizable and non-factorizable corrections; otherwise the splitting into those parts would not be gauge-invariant.
In RacoonWW the virtual corrections are treated in DPA, including the full set of factorizable and non-factorizable $`𝒪(\alpha )`$ corrections. The real corrections are calculated from full matrix elements for $`\mathrm{e}^+\mathrm{e}^{}4f\gamma `$, as described in Ref. ?, i.e. the DPA is not used in this part. In this way, we avoid potential problems in the definition of the DPA for the emission of photons with energies $`E_\gamma \mathrm{\Gamma }_\mathrm{W}`$. On the other hand, this asymmetry in the calculation of virtual and real corrections requires particular care concerning the structure of IR and mass singularities. The singularities have the form of a universal radiator function convoluted with the respective lowest-order matrix element squared $`|_0|^2`$ of the non-radiative process. Since the virtual corrections are calculated in DPA, but the full matrix element is used for the real photons, a simple summation of virtual and real corrections would lead to a mismatch in the singularity structure and eventually to totally wrong results. Therefore, we extract those singular parts from the real photon contribution that exactly match the singular parts of the virtual photon contribution, then replace in these terms the full $`|_0|^2`$ by the one calculated in DPA, and finally add this modified part to the virtual corrections. This modification is allowed within DPA and leads to a proper matching of all IR and mass singularities. This treatment has been carried out in two different ways, once following phase-space slicing, once using the subtraction formalism of Ref. ?.
Beyond $`𝒪(\alpha )`$, RacoonWW includes soft-photon exponentiation and leading higher-order ISR effects in the structure-function approach. Using $`G_\mu `$ as input parameter instead of $`\alpha (0)`$, also the leading effects from $`\mathrm{\Delta }\alpha `$ and $`\mathrm{\Delta }\rho `$ are absorbed and partially resummed in the lowest order.
## 3 Phenomenological results
A survey of numerical results obtained with RacoonWW has already been presented in Ref. ? for LEP2 and linear-collider energies. Here we only review the W invariant-mass distribution given there and extend the results for the total cross section.
Figure 2 (l.h.s.) shows the invariant-mass distribution of the $`\mathrm{d}\overline{\mathrm{u}}`$ quark pair for the semi-leptonic channel $`\mathrm{e}^+\mathrm{e}^{}\nu _\mu \mu ^+\mathrm{d}\overline{\mathrm{u}}`$ at $`\sqrt{s}=200\mathrm{GeV}`$ at tree-level and with electroweak $`𝒪(\alpha )`$ corrections for two different recombination cuts, $`M_{\mathrm{rec}}=5`$ and $`25\mathrm{GeV}`$.
The recombination of photons with final-state charged fermions is performed as described in Ref. ?: we first determine the lowest invariant mass $`M_{\gamma f}`$ built by an emitted photon and a charged final-state fermion. If $`M_{\gamma f}`$ is smaller than $`M_{\mathrm{rec}}`$, the photon momentum is added to the one of the corresponding fermion $`f`$. The maxima of the corrected line shapes differ by up to 30–40$`\mathrm{MeV}`$ for the two values of $`M_{\mathrm{rec}}`$. As expected, there is a tendency to shift the maxima to larger invariant masses if more and more photons are recombined. In Fig. 2 (r.h.s.) we display the relative corrections $`\delta =\mathrm{d}\sigma /\mathrm{d}\sigma _01`$ for the two values of $`M_{\mathrm{rec}}`$, which illustrates the strong dependence of the corrected invariant-mass distributions on the treatment of the real photons. We obtain consistent results for the phase-space “slicing” and the “subtraction” methods. The size of the shown effects demonstrates that a careful treatment of real photons is mandatory in the W-mass reconstruction at LEP2 accuracy.
Figure 3 shows a comparison of RacoonWW results and of other predictions with recent LEP2 data, as given by the LEP Electroweak Working Group $`^{\mathrm{?},\mathrm{?}}`$.
The data are in good agreement with the predictions of RacoonWW and YFSWW3 $`^\mathrm{?}`$. The predictions of these two generators differ between 0.5–0.7%.<sup>§</sup><sup>§</sup>§Meanwhile the dominant source of this difference has been found, and the new results of YFSWW3 are closer to the results of RacoonWW. Details on the new YFSWW3 predictions can be found in Ref. ?. More details on the conceptual differences of the two generators, as well as a detailed comparison of numerical results, can be found in Ref. ?. Figure 3 also includes the prediction provided by GENTLE $`^\mathrm{?}`$, which differs from the RacoonWW and YFSWW3 results by 2–2.5%. This difference is due to the neglect of non-leading, non-universal $`𝒪(\alpha )`$ corrections in GENTLE. In summary, the comparison between SM predictions with the precise measurements of the W-pair production cross section at LEP2 reveals evidence of non-leading electroweak radiative corrections beyond the level of universal effects.
## References
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# Probability distribution of the free energy of a directed polymer in a random medium
## 1 Introduction
Directed polymers in a random medium is one of the simplest systems for which the effect of strong disorder can be studied. At the mean field level, it possesses a low temperature phase, with a broken symmetry of replica similar to mean field spin glasses. The problem is however much better understood than spin glasses; in particular one can write closed expressions of the mean field free energy and one can predict the existence of phase transitions in all dimensions $`d+1>2+1`$. It is also an interesting system from the point of view of non-equilibrium phenomena: through the Kardar-Parisi-Zhang (KPZ) equation, it is related to ballistic growth models and, in $`1+1`$ dimensions, to the asymmetric simple exclusion process (ASEP).
In the theory of disordered systems, the replica approach plays a very special role. On the one hand, it is one of the most powerful theoretical tools and often the only possible approach to study some strongly disordered systems. On the other hand it is difficult to tell in advance whether the predictions of the replica approach are correct or not. When it does not work, one can always try to break the symmetry of replica: this usually makes the calculations much more complicated without being certain that the results become correct. In the replica approach, the calculation usually starts with an integer number $`n`$ of replica. Then, as the limit of physical interest is the limit $`n0`$, one has to extend to non-integer $`n`$ results obtained for integer $`n`$. This is in fact the big difficulty of the replica approach, so it is useful to look at simple examples for which the $`n`$ dependence can be studied in detail.
This is one of the motivations of the present work, where we show how to calculate integer and non-integer moments $`Z^n`$ of the partition function $`Z`$ of a directed polymer in $`1+1`$ dimensions. The geometry we consider is a cylinder infinite in the $`t`$ direction and periodic, of size $`L`$, in the $`x`$ direction (i.e. $`x+Lx`$). The partition function $`Z(x,t)`$ of a directed polymer joining the points $`(0,0)`$ and $`(x,t)`$ on this cylinder is given by the path integral
$`Z(x,t)={\displaystyle _{(0,0)}^{(x,t)}}𝒟y(s)\mathrm{exp}(`$ $`{\displaystyle _0^t}ds[{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{dy(s)}{ds}}\right)^2+`$
$`\eta (y(s),s)]),`$ (1)
where the random medium is characterised by a Gaussian white noise $`\eta (x,t)`$
$$\eta (x,t)\eta (x^{},t^{})=\gamma \delta (xx^{})\delta (tt^{}).$$
(2)
One of the main goals of the present work is to calculate the cumulants $`lim_t\mathrm{}\mathrm{ln}^kZ(t)_c/t`$ of the free energy per unit length of the directed polymer. These cumulants are the coefficients of the small $`n`$ expansion of $`E(n,L,\gamma )`$ defined as
$$E(n,L,\gamma )=\underset{t\mathrm{}}{lim}\frac{1}{t}\mathrm{ln}\left[\frac{Z^n(x,t)}{Z(x,t)^n}\right].$$
(3)
This $`E(n,L,\gamma )`$ was calculated exactly by Kardar for integer $`n`$ and $`L=\mathrm{}`$. His closed expression $`E(n,\mathrm{},\gamma )=n(n^21)\gamma ^2/24`$ cannot however be continued to all values of $`n`$, in particular to negative $`n`$, as it would violate the fact that $`^2E(n,L,\gamma )/n^2`$ is negative. Therefore one does not know the range of validity of this expression.
The second motivation of the present work is to test the universality class of the KPZ equation. The problem (1) of a directed polymer in a random medium is described by the KPZ equation as several other problems such as growing interfaces or exclusion processes. For certain models of this class, the asymmetric exclusion processes, the distribution of the total current $`Y_t`$ integrated over time $`t`$ has been calculated exactly in the long time limit. For large $`t`$, the generating function of this integrated current $`Y_t`$ on a ring of $`L`$ sites takes the form
$$\mathrm{ln}e^{\alpha Y_t}\mathrm{\Lambda }_{\text{max}}(\alpha )t,$$
(4)
and it was shown, when $`L`$ is large and when the parameter $`\alpha `$ in (4) is of order $`L^{3/2}`$, that $`\mathrm{\Lambda }_{\text{max}}(\alpha )`$ takes the following scaling form
$$\mathrm{\Lambda }_{\text{max}}(\alpha )\alpha K_1=K_2G(\alpha K_3)$$
(5)
where $`K_1`$, $`K_2`$ and $`K_3`$ are three constants which depend on the system size $`L`$, the density of particles and the asymmetry.
The interesting aspect of (5) is that the function $`G(\beta )`$ is universal in the sense that it does not depend on any of the microscopic parameters which define the model. It is given (in a parametric form) by
$`\beta `$ $`={\displaystyle \underset{p=1}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{ϵ^p}{p^{3/2}}},`$ (6)
$`G(\beta )`$ $`={\displaystyle \underset{p=1}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{ϵ^p}{p^{5/2}}}.`$ (7)
In the correspondence between the directed polymer problem and the asymmetric exclusion process through the KPZ equation, the role played by $`\mathrm{ln}(Z(t))`$ is the ratio $`Y_t/L`$. Comparing $`\mathrm{exp}(\alpha Y_t)`$ and $`Z^n(t)`$ in equations (3, 4), we see that $`n`$ corresponds to $`\alpha L`$ and $`E(n,L,\gamma )`$ to $`\mathrm{\Lambda }_{\text{max}}(\alpha )`$. If the function $`G(\beta )`$ is characteristic of systems described by the KPZ equation, we expect in the scaling regime (large $`L`$ and $`nL^{1/2}`$), a relation similar to (5) between $`E(n,L,\gamma )`$ (defined by (3)) and $`n`$. This is indeed one of the main results of the present work: when $`L`$ is large and $`nL^{1/2}`$, we find
$$E(n,L,\gamma )=\frac{n\gamma ^2}{24}\frac{\sqrt{\gamma }}{2\sqrt{2\pi }L^{3/2}}G(n\sqrt{2\pi L\gamma }).$$
(8)
It is clear that in order to establish this relation we have to calculate non-integer moments of the partition function.
The paper is organised as follows. In section 2, we recall how the replica approach of (1) can be formulated as a quantum problem with $`n`$ particles on a ring and how this problem can be solved by the Bethe ansatz when the noise is $`\delta `$ correlated as in (2). In section 3, we write an integral equation (26) which, together with some symmetry conditions (27, 28), allows to solve the Bethe equations of section 2. The main advantage of (26) is that the strength $`c`$ of the disorder (where $`c=\gamma L/2`$) and the number of replica appear as continuous parameters. We show how expansions in powers of $`c`$ or in powers of the number $`n`$ of replica can be obtained from this integral equation. In the expansion of the energy $`E(n,L,\gamma )`$ in powers of $`c`$, all the coefficients are polynomials in $`n`$. This allows us to define $`E(n,L,\gamma )`$ for a non-integer $`n`$ at least perturbatively in $`c`$. At the end of section 3, we show how to generate a small $`n`$ expansion which solves the integral equation (26). We also give explicit expressions up to order $`n^3`$ and we notice that in this small $`n`$ expansion of the energy, we have to deal with coefficients that are functions of $`c`$ with a zero radius of convergence. The content of sections 2 and 3 is essentially a recall of a method developed in our previous work. In section 4, we show that the recursion of section 3 which generates all the terms of the small $`n`$ expansion simplifies greatly in the scaling regime ($`c`$ large and $`nc^{1/2}`$) allowing to calculate all the terms of the expansion and to establish (8).
## 2 A quantum system of $`n`$ particles with $`\delta `$ interactions
Let us start with a case slightly more general than (2) where the noise $`\eta (x,t)`$ in (1) is a Gaussian noise $`\delta `$-correlated in time but with some given correlation $`v`$ in space
$$\eta (x,t)\eta (x^{},t^{})=\gamma v(xx^{})\delta (tt^{}).$$
(9)
If we consider the correlation function $`Z(x_1,t)Z(x_2,t)\mathrm{}Z(x_n,t)`$ of the partition function $`Z(x,t)`$ at points $`x_1`$, $`x_2`$,…, $`x_n`$, one can check from (1, 9) that it satisfies
$`{\displaystyle \frac{d}{dt}}Z(x_1,t)Z(x_2,t)\mathrm{}Z(x_n,t)=`$
$`\stackrel{~}{}Z(x_1,t)Z(x_2,t)\mathrm{}Z(x_n,t).`$ (10)
where the Hamiltonian $`\stackrel{~}{}`$ is given by
$$\stackrel{~}{}=\frac{1}{2}\underset{\alpha }{}\frac{^2}{x_\alpha ^2}\gamma \underset{\alpha <\beta }{}v(x_\alpha x_\beta )\gamma \frac{n}{2}v(0),$$
(11)
and where, because of the cylinder geometry in the directed polymer problem, we have $`x_\alpha x_\alpha +L`$ for $`1\alpha n`$.
This implies that in the long time limit,
$$Z(x_1,t)Z(x_2,t)\mathrm{}Z(x_n,t)e^{t\stackrel{~}{E}(n,L,\gamma )},$$
(12)
where $`\stackrel{~}{E}(n,L,\gamma )`$ is the ground state energy of (11).
If one takes the limit $`v(xx^{})\delta (xx^{})`$, the energy $`\stackrel{~}{E}(n,L,\gamma )`$ becomes infinite because of the constant part $`nv(0)/2`$ in (11). This divergence disappears, however, if we consider the ratio $`Z(x_1,t)Z(x_2,t)\mathrm{}Z(x_n,t)/_\alpha Z(x_\alpha ,t)`$, and one can see that in the long time limit,
$$\frac{Z(x_1,t)Z(x_2,t)\mathrm{}Z(x_n,t)}{Z(x_1,t)Z(x_2,t)\mathrm{}Z(x_n,t)}e^{tE(n,L,\gamma )},$$
(13)
where $`E(n,L,\gamma )`$ is the ground state energy of the Hamiltonian
$$=\frac{1}{2}\underset{\alpha }{}\frac{^2}{x_\alpha ^2}\gamma \underset{\alpha <\beta }{}\delta (x_\alpha x_\beta ),$$
(14)
where the positions $`x_\alpha `$ of the $`n`$ particles are on a ring of length $`L`$.
Lieb and Liniger have shown that the Bethe ansatz allows to calculate the ground state energy $`E(n,L,\gamma )`$ of this one dimensional quantum Hamiltonian exactly. The Bethe ansatz consists in looking for a ground state wave function $`\mathrm{\Psi }(x_1,\mathrm{},x_n)`$ of (14) of the form
$$\mathrm{\Psi }(x_1,\mathrm{},x_n)=\underset{P}{}a_Pe^{2(q_1x_{P(1)}+\mathrm{}+q_nx_{P(n)})/L}$$
(15)
in the region $`0x_1\mathrm{}x_nL`$. The sum in (15) runs over all the permutations $`P`$ of $`\{1,\mathrm{},n\}`$ and the value of $`\mathrm{\Psi }`$ in other regions can be deduced from (15) by symmetries. One can show that (15) is the ground state wave function of (14) at energy
$$E(n,L,\gamma )=\frac{2}{L^2}\underset{1\alpha n}{}q_\alpha ^2,$$
(16)
if the $`q_\alpha `$ are the solutions of the $`n`$ coupled equations
$$e^{2q_\alpha }=\underset{\beta \alpha }{}\frac{q_\alpha q_\beta +c}{q_\alpha q_\beta c},$$
(17)
obtained by continuity from the solution $`\{q_\alpha \}=\{0\}`$ at $`c=0`$ where
$$c=\frac{\gamma L}{2}.$$
(18)
Moreover, the $`q_\alpha `$ are all different and the ground state is symmetric ($`\{q_\alpha \}=\{q_\alpha \}`$). (See for instance . Note that $`ik_j`$ and $`c`$ in are here $`\frac{2}{L}q_j`$ and $`\gamma `$; so our $`c`$ defined by (18) and the $`c`$ in are different.)
If we introduce the polynomial $`P(X)`$
$$P(X)=\underset{q_\alpha }{}(Xq_\alpha ),$$
(19)
the system of equations (17) becomes
$$e^{q_\alpha }P(q_\alpha c)+e^{q_\alpha }P(q_\alpha +c)=0,$$
(20)
for any $`1\alpha n`$, and we have from the symmetry of the ground state
$$P(X)=(1)^nP(X).$$
(21)
The knowledge of the polynomial $`P(X)`$ determines the energy (16) as
$$P(X)=X^n\frac{1}{2}\left(\underset{1\alpha n}{}q_\alpha ^2\right)X^{n2}+\mathrm{}$$
(22)
(using (19) and the fact that $`q_\alpha =0`$.)
For small $`c`$, it is possible to solve directly (20) and to determine the $`q_\alpha `$ (see appendix D). This leads to the following expression of the ground state energy (16)
$`E(n,L,\gamma )={\displaystyle \frac{2}{L^2}}n(n1)\left({\displaystyle \frac{c}{2}}+{\displaystyle \frac{c^2}{12}}+{\displaystyle \frac{nc^3}{180}}+O(c^4)\right).`$ (23)
We see that the first coefficients of the small $`c`$ expansion are polynomial in $`n`$. In fact, following the approach of appendix D, one can see that each coefficient of the small $`c`$ expansion of $`E(n,L,\gamma )`$ is polynomial in $`n`$, allowing to define, at least perturbatively in $`c`$, the ground state energy $`E(n,L,\gamma )`$ for non-integer $`n`$. The approach of appendix D becomes however quickly complicated. This is why in the next section we develop a different approach based on the integral equation (26).
## 3 Solution of the Bethe ansatz using an integral equation
In this section we recall the approach developed in our previous work, which consists in writing an integral equation where $`c`$ and $`n`$ appear as continuous parameters and which allows to expand the energy in powers of $`c`$ as well as in powers of $`n`$.
Let us introduce the following function of $`\{q_\alpha \}`$:
$$B(u)=\frac{1}{n}e^{c(u^21)/4}\underset{q_\alpha }{}\rho (q_\alpha )e^{q_\alpha (u1)},$$
(24)
where the parameters $`\rho (q_\alpha )`$ are defined by
$$\rho (q_\alpha )=\underset{q_\beta q_\alpha }{}\frac{q_\alpha q_\beta +c}{q_\alpha q_\beta }.$$
(25)
If the $`\{q_\alpha \}`$ are given by the solution of (17) which corresponds to the ground state, one can show (see appendix A) that the function $`B(u)`$ satisfies the integral equation
$`B(1+u)B(1u)=`$
$`nc{\displaystyle _0^u}𝑑ve^{c(v^2uv)/2}B(1v)B(1+uv).`$ (26)
and the following two conditions
$`B(1)`$ $`=1,`$ (27)
$`B(u)`$ $`=B(u).`$ (28)
Moreover, the energy (16) can be extracted from the knowledge of $`B(u)`$ through
$$E(n,L,\gamma )=\frac{2}{L^2}\left[\frac{n^3c^2}{6}+\frac{nc^2}{12}+\frac{nc}{2}nB^{\prime \prime }(1)\right].$$
(29)
The derivation of (26, 27, 28, 29) is given in appendix A. We are now going to see how one can find perturbatively in $`c`$ or in $`n`$ the solution of (26, 27, 28) and, consequently, the ground state energy (29).
### 3.1 Expansion in powers of $`c`$
To obtain the small $`c`$ expansion of $`B(u)`$ for arbitrary $`n`$, we write
$$B(u)=B_0(u)+cB_1(u)+c^2B_2(u)+\mathrm{}$$
(30)
Conditions (27) and (28) impose that $`B_0(0)=1`$ and all $`B_k(1)=0`$ for $`k>0`$, and that the $`B_k(u)`$ are all even. Moreover, as can be seen directly from (17), the $`q_\alpha `$ scale like $`\sqrt{c}`$ when $`c`$ is small. (Appendix D shows how to obtain the small $`c`$ expansion of the $`q_\alpha `$.) This implies from the definition (24) of $`B(u)`$ that all the $`B_k(u)`$ are polynomials in $`u`$.
At zero-th order in $`c`$, (26) becomes:
$$B_0(1+u)B_0(1u)=0.$$
(31)
The only polynomial solution of (31) consistent with (27, 28), i.e. $`B_0(u)=B_0(u)`$ and $`B_0(1)=1`$ is simply
$$B_0(u)=1$$
(32)
for any $`u`$. We put this back into (26) and we get at first order in $`c`$
$$B_1(1+u)B_1(1u)=nu.$$
(33)
Again, there is a unique polynomial solution which satisfies the facts that $`B_1(u)`$ is even and that $`B_1(1)=0`$:
$$B_1(u)=\frac{n}{4}(u^21).$$
(34)
It is easy to see from (26) that at any order in $`c`$, we have to solve
$`B_k(1+u)B_k(1u)=\varphi _k(u),`$ (35)
where $`\varphi _k(u)`$ is a polynomial odd in $`u`$. There is a unique even polynomial $`B_k(u)`$ solution of (35) satisfying $`B_k(1)=0`$: it is one degree higher than $`\varphi _k(u)`$ and can be determined by equating each power of $`u`$ on both sides of (35). $`[`$ Alternatively, we found a way of writing the solution for any $`\varphi _k(u)`$:
$`B_k(u)=`$ $`[{\displaystyle \frac{s_0}{2}}{\displaystyle _1^u}dv\varphi _k(v)+{\displaystyle \frac{s_1}{2}}(\varphi _k^{}(u)\varphi _k^{}(1))`$
$`+{\displaystyle \frac{s_2}{2}}\left(\varphi _k^{\prime \prime \prime }(u)\varphi _k^{\prime \prime \prime }(1)\right)+\mathrm{}`$
$`+{\displaystyle \frac{s_p}{2}}(\varphi _k^{(2p1)}(u)\varphi _k^{(2p1)}(1))+\mathrm{}]/2`$ (36)
where the $`s_k`$ are the coefficients of the expansion of $`x/\mathrm{sinh}x`$ in powers of $`x`$ (i.e. as $`x/\mathrm{sinh}x=1x^2/6+7x^4/360+\mathrm{}`$, one has $`s_0=1`$, $`s_1=1/6`$, $`s_2=7/360`$, …). $`]`$
This procedure gives for the first terms
$`B(u)=1+{\displaystyle \frac{cn(u^21)}{4}}+{\displaystyle \frac{c^2n(2n+1)(u^21)^2}{96}}`$
$`+{\displaystyle \frac{c^3n(u^21)^2\left(\genfrac{}{}{0pt}{}{5n^2\left(u^21\right)+4n\left(2u^21\right)}{+2\left(u^23\right)}\right)}{5760}}`$
$`+O(c^4).`$ (37)
The energy can then be deduced from (29):
$`E(n,L,\gamma )=2{\displaystyle \frac{n(n1)}{L^2}}[`$ $`{\displaystyle \frac{c}{2}}+{\displaystyle \frac{c^2}{12}}+{\displaystyle \frac{n}{180}}c^3`$ (38)
$`+({\displaystyle \frac{n^2}{1512}}{\displaystyle \frac{n}{1260}})c^4+\mathrm{}].`$
(For (38), we used more terms than given above in $`B(u)`$.) Of course, this expression agrees with (23) obtained directly by expanding the $`q_\alpha `$.
### 3.2 Expansion in powers of $`n`$
The number of particles $`n`$ is *a priori* an integer. However, when we look at the small $`c`$ expansion (3.1) of $`B(u)`$ or (38) of the energy, we see that at any given order in $`c`$ the expression is polynomial in $`n`$. Therefore, one can extend the definition of the small $`c`$ expansion of $`B(u)`$ or of $`E(n,L,\gamma )`$ to non-integer $`n`$. We can also collect in the small $`c`$ expansion of $`B(u)`$ all the terms proportional to $`n`$ and call this series $`b_1(u)`$. From (3.1) we see that
$`b_1(u)=`$ $`{\displaystyle \frac{(u^21)}{4}}c+{\displaystyle \frac{(u^21)^2}{96}}c^2`$
$`+{\displaystyle \frac{(u^21)^2(u^23)}{2880}}c^3+O(c^4).`$ (39)
More generally, we can collect all the terms proportional to $`n^k`$ in the small $`c`$ expansion and call the series $`b_k(u)`$. This means that we can write $`B(u)`$ as a power series in $`n`$
$$B(u)=1+nb_1(u)+n^2b_2(u)+\mathrm{},$$
(40)
where all the $`b_k(u)`$ are defined perturbatively in $`c`$. Conditions (27, 28) impose that all the $`b_k(u)`$ are even and that $`b_k(1)=0`$ for all $`k1`$. We define $`b_0(u)=1`$ for consistency. (It is easy to see in the small $`c`$ expansion that if $`n=0`$, then $`B(u)=1`$.)
We are now going to describe the procedure we used to determine the whole function $`b_1(u)`$ and eventually all the $`b_k(u)`$. If we insert (40) into (26) we get, at first order in $`n`$,
$$b_1(1+u)b_1(1u)=c_0^ue^{c(v^2uv)/2}𝑑v.$$
(41)
It is easy to check that a solution of (41) compatible with the conditions $`b_1(1)=0`$ and $`b_1(u)=b_1(u)`$ is
$$b_1(u)=\sqrt{c}_0^+\mathrm{}𝑑\lambda \frac{\mathrm{cosh}\left(\frac{\lambda u\sqrt{c}}{2}\right)\mathrm{cosh}\left(\frac{\lambda \sqrt{c}}{2}\right)}{\mathrm{sinh}\left(\frac{\lambda \sqrt{c}}{2}\right)}e^{\lambda ^2/2}.$$
(42)
There are however many other solutions of (41), which can be obtained by adding to (42) an arbitrary function $`F(u,c)`$ even and periodic in $`u`$ of period 2 and vanishing at $`u=1`$. If we require that each term in the small $`c`$ expansion of $`b_1(u)`$ is polynomial in $`u`$ (as justified in section 3.1), we see that all the terms of the small $`c`$ expansion of $`F(u,c)`$ must be identically zero. This already shows that (42) has the same small $`c`$ expansion (3.2) as what one would get by collecting all the terms proportional to $`n`$ in the small $`c`$ expansion of section 3.1.
If the solution (42) of (41) had a non-zero radius of convergence in $`c`$, it would be natural to choose this solution and set $`F(u,c)=0`$. However it is easy to see that (42) has a zero radius of convergence in $`c`$: by making the change of variable $`\lambda ^2=2\nu `$, it is easy to see that (42) is the Borel sum of a divergent series.
Apart from being the Borel sum of its expansion in powers of $`c`$, we did not find definitive reasons why (42) is the solution of (41) we should select. However, we can notice that for integer $`n`$, all the $`q_\alpha `$ are real and $`B(u)`$ defined by (24) is analytic in $`u`$ and remains bounded as $`|\text{Im}u|\mathrm{}`$. The solution $`b_1(u)`$ given by (42) is also analytic in $`u`$ and grows as $`\mathrm{ln}(u)`$ as $`|\text{Im}u|\mathrm{}`$. Adding any function $`F(u,c)`$ periodic and analytic in $`u`$ to (42) would produce a much faster growth.
If we insert (40) into (26), we have to solve at order $`n^k`$
$$b_k(1+u)b_k(1u)=\phi _k(u),$$
(43)
where $`\phi _k(u)`$ is some function odd in $`u`$ which can be calculated if we know the previous orders $`b_1(u),\mathrm{},b_{k1}(u)`$.
$`\phi _k(u)=c{\displaystyle \underset{i=0}{\overset{k1}{}}}{\displaystyle _0^u}𝑑v`$ $`e^{c(v^2uv)/2}b_i(1v)\times `$
$`b_{ki1}(1+uv).`$ (44)
We see that the difficulty of selecting a solution of a difference equation appears at all orders in the expansion in powers of $`n`$, and we are now going to explain the procedure we have used to select one solution.
If we write, as $`\phi _k(u)`$ is an odd function of $`u`$,
$$\phi _k(u)=2_0^+\mathrm{}𝑑\lambda \mathrm{sinh}\left(\frac{\lambda u\sqrt{c}}{2}\right)a_k(\lambda ),$$
(45)
which is equivalent, by inverting when $`u`$ is imaginary the Fourier transform in (45), to define $`a_k(\lambda )`$ by
$$a_k(\lambda )=\frac{1}{2i\pi }_0^+\mathrm{}𝑑u\mathrm{sin}\left(\frac{\lambda u}{2}\right)\phi _k\left(\frac{iu}{\sqrt{c}}\right),$$
(46)
then the solution for $`b_k(u)`$ we select is given by
$$b_k(u)=_0^+\mathrm{}𝑑\lambda \frac{\mathrm{cosh}\left(\frac{\lambda u\sqrt{c}}{2}\right)\mathrm{cosh}\left(\frac{\lambda \sqrt{c}}{2}\right)}{\mathrm{sinh}\left(\frac{\lambda \sqrt{c}}{2}\right)}a_k(\lambda ).$$
(47)
Indeed, $`b_k(u)`$ is an even function, vanishes at $`u=1`$ and one can check using (45) that (47) solves (43).
The integrals in (4547) are convergent and equations (3.2, 46, 47) give an automatic way of calculating the $`b_k(u)`$ up to any desired order.
This procedure is the direct generalisation of the choice (42) we did to solve (41). In fact, for $`k=1`$, equations (3.2, 46) give (for $`\lambda 0`$) $`a_1(\lambda )=\sqrt{c}\mathrm{exp}(\lambda ^2/2)`$ and (47) is identical to (42).
As for (42), the solution (47) is not the only solution of (43). At any order $`k`$, we could add an arbitrary even periodic function $`F(u,c)`$ of period 2, the expansion of which vanishes to all order in $`c`$. As for $`b_1(u)`$, we did not find an unquestionable justification of our choice. One can notice nevertheless that (47) is the solution of (43) analytic in $`u`$ and with the slowest growth when $`|\text{Im}u|\mathrm{}`$.
At order $`n^2`$, the procedure (3.2, 46) gives
$`a_2(\lambda )=ce^{\lambda ^2/2}[`$ $`{\displaystyle _0^\lambda }𝑑\mu e^{\mu ^2/2}{\displaystyle \frac{2\mathrm{cosh}\left(\frac{\lambda \mu }{2}\right)2}{\mathrm{tanh}\left(\frac{\mu \sqrt{c}}{2}\right)}}`$ (48)
$`+{\displaystyle _\lambda ^+\mathrm{}}d\mu e^{\mu ^2/2}{\displaystyle \frac{e^{\lambda \mu /2}2}{\mathrm{tanh}\left(\frac{\mu \sqrt{c}}{2}\right)}}],`$
with $`b_2(u)`$ given by (47). Writing down $`b_3(u)`$ or $`a_3(u)`$ would take here about half a column.
We can now give the first terms in the small $`n`$ expansion of the energy. Using relation (29), we find
$`{\displaystyle \frac{L^2}{2}}E(n,L,\gamma )=`$ $`n\left({\displaystyle \frac{c}{2}}+{\displaystyle \frac{c^2}{12}}\right)`$ (49)
$`n^2{\displaystyle \frac{c^{3/2}}{4}}{\displaystyle _0^+\mathrm{}}𝑑\lambda {\displaystyle \frac{\lambda ^2}{\mathrm{tanh}\left(\frac{\lambda \sqrt{c}}{2}\right)}}e^{\lambda ^2/2}`$
$`n^3{\displaystyle \frac{c^2}{4}}{\displaystyle _0^+\mathrm{}}d\lambda {\displaystyle \frac{\lambda ^2}{\mathrm{tanh}\left(\frac{\lambda \sqrt{c}}{2}\right)}}e^{\lambda ^2/2}(`$
$`{\displaystyle _0^\lambda }𝑑\mu e^{\mu ^2/2}{\displaystyle \frac{2\mathrm{cosh}\left(\frac{\lambda \mu }{2}\right)2}{\mathrm{tanh}\left(\frac{\mu \sqrt{c}}{2}\right)}}+`$
$`{\displaystyle _\lambda ^+\mathrm{}}d\mu e^{\mu ^2/2}{\displaystyle \frac{e^{\lambda \mu /2}2}{\mathrm{tanh}\left(\frac{\mu \sqrt{c}}{2}\right)}})+{\displaystyle \frac{n^3c^2}{6}}`$
$`+O(n^4).`$
By making the change of variable $`\lambda ^2=2\nu `$, the terms of order $`n^2`$ and $`n^3`$ appear as Borel transforms of series in $`c`$ with a finite radius of convergence. We conclude that these terms have both a zero radius of convergence in $`c`$.
This small $`n`$ expansion gives quickly very complicated expressions of $`b_k(u)`$. It turns out, as we shall see in the next section, that for large $`c`$, the expressions of the $`b_k(u)`$ get simpler and the energy $`E(n,L,\gamma )`$ can be calculated to all orders in powers of $`n`$.
## 4 Expansion in powers of $`n`$ in the regime $`c\mathrm{}`$
In the previous section, we have developed a procedure allowing to get the small $`n`$ expansion of the energy by solving the problem (2628). Here, we show how this procedure gets greatly simplified for large $`c`$.
The expansion in powers of $`n`$ of the previous section can be summarised as follows: if we use (40) and we write
$$a(\lambda )=na_1(\lambda )+n^2a_2(\lambda )+\mathrm{},$$
(50)
the $`b_k(u)`$ and $`a_k(\lambda )`$ can be obtained by expanding in powers of $`n`$ the following two equations
$$B(u)=1+_0^+\mathrm{}𝑑\lambda \frac{\mathrm{cosh}\left(\frac{\lambda u\sqrt{c}}{2}\right)\mathrm{cosh}\left(\frac{\lambda \sqrt{c}}{2}\right)}{\mathrm{sinh}\left(\frac{\lambda \sqrt{c}}{2}\right)}a(\lambda ),$$
(51)
(this is a rewriting of (47)), and
$`a(\lambda )={\displaystyle \frac{nc}{2i\pi }}{\displaystyle _0^+\mathrm{}}du\mathrm{sin}\left({\displaystyle \frac{\lambda u}{2}}\right)\times `$ (52)
$`{\displaystyle _0^{\frac{iu}{\sqrt{c}}}}𝑑ve^{c(v^2iuv/\sqrt{c})/2}B(1v)B(1+{\displaystyle \frac{iu}{\sqrt{c}}}v).`$
(This is a rewriting of (3.2, 46).) It will be convenient in the following to replace (52) by its Fourier transform
$`2{\displaystyle _0^+\mathrm{}}𝑑\lambda \mathrm{sinh}\left({\displaystyle \frac{\lambda u\sqrt{c}}{2}}\right)a(\lambda )=`$ (53)
$`nc{\displaystyle _0^u}𝑑ve^{c(v^2uv)/2}B(1v)B(1+uv).`$
(This is a rewriting of (3.2, 45).)
We are going to see how one can simplify (5153) when $`c`$ is large. First we observe that for large $`c`$ and $`u`$ fixed of order 1, the expression $`b_1(u)`$ takes the scaling form
$$b_1(1+\frac{u}{\sqrt{c}})\sqrt{c}_0^+\mathrm{}(e^{\lambda u/2}1)e^{\lambda ^2/2}𝑑\lambda .$$
(54)
One can check from (3.2, 46, 47) that this scaling form is present at any order in the small $`n`$ expansion. Indeed, (51) becomes in the large $`c`$ limit
$$B(1+\frac{u}{\sqrt{c}})=1+_0^+\mathrm{}𝑑\lambda (e^{\lambda u/2}1)a(\lambda ),$$
(55)
and using (53), we find
$`2{\displaystyle _0^+\mathrm{}}𝑑\lambda \mathrm{sinh}\left({\displaystyle \frac{\lambda u}{2}}\right)a(\lambda )=`$ (56)
$`n\sqrt{c}{\displaystyle _0^u}𝑑ve^{(v^2uv)/2}B\left(1{\displaystyle \frac{v}{\sqrt{c}}}\right)B\left(1+{\displaystyle \frac{uv}{\sqrt{c}}}\right).`$
It is apparent on (55) and (56) that in the large $`c`$ limit the function $`B(1+u/\sqrt{c})`$ depends only on $`u`$ and $`n\sqrt{c}`$, and $`a(\lambda )`$ depends only on $`\lambda `$ and $`n\sqrt{c}`$. Let us introduce the constant $`K`$
$$K=1_0^+\mathrm{}𝑑\lambda a(\lambda ).$$
(57)
Equation (55) becomes
$$B(1+\frac{u}{\sqrt{c}})=K+_0^+\mathrm{}𝑑\lambda e^{\lambda u/2}a(\lambda ).$$
(58)
In (56), if we write the integral from 0 to $`u`$ as the difference between an integral from 0 to $`+\mathrm{}`$ and an integral from $`u`$ to $`+\mathrm{}`$, and if we change the variable in the second integral to shift it to 0 to $`+\mathrm{}`$, we obtain
$`2{\displaystyle _0^+\mathrm{}}𝑑\lambda \mathrm{sinh}{\displaystyle \frac{\lambda u}{2}}a(\lambda )=`$ (59)
$`n\sqrt{c}{\displaystyle _0^+\mathrm{}}dve^{v^2/2}B(1{\displaystyle \frac{v}{\sqrt{c}}})[e^{uv/2}B(1+{\displaystyle \frac{uv}{\sqrt{c}}})`$
$`e^{uv/2}B(1{\displaystyle \frac{u+v}{\sqrt{c}}})].`$
If we replace $`B(1+(uv)/\sqrt{c})`$ and $`B(1(u+v)/\sqrt{c})`$ by their expression (58), we get after some rearrangements
$`2{\displaystyle _0^+\mathrm{}}𝑑\lambda \mathrm{sinh}{\displaystyle \frac{\lambda u}{2}}a(\lambda )=`$ (60)
$`n\sqrt{c}{\displaystyle _0^+\mathrm{}}dve^{v^2/2}B(1{\displaystyle \frac{v}{\sqrt{c}}})[2K\mathrm{sinh}\left({\displaystyle \frac{uv}{2}}\right)+`$
$`{\displaystyle _0^+\mathrm{}}d\mu a(\mu )e^{\mu v/2}2\mathrm{sinh}\left(u{\displaystyle \frac{v+\mu }{2}}\right)].`$
Taking the Fourier transform of this expression for imaginary $`u`$, we get for $`\lambda 0`$
$`a(\lambda )=`$ $`n\sqrt{c}{\displaystyle _0^+\mathrm{}}dve^{v^2/2}B(1{\displaystyle \frac{v}{\sqrt{c}}})[K\delta (\lambda v)+`$
$`{\displaystyle _0^+\mathrm{}}d\mu a(\mu )e^{\mu v/2}\delta (\lambda v\mu )].`$ (61)
This last expression can be used to calculate $`B(1+u/\sqrt{c})`$ using (58):
$`B(1+{\displaystyle \frac{u}{\sqrt{c}}})=K+`$ (62)
$`n\sqrt{c}{\displaystyle _0^+\mathrm{}}dve^{v^2/2}B(1{\displaystyle \frac{v}{\sqrt{c}}})[Ke^{vu/2}+`$
$`{\displaystyle _0^+\mathrm{}}d\mu a(\mu )e^{\mu v/2}e^{(v+\mu )u/2}].`$
Finally, using (58), we recognise the relation
$`B(1+{\displaystyle \frac{u}{\sqrt{c}}})=K+`$ (63)
$`n\sqrt{c}{\displaystyle _0^+\mathrm{}}𝑑ve^{v^2/2}B(1{\displaystyle \frac{v}{\sqrt{c}}})e^{vu/2}B(1+{\displaystyle \frac{uv}{\sqrt{c}}}).`$
We see that, in the large $`c`$ limit, (51, 52) reduce to this single equation (63). We are now going to see that (63) can be solved to all order in the parameter $`n\sqrt{c}`$. If we introduce the function $`\beta (u)`$ and the parameter $`ϵ`$ defined by
$$\beta (u)=\frac{1}{2K\sqrt{\pi }}e^{u^2/4}B(1+\frac{u}{\sqrt{c}}),$$
(64)
and
$$ϵ=2nK\sqrt{\pi c},$$
(65)
then (63) simply becomes
$$\beta (u)=\frac{1}{2\sqrt{\pi }}e^{u^2/4}+ϵ_0^+\mathrm{}𝑑v\beta (uv)\beta (v).$$
(66)
Using (27, 29, 64), we can express the ground state energy $`E(n,L,\gamma )`$ in terms of $`\beta (u)`$:
$$E(n,L,\gamma )=\frac{2}{L^2}\left[\frac{n^3c^2}{6}+\frac{nc^2}{12}nc\frac{\beta ^{\prime \prime }(0)}{\beta (0)}\right].$$
(67)
It is clear that relation (66) alone determines $`\beta (u)`$, at least perturbatively in $`ϵ`$. So, from (67), we only need to extract $`\beta (0)`$ and $`\beta ^{\prime \prime }(0)`$ from (66).
It is easy to do it for the first orders in $`ϵ`$ directly from equation (66). Moreover, we have found a way of calculating $`\beta (0)`$ and $`\beta ^{\prime \prime }(0)`$, and hence the energy, to all orders in $`ϵ`$. The details of the calculation are given in appendix B. The final result can be written as
$`n\sqrt{c}`$ $`={\displaystyle \frac{1}{2\sqrt{\pi }}}{\displaystyle \underset{k=1}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{ϵ^k}{k^{3/2}}},`$ (68)
$`E(n,L,\gamma )`$ $`={\displaystyle \frac{2}{L^2}}\left[{\displaystyle \frac{nc^2}{12}}+{\displaystyle \frac{\sqrt{c}}{4\sqrt{\pi }}}{\displaystyle \underset{k=1}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{ϵ^k}{k^{5/2}}}\right].`$ (69)
We see that the energy is defined in an implicit way: expression (68) allows to calculate $`ϵ`$ as a function of $`n\sqrt{c}`$, and (69) gives the energy as a function of $`ϵ`$. If we substitute $`c`$ using (18), we obtain the result announced in (8).
For small $`n`$, one can eliminate $`ϵ`$ from (68) and (69). We get
$`{\displaystyle \frac{L^2}{2}}E(n,L,\gamma ){\displaystyle \frac{nc^2}{12}}={\displaystyle \frac{\sqrt{c}}{4\sqrt{\pi }}}(2n\sqrt{c\pi }`$ (70)
$`{\displaystyle \frac{\sqrt{2}}{8}}(2n\sqrt{c\pi })^2+\left({\displaystyle \frac{1}{8}}{\displaystyle \frac{2\sqrt{3}}{27}}\right)(2n\sqrt{c\pi })^3`$
$`+O\left((n\sqrt{c})^4\right)).`$
## 5 Conclusion
In this paper, we have calculated, using the replica method, the first cumulants (13, 49) of the free energy of a directed polymer in a random medium (1) for a cylinder geometry. We used the integral equation (26) of which together with conditions (27, 28) allowed us to expand the moments $`Z^n`$ of the partition function in powers of the strength $`c`$ of the disorder or in powers of the number $`n`$ of replica. All the coefficients of the small $`c`$ expansion (38) are polynomial in $`n`$ allowing to define the expansions for non-integer $`n`$. On the other hand, the coefficients of the expansion (49) in powers of $`n`$ are complicated functions of $`c`$, with in general a zero radius of convergence at $`c=0`$. As already mentioned in , we think that weak disorder expansions of the moments $`Z^n`$ have generically a zero radius of convergence for non-integer $`n`$ when the disorder is Gaussian; this is already the case for a single Ising spin in a Gaussian random field.
To obtain our small $`n`$ expansion, we solved a difference equation (26) which at each order in powers of $`n`$ has several solutions. We selected the particular solution which has the slowest growth in the imaginary $`u`$ direction and has the right small $`c`$ expansion, but we could not exclude other solutions. A different approach, with a direct calculation of the first cumulants of the free energy, and not based on replica would therefore be very useful to test the validity of our expressions (49) which we have been able to derive only perturbatively to all orders in $`c`$.
Although our expansion in powers of $`n`$ becomes quickly very complicated, it simplifies when $`c`$ is large and we could write in this limiting case, all the terms of the small $`n`$ expansion (68, 69). The expression (8) we obtain of the energy $`E(n,L,\gamma )`$ (that is, through (3), the expression of $`Z^n`$) is given exactly by the same scaling function as found for the ASEP. The present work therefore gives additional evidence that the scaling function $`G(\beta )`$ given by (6, 7) is characteristic of the long time behaviour of the KPZ equation in 1+1 dimensions on a ring and that the probability distribution of the free energy for a very long directed polymer on a ring should have an universal shape in the range where the fluctuations per unit length of the free energy are of order $`1/L`$. Other universal distributions for the free energy of a directed polymer have been found recently for different geometries. Our present approach, based on the Bethe ansatz is, at the moment, unable to recover these other distributions. One can try however to extend it to open boundary conditions (in this case too, the Bethe ansatz can be used) instead of periodic boundary conditions and see how this change of boundary conditions affects the distribution of $`\mathrm{ln}Z`$. Of course, it would be very nice to find a simpler approach which would somehow unify all these results and allow to relate all these universal distributions corresponding to the possible geometries, in the spirit of critical phenomena in two dimensions where conformal invariance allows to connect the properties of different geometries.
Technically, the approach followed in the present work is simply to try to find the $`q_\alpha `$ solution of (17) and to calculate the energy (16) which is a symmetric function of the roots $`q_\alpha `$, in such a way that $`n`$ becomes a continuous variable. One could do the same in all kinds of situations. For example, in appendix C, we show how to define and calculate symmetric functions of the roots of Hermite polynomials when the degree of the polynomial becomes non-integer.
Another interesting extension of the present work would be to consider more general correlations of the noise (9). The corresponding quantum problem becomes then the general problem of quantum particles interacting with an arbitrary pair potential. If the interactions are short ranged, one expects the universality class of the KPZ equation to hold, so one could try to repeat our expansion in powers of $`c`$ for a general potential (without the use of the Bethe ansatz) simply by a standard perturbation theory in the strength of the potential. We believe that at any order in the strength of the potential, the ground state energy is polynomial in $`n`$ allowing to define the perturbation expansion for non-integer $`n`$ as we did here. If, with such an approach based on perturbation theory, one could recover the scaling function $`G`$ of (6, 7), one could try to extend the approach to higher dimension as the relation between the directed polymer problem and the quantum hamiltonian is valid in any dimension.
Acknowledgements We thank François David, Michel Gaudin, Vincent Pasquier, Herbert Spohn and André Voros for interesting discussions.
## A Derivation of (26, 27, 28, 29)
Let us first establish some useful properties of the numbers $`\rho (q_\alpha )`$ defined by (25). If the $`q_\alpha `$ are the $`n`$ roots of the polynomial $`P(X)`$
$$P(X)=\underset{q_\alpha }{}(Xq_\alpha ),$$
(A.1)
it is easy to see that the $`\rho (q_\alpha )`$ defined in (25) satisfy
$$\frac{P(X+c)}{P(X)}=1+c\underset{q_\alpha }{}\frac{\rho (q_\alpha )}{Xq_\alpha }.$$
(A.2)
(The two sides have the same poles with the same residues and coincide at $`X\mathrm{}`$.) Expanding the right hand side of (A.2) for large $`X`$, we get
$`{\displaystyle \frac{P(X+c)}{P(X)}}=`$ $`1+c{\displaystyle \underset{q_\alpha }{}}{\displaystyle \frac{\rho (q_\alpha )}{X}}\left(1+{\displaystyle \frac{q_\alpha }{X}}+{\displaystyle \frac{q_\alpha ^2}{X^2}}\right)`$
$`+O\left({\displaystyle \frac{1}{X^4}}\right).`$ (A.3)
On the other hand, using (16, A.1) and the symmetry $`\{q_\alpha \}=\{q_\alpha \}`$ we have
$`P(X)=`$ $`X^n+{\displaystyle \frac{L^2}{4}}E(n,L,\gamma )X^{n2}`$
$`+O(X^{n4}),`$ (A.4)
so that
$`{\displaystyle \frac{P(X+c)}{P(X)}}=1+{\displaystyle \frac{nc}{X}}+{\displaystyle \frac{c^2\left(\genfrac{}{}{0pt}{}{n}{2}\right)}{X^2}}`$ (A.5)
$`+{\displaystyle \frac{c^3\left(\genfrac{}{}{0pt}{}{n}{3}\right)cE(n,L,\gamma )L^2/2}{X^3}}+O\left({\displaystyle \frac{1}{X^4}}\right).`$
Comparing (A.3) and (A.5), we get the relations
$`{\displaystyle \underset{q_\alpha }{}}\rho (q_\alpha )`$ $`=n,`$ (A.6)
$`{\displaystyle \underset{q_\alpha }{}}q_\alpha \rho (q_\alpha )`$ $`=c\left({\displaystyle \genfrac{}{}{0pt}{}{n}{2}}\right),`$ (A.7)
$`{\displaystyle \underset{q_\alpha }{}}q_\alpha ^2\rho (q_\alpha )`$ $`=c^2\left({\displaystyle \genfrac{}{}{0pt}{}{n}{3}}\right){\displaystyle \frac{E(n,L,\gamma )L^2}{2}}.`$ (A.8)
Moreover, by letting $`X=\pm q_\beta c`$ in (A.2) we get for any $`q_\beta `$ root of $`P(X)`$
$$\frac{1}{c}=\underset{q_\alpha }{}\frac{\rho (q_\alpha )}{q_\alpha q_\beta +c}=\underset{q_\alpha }{}\frac{\rho (q_\alpha )}{q_\alpha +q_\beta +c}.$$
(A.9)
Lastly using the symmetry $`\{q_\alpha \}=\{q_\alpha \}`$ and the definition (25), the Bethe ansatz equations (17) reduce to
$$e^{q_\alpha }\rho (q_\alpha )e^{q_\alpha }\rho (q_\alpha )=0.$$
(A.10)
From the definition (24) of $`B(u)`$ and the properties (A.6A.10), it is straightforward to establish (2629): the integral equation (26) is a direct consequence of (24) and (A.9). Properties (27, 28) follow from (24, A.6) and (24, A.10) respectively. Lastly (29) is a consequence of (24, A.6A.8).
## B The energy in the scaling regime
In this appendix, we show how to calculate the energy from the integral equation (66). This equation is of the form
$$\beta (u)=H(u)+ϵ_0^+\mathrm{}𝑑v\beta (uv)\beta (v),$$
(B.1)
where, in our case, $`H(u)`$ is given by
$$H(u)=\frac{1}{2\sqrt{\pi }}e^{\frac{u^2}{4}}.$$
(B.2)
We are going to do our calculations for an arbitrary function $`H(u)`$, even in $`u`$ and decreasing fast enough (to make all the integrals converge) when $`|u|\mathrm{}`$.
To find the energy, we see from (67), that we have to calculate from (B.1) the quantities $`\beta (0)`$ and $`\beta ^{\prime \prime }(0)`$ as functions of $`ϵ`$. We first show that (B.1) is equivalent to
$$\beta (u)=H(u)+ϵ_0^+\mathrm{}𝑑vH(uv)\beta (v),$$
(B.3)
as long as $`H(u)`$ is even and decreases fast enough. Then, we will introduce a new function $`\beta ^{}(u)`$ which is easy to calculate, and relate the derivatives of $`\beta (u)`$ and $`\beta ^{}(u)`$ at $`u=0`$.
### B.1 Equivalence between (B.1) and (B.3)
The solution of (B.3) can be written as
$$\beta (u)=\beta _0(u)+ϵ\beta _1(u)+ϵ^2\beta _2(u)+\mathrm{},$$
(B.4)
where
$`\beta _0(u)`$ $`=H(u),`$ (B.5)
$`\beta _1(u)`$ $`={\displaystyle _0^+\mathrm{}}H(uv_1)H(v_1)𝑑v_1,`$
$`\beta _2(u)`$ $`={\displaystyle _0^+\mathrm{}}H(uv_1)H(v_1v_2)H(v_2)𝑑v_1𝑑v_2,`$
$`\mathrm{}`$
$`\beta _k(u)`$ $`={\displaystyle \mathrm{}_0^+\mathrm{}H(uv_1)H(v_1v_2)\mathrm{}}`$
$`\mathrm{}H(v_{k1}v_k)H(v_k)dv_1\mathrm{}dv_k.`$
For a given $`k>0`$, the integration range of $`\beta _k(u)`$ can be divided into $`k`$ parts: the region where $`v_1`$ has the lowest value of all the $`\{v_i\}`$, the region where $`v_2`$ has the lowest value, …, the region where $`v_k`$ has the lowest value. Let us consider, for some $`j`$ such that $`1jk`$, the region where $`v_j`$ has the lowest value. All the other integrals then run from $`v_j`$ to $`+\mathrm{}`$. If we translate those to integrals running from $`0`$ to $`+\mathrm{}`$ by changing $`v_i`$ into $`v_i+v_j`$, we get:
$`{\displaystyle _0^+\mathrm{}}𝑑v_j{\displaystyle _0^+\mathrm{}}𝑑v_1\mathrm{}𝑑v_{j1}`$ $`H(uv_1v_j)`$ (B.6)
$`H(v_1v_2)\mathrm{}H(v_{j1})\times `$
$`{\displaystyle _0^+\mathrm{}}𝑑v_{j+1}\mathrm{}𝑑v_k`$ $`H(v_{j+1})`$
$`H(v_{j+1}v_{j+2})\mathrm{}H(v_k+v_j)`$
Using the fact that $`H(u)=H(u)`$, we see that (B.6) is equal to
$$_0^+\mathrm{}𝑑v_j\beta _{j1}(uv_j)\beta _{kj}(v_j).$$
(B.7)
By summing over $`j`$, we therefore have
$$\beta _k(u)=_0^+\mathrm{}𝑑v\underset{j=1}{\overset{k}{}}\beta _{j1}(uv)\beta _{kj}(v).$$
(B.8)
Finally, if we multiply by $`ϵ^k`$ and if we sum over $`k`$ all these terms (keeping apart the term for $`k=0`$), we obtain equation (B.1).
The equations (B.1) and (B.3) are thus equivalent and (B.4, B.5) give the solution of (B.1) to any order in $`ϵ`$.
### B.2 Calculation of the derivatives of $`\beta (u)`$
If we look at the expression (B.5) of $`\beta (u)`$ in powers of $`ϵ`$, the calculation of $`\beta (0)`$ and $`\beta ^{\prime \prime }(0)`$ looks simple, especially when $`H(u)`$ is given by (B.2). However, when we try to actually do the calculation, the expressions become quickly complicated with error-functions, primitives of error-functions, etc. It would be much easier if the integrals in (B.5) were running from $`\mathrm{}`$ to $`+\mathrm{}`$ instead of $`0`$ to $`+\mathrm{}`$. This is why we introduce the even function
$$\beta ^{}(u)=\beta _0^{}(u)+ϵ\beta _1^{}(u)+ϵ^2\beta _2^{}(u)+\mathrm{},$$
(B.9)
where, for $`k>0`$,
$$\beta _k^{}(u)=\frac{1}{k+1}\mathrm{}_{\mathrm{}}^+\mathrm{}H(uv_1)\mathrm{}H(v_k)𝑑v_1\mathrm{}𝑑v_k,$$
(B.10)
and $`\beta _0^{}(u)=H(u)`$. One can see easily that
$$\beta ^{}(u)=\frac{1}{2\pi ϵ}_{\mathrm{}}^+\mathrm{}𝑑qe^{iqu}\mathrm{ln}(1ϵ\widehat{H}(q)),$$
(B.11)
where we have defined
$$\widehat{H}(q)=_{\mathrm{}}^+\mathrm{}𝑑ue^{iqu}H(u).$$
(B.12)
The Wiener-Hopf technique allows to relate $`\beta (u)`$ and $`\beta ^{}(u)`$. More specifically, we are going to show that for any $`X>0`$,
$$ϵ_0^+\mathrm{}𝑑ue^{uX}\beta ^{}(u)=\mathrm{ln}\left(1+ϵ_0^+\mathrm{}𝑑ue^{uX}\beta (u)\right).$$
(B.13)
This relation allows to relate the derivatives of $`\beta (u)`$ and $`\beta ^{}(u)`$ at $`u=0`$: indeed, if $`X`$ is large in (B.13), we get
$$_0^+\mathrm{}𝑑ue^{uX}\beta (u)=\frac{\beta (0)}{X}+\frac{\beta ^{}(0)}{X^2}+\frac{\beta ^{^{\prime \prime }}(0)}{X^3}+\mathrm{},$$
(B.14)
and a similar expression for $`\beta ^{}(u)`$. Comparing both sides of (B.13) gives
$`\beta (0)=`$ $`\beta ^{}(0),`$ (B.15)
$`\beta ^{}(0)=`$ $`{\displaystyle \frac{ϵ}{2}}\beta (0)^2,`$
$`\beta ^{\prime \prime }(0)=`$ $`\beta _{}^{}{}_{}{}^{\prime \prime }(0)+{\displaystyle \frac{ϵ^2}{6}}\beta (0)^3.`$
(We have used the fact that $`\beta _{}^{}{}_{}{}^{}(0)=0`$ because $`\beta ^{}(u)`$ is an even function.)
In order to prove (B.13), the first thing to note is that, as $`H(u)`$ decreases fast when $`u\pm \mathrm{}`$, then also does $`\beta (u)`$. This allows to define the two “partial” Fourier transforms
$`\widehat{\beta }_+(q)`$ $`={\displaystyle _0^+\mathrm{}}𝑑ue^{iqu}\beta (u),`$ (B.16)
$`\widehat{\beta }_{}(q)`$ $`={\displaystyle _{\mathrm{}}^0}𝑑ue^{iqu}\beta (u).`$ (B.17)
It is easy to see that $`\widehat{\beta }_+(q)`$ is analytic in the upper half-plane ($`\text{Im}q0`$). Moreover, in this half-plane, $`\widehat{\beta }_+(q)`$ is bounded and vanishes when $`|q|\mathrm{}`$. Conversely, $`\widehat{\beta }_{}(q)`$ is analytic, bounded and decreases to 0 at infinity when $`\text{Im}q0`$.
The function $`\beta (u)`$ can be written in terms of $`\widehat{\beta }_+(q)`$ and $`\widehat{\beta }_{}(q)`$:
$$\beta (u)=\frac{1}{2\pi }_{\mathrm{}}^+\mathrm{}𝑑qe^{iqu}(\widehat{\beta }_+(q)+\widehat{\beta }_{}(q)),$$
(B.18)
which allows to express the right-hand side of (B.13) when $`X`$ is positive:
$`{\displaystyle _0^+\mathrm{}}𝑑ue^{uX}\beta (u)=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑q{\displaystyle \frac{\widehat{\beta }_+(q)}{X+iq}}`$ (B.19)
$`+{\displaystyle \frac{1}{2\pi }}{\displaystyle _{\mathrm{}}^+\mathrm{}}𝑑q{\displaystyle \frac{\widehat{\beta }_{}(q)}{X+iq}}.`$
We calculate the two integrals in the right hand side of (B.19) by the residue theorem. As $`\widehat{\beta }_+(q)`$ is analytic and decreases at infinity in the upper half-plane, the first integral can be written when $`X>0`$ as a contour integral around the upper half-plane. The only contribution to the first integral comes, using Cauchy’s theorem, from the pole $`q=iX`$. One can also check that the second integral vanishes (using a contour around the lower half-plane and the fact that $`\widehat{\beta }_{}(q)`$ has no pole). Therefore, (B.19) gives
$$_0^+\mathrm{}𝑑ue^{uX}\beta (u)=\widehat{\beta }_+(iX).$$
(B.20)
Now, if we multiply (B.3) by $`\mathrm{exp}(iqu)`$ and if we integrate over $`u`$, we easily get for any real $`q`$
$$\widehat{\beta }_+(q)+\widehat{\beta }_{}(q)=\widehat{H}(q)+ϵ\widehat{H}(q)\widehat{\beta }_+(q).$$
(B.21)
This relation between $`\widehat{H}(q)`$, $`\widehat{\beta }_{}(q)`$ and $`\widehat{\beta }_+(q)`$, together with (B.11) gives
$`\beta ^{}(u)={\displaystyle \frac{1}{2\pi ϵ}}{\displaystyle _{\mathrm{}}^+\mathrm{}}dqe^{iqu}(`$ $`\mathrm{ln}(1+ϵ\widehat{\beta }_+(q))`$ (B.22)
$`\mathrm{ln}(1ϵ\widehat{\beta }_{}(q)).`$
Using again that, in the upper half-plane, $`\widehat{\beta }_+(u)`$ is analytic and vanishes at infinity, we see that, for a *small enough* $`ϵ`$, the quantity $`\mathrm{ln}(1+ϵ\widehat{\beta }_+(q))`$ is also analytic and decreases to 0 at infinity when $`\text{Im}q0`$. Similarly, $`\mathrm{ln}(1ϵ\widehat{\beta }_{}(q))`$ has the same properties for $`\text{Im}q0`$. This allows to calculate the left hand side of (B.13) as we did for the right hand side. We find
$$_0^+\mathrm{}𝑑ue^{uX}\beta ^{}(u)=\frac{1}{ϵ}\mathrm{ln}(1+ϵ\widehat{\beta }_+(iX)).$$
(B.23)
Comparing (B.20) and (B.23) completes the proof of (B.13).
We can now give an expression of the energy. If we use the definition (B.2) of $`H(u)`$ in (B.11, B.12), we find
$$\beta ^{}(u)=\frac{1}{2\sqrt{\pi }}\underset{k=0}{\overset{+\mathrm{}}{}}\frac{ϵ^k}{(k+1)^{3/2}}e^{u^2/[4(k+1)]}.$$
(B.24)
This gives
$`\beta ^{}(0)=`$ $`{\displaystyle \frac{1}{2\sqrt{\pi }}}{\displaystyle \underset{k=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{ϵ^k}{(k+1)^{3/2}}},`$ (B.25)
$`\beta _{}^{}{}_{}{}^{\prime \prime }(0)=`$ $`{\displaystyle \frac{1}{4\sqrt{\pi }}}{\displaystyle \underset{k=0}{\overset{+\mathrm{}}{}}}{\displaystyle \frac{ϵ^k}{(k+1)^{5/2}}},`$ (B.26)
and, together with (B.15), these equations allows to give an expression of $`\beta (0)`$ and $`\beta ^{\prime \prime }(0)`$.
From (27, 64, 65), we see that
$$ϵ\beta (0)=n\sqrt{c}.$$
(B.27)
Then, using (B.15) we get
$`ϵ\beta ^{}(0)=`$ $`n\sqrt{c},`$ (B.28)
$`{\displaystyle \frac{\beta ^{\prime \prime }(0)}{\beta (0)}}=`$ $`{\displaystyle \frac{ϵ}{n\sqrt{c}}}\beta _{}^{}{}_{}{}^{\prime \prime }(0)+{\displaystyle \frac{n^2c}{6}}.`$
The energy is given by (67). We get:
$$E(n,L,\gamma )=\frac{2}{L^2}\left[\frac{nc^2}{12}ϵ\sqrt{c}\beta _{}^{}{}_{}{}^{\prime \prime }(0)\right].$$
(B.29)
And, finally, using relation (B.25, B.26), we obtain (68, 69).
## C Hermite polynomials with a non-integer number of roots
What we try to do in this whole paper is essentially to calculate $`_\alpha q_\alpha ^2`$ (the energy) where $`\{q_\alpha \}`$ is solution of (17), in such a way that $`n`$ appears as a continuous parameter. This allows us to obtain expressions of the energy for non-integer $`n`$.
One can use the same procedure in other kinds of situations. A simple example which illustrates our calculations is the case of the zeroes of Hermite polynomials.
The $`n`$-th Hermite polynomial $`H_n(X)`$ is the solution polynomial in $`X`$ with leading coefficient 1 of the differential equation
$$\frac{1}{2}H_n^{\prime \prime }(X)XH_n^{}(X)+nH_n(X)=0.$$
(C.1)
The polynomial $`H_n(x)`$ is of degree $`n`$ and has the symmetry $`H_n(X)=()^nH_n(X)`$. For example, we have $`H_4(X)=X^43X^2+\frac{3}{4}`$. The $`n`$ roots $`\{h_\alpha \}`$ ($`1\alpha n`$) of $`H(X)`$ are real and distinct.
By deriving (C.1) $`p`$ times with respect to $`X`$, we see that, for all $`p`$,
$$XH_n^{(p+1)}(X)=\frac{1}{2}H_n^{(p+2)}(X)+(np)H_n^{(p)}(X).$$
(C.2)
This shows that the $`(np)`$-th Hermite polynomial is, up to a constant factor, equal to the $`p`$-th derivative of $`H_n(X)`$. (This property will be used a lot in appendix D).
Equation (C.1) can be used directly to calculate the first coefficients of $`H_n(X)`$
$$H_n(X)=X^n\frac{1}{2}\left(\genfrac{}{}{0pt}{}{n}{2}\right)X^{n2}+\frac{3}{4}\left(\genfrac{}{}{0pt}{}{n}{4}\right)X^{n4}+\mathrm{}$$
(C.3)
Using (C.3), the symmetry of $`H(X)`$ and the large $`X`$ expansion
$$\frac{H_n^{}(X)}{H_n(X)}=\underset{p0}{}\frac{1}{X^{p+1}}\left(\underset{\alpha }{}h_\alpha ^p\right),$$
(C.4)
we can calculate the moments of the roots $`\{h_\alpha \}`$ of $`H(X)`$:
$`{\displaystyle \underset{\alpha }{}}h_\alpha ^2`$ $`={\displaystyle \frac{n(n1)}{2}},`$ (C.5)
$`{\displaystyle \underset{\alpha }{}}h_\alpha ^4`$ $`={\displaystyle \frac{n(n1)}{4}}(2n3),`$ (C.6)
and so on. These moments are a priori defined only for integer $`n`$ but as the expressions are polynomial in $`n`$, one can obviously extend their definition to non-integer $`n`$ (similarly to what we do in the small $`c`$ expansion of $`B(u)`$ in section (3.2)).
To generate all the moments of the roots $`h_\alpha `$, it is convenient to consider the generating function
$$Q(u)=\underset{h_\alpha }{}e^{h_\alpha u},$$
(C.7)
which is reminiscent of the quantity $`\beta (u)`$ defined in our quantum problem. (Using (24) and (64) we can check that $`\beta (u)\mathrm{exp}(u\sqrt{c}/2)\rho (q_\alpha )\mathrm{exp}(q_\alpha u/\sqrt{c})`$.)
The function $`Q(u)`$ is hard to calculate for general $`n`$ but we can expand it in powers of $`n`$. This can be done by considering
$$\mathrm{\Psi }(X)=\frac{H_n^{}(X)}{H_n(X)}=_0^+\mathrm{}𝑑uQ(u)e^{uX},$$
(C.8)
which is defined only for $`X`$ positive and large enough to make the integral converges. This function $`\mathrm{\Psi }(X)`$ is solution of a differential equation which follows from (C.1):
$$\frac{1}{2}\mathrm{\Psi }^{}(X)+\frac{1}{2}\mathrm{\Psi }(X)^2X\mathrm{\Psi }(X)+n=0.$$
(C.9)
To obtain an expansion in powers of $`n`$, we write
$$\mathrm{\Psi }(X)=n\mathrm{\Psi }_1(X)+n^2\mathrm{\Psi }_2(X)+\mathrm{}$$
(C.10)
Thus $`\mathrm{\Psi }_1(X)`$ satisfies
$$\frac{1}{2}\mathrm{\Psi }_1^{}(X)X\mathrm{\Psi }_1(X)+1=0.$$
(C.11)
This differential equation can easily be solved, and the integration constant can be fixed using the requirement (C.8) that, for large $`X`$, $`\mathrm{\Psi }(X)n/X`$
$$\mathrm{\Psi }_1(X)=_0^+\mathrm{}𝑑ue^{uXu^2/4}.$$
(C.12)
Then order $`n^2`$ of (C.9) gives
$$\frac{1}{2}\mathrm{\Psi }_2^{}(X)X\mathrm{\Psi }_2(X)+\frac{1}{2}\mathrm{\Psi }_1(X)^2=0,$$
(C.13)
the solution of which can be written as
$$\mathrm{\Psi }_2(X)=2_0^+\mathrm{}𝑑ue^{uXu^2/4}_0^+\mathrm{}𝑑t\frac{\mathrm{cosh}\left(\frac{ut}{\sqrt{2}}\right)1}{t}e^{t^2}.$$
(C.14)
The procedure can be iterated to any order in $`n`$ (of course expressions become more and more complicated). Using (C.8) and the expressions of $`\mathrm{\Psi }_1(X)`$ and $`\mathrm{\Psi }_2(X)`$ we can give an expression of $`Q(u)`$:
$`Q(u)=`$ $`ne^{u^2/4}+`$ (C.15)
$`2n^2e^{u^2/4}{\displaystyle _0^+\mathrm{}}𝑑t{\displaystyle \frac{\mathrm{cosh}\left(\frac{ut}{\sqrt{2}}\right)1}{t}}e^{t^2}+O(n^3).`$
Expanding this expression in powers of $`u`$, one calculate from this expression and from (C.7) the terms linear and quadratic in $`n`$ of all the moments of the $`h_\alpha `$. (The results agree for the second and the fourth moments with (C.5, C.6).)
We noticed that for small $`n`$, the expression (C.15) corresponds to $`n`$ roots $`h_\alpha `$ distributed along the imaginary axis with a Gaussian distribution. We do not know whether this is general and whether there exists, for general non-integer $`n`$, a distribution of the roots $`h_\alpha `$ in the complex plane plane which gives all moments calculated as in (C.5, C.6).
It is interesting to notice the similarity between $`Q(u)`$ and $`\beta (u)`$ defined in section 4.
## D The expansion in powers of $`c`$ using Hermite polynomials
In this appendix we show how to expand the solution $`\{q_\alpha \}`$ of (17) in powers of $`c`$ for integer $`n`$. One can see from (17) that the roots $`q_\alpha `$ scale for small $`c`$ like $`\sqrt{c}`$. It is thus convenient to rescale the polynomial $`P(X)`$ defined in (19) and the $`q_\alpha `$ in the following way:
$`q_\alpha `$ $`=r_\alpha \sqrt{c},`$ (D.1)
$`P(X\sqrt{c})`$ $`=c^{n/2}R(X).`$
($`\{r_\alpha \}`$ are thus the roots of $`R(X)`$.) With these new variables, equation (20) becomes
$$e^{r_\alpha \sqrt{c}}R(r_\alpha \sqrt{c})+e^{r_\alpha \sqrt{c}}R(r_\alpha +\sqrt{c})=0.$$
(D.2)
As the roots $`r_\alpha `$ of $`R(X)`$ are all distinct, this equation is equivalent to
$`e^{X\sqrt{c}}R(X\sqrt{c})+e^{X\sqrt{c}}R(X+\sqrt{c})=`$ (D.3)
$`2(\mathrm{cosh}X\sqrt{c}+f(X))R(X),`$
where $`f(X)`$ is *analytic* (this follows from the fact that as $`R(X)`$ is polynomial, $`f(X)`$ defined by (D.3) is obviously meromorphic; moreover as the left hand side of (D.3) vanishes at all the roots of $`R(X)`$, $`f(X)`$ has no pole.) We are now going to solve (D.3) as a power series in $`c`$ (i.e. find both $`f(X)`$ and $`R(X)`$ as power series in $`c`$).
### D.1 Expansion of the polynomial $`R(X)`$
We only have the single equation (D.3) to obtain two quantities ($`R(X)`$ and $`f(X)`$); however, using the fact that $`f(X)`$ has no pole and $`R(X)`$ is a polynomial, both quantities can be determined in a small $`c`$ expansion. Let us write
$`R(X)`$ $`=R_0(X)+cR_1(X)+c^2R_2(X)+\mathrm{},`$ (D.4)
$`f(X)`$ $`=cf_1(X)+c^2f_2(X)+\mathrm{},`$
where the $`f_i(X)`$ have no pole, $`R_0(X)`$ is a polynomial of degree $`n`$ (the term of highest degree in $`R_0(X)`$ is $`X^n`$) and all the $`R_i(X)`$ (for $`i1`$) are polynomials of degree less than $`n`$. At first order in $`c`$, we find that (D.3) gives:
$$\frac{1}{2}R_0^{\prime \prime }XR_0^{}=f_1R_0.$$
(D.5)
As $`f_1(X)`$ has no pole, it must be a polynomial. Because $`R_0(X)`$ is of degree $`n`$, we see by looking at both sides of (D.5) that, necessarily, $`f_1(X)=n`$. We recognise then the differential equation (C.1) that defines Hermite polynomials. Therefore
$`f_1(X)=`$ $`n,`$ (D.6)
$`R_0(X)=`$ $`H(X).`$
We recover that way that the $`r_\alpha `$ are the zeroes of the $`n`$-th Hermite polynomial when $`c`$ is very small.
At next order in $`c`$, equation (D.3) gives
$`{\displaystyle \frac{1}{2}}R_1^{\prime \prime }XR_1^{}+nR_1f_2H=`$ (D.7)
$`{\displaystyle \frac{X^3}{6}}H^{}{\displaystyle \frac{X^2}{4}}H^{\prime \prime }+{\displaystyle \frac{X}{6}}H^{(3)}{\displaystyle \frac{1}{24}}H^{(4)}.`$
As $`R_1`$ and $`H`$ are polynomials, (D.7) tells us that $`f_2H`$ is a polynomial too. We also know that $`f_2(X)`$ has no pole, thus it must be a polynomial. $`R_1(X)`$ is of degree strictly less than $`n`$, so the expression $`R_1^{\prime \prime }/2XR_1+nR_1`$ is of degree strictly less than $`n`$. As $`H`$ is of degree $`n`$, we recognise in (D.7) an euclidian division of polynomials: $`f_2(X)`$ is the quotient of the right hand side of equation (D.7) divided by $`H(X)`$, and the terms involving $`R_1(X)`$ form the remainder of this division. This ensures that there is only one possible function $`f_2(X)`$ which verifies (D.7).
In practise, to perform this euclidian division we can use the property (C.2) of the Hermite polynomials as many times as needed in the right hand side of (D.7): for instance, we transform the term $`X^3H^{}/6`$ into $`nX^2H/6+X^2H^{\prime \prime }/12`$. We cannot change $`X^2H`$ anymore, but we can apply (C.2) to the term $`X^2H^{\prime \prime }`$. When no more transformation is possible, we are left with:
$`{\displaystyle \frac{X^3}{6}}H^{}{\displaystyle \frac{X^2}{4}}H^{\prime \prime }+{\displaystyle \frac{X}{6}}H^{(3)}{\displaystyle \frac{1}{24}}H^{(4)}=`$ (D.8)
$`\left({\displaystyle \frac{n}{6}}X^2{\displaystyle \frac{n(n1)}{6}}\right)H{\displaystyle \frac{1}{12}}H^{\prime \prime }.`$
The Euclidian division is then easy to perform
$`f_2(X)={\displaystyle \frac{n}{6}}X^2+{\displaystyle \frac{n(n1)}{6}},`$ (D.9)
$`{\displaystyle \frac{1}{2}}R_1^{\prime \prime }XR_1^{}+nR_1={\displaystyle \frac{1}{12}}H^{\prime \prime }.`$
Using again (C.2), the differential equation on $`R_1`$ can be solved; we find
$$R_1(X)=\frac{1}{24}H^{\prime \prime }(X).$$
(D.10)
As $`R_1(X)`$ is simply a derivative of $`H(X)`$, and as $`f_1(X)`$ is a known polynomial of $`X`$, we see that at the next order in $`c`$ we will have to solve an equation of the form
$$\frac{1}{2}R_2^{\prime \prime }XR_2+nR_2f_3H=X^jH^{(k)}.$$
(D.11)
Using many times equation (C.2) the right hand side can be written in a “canonical form”:
$$X^jH^{(k)}=X^jH+H^{(k)},$$
(D.12)
which allows to write $`f_3`$ as a polynomial in $`X`$ and $`R_2`$ as a sum of derivatives of $`H(X)`$. It is easy to see recursively that at any order $`c^k`$ in the expansion we can repeat this procedure to calculate $`f_k(X)`$ and $`R_{k1}(X)`$. As a result we see that $`f_k`$ is a polynomial in $`X`$ and that $`R_{k1}`$ can be written as a sum of derivatives of $`H(X)`$.
It is worth noting that at each order the variable $`n`$ comes from the previous orders and from transformations of the kind $`XH^{}(X)\frac{1}{2}H^{\prime \prime }(X)+nH(X)`$. Because those are the two only mechanisms by which $`n`$ appears, it is easy to see that at each order the coefficients of the sum of derivatives of $`H(X)`$ that constitutes $`R_{k1}(X)`$ are all *polynomials in $`n`$*.
A computer can easily do this tedious but straightforward task to any desired order. Up to $`c^3`$, we find:
$`R=`$ $`H{\displaystyle \frac{c}{24}}H^{\prime \prime }c^2\left({\displaystyle \frac{n}{360}}H^{\prime \prime }{\displaystyle \frac{7}{5760}}H^{(4)}\right)+`$
$`c^3(({\displaystyle \frac{n}{2520}}{\displaystyle \frac{n^2}{3024}})H^{\prime \prime }+{\displaystyle \frac{11n}{60480}}H^{(4)}`$
$`{\displaystyle \frac{31}{967680}}H^{(6)})+O(c^4).`$ (D.13)
### D.2 Expansion of the roots $`r_\alpha `$ of $`R(X)`$
As seen in (D.1), the polynomial $`R(X)`$ is to leading order in $`c`$ given by $`H(X)`$. It is thus natural to write the roots $`r_\alpha `$ of $`R(X)`$ as
$$r_\alpha =h_\alpha +cx_\alpha +O(c^2).$$
(D.14)
($`\{h_\alpha \}`$ are the roots of $`H`$). Inserting (D.14) into (D.1), we find, at first order in $`c`$,
$$x_\alpha H^{}(h_\alpha )\frac{1}{24}H^{\prime \prime }(h_\alpha )=0.$$
(D.15)
Using the definition (C.1) of Hermite polynomials, we have $`H^{\prime \prime }(h_\alpha )=2h_\alpha H^{}(h_\alpha )`$. This gives in turn $`x_\alpha =\frac{1}{12}h_\alpha `$. Repeating this procedure to any order in $`c`$, we generate terms of the form $`h_\alpha ^jH^{(k)}(h_\alpha )`$ which can be reduced to terms of the form $`h_\alpha ^lH^{}(h_\alpha )`$ by using (C.2) as many times as necessary. It is then possible to divide the expression by $`H^{}(h_\alpha )`$ and we are left with an equation giving each new term in the expansion of $`r_\alpha `$ as a *polynomial in $`h_\alpha `$*. Again, this can be programmed, and we get, up to the order $`c^2`$:
$`r_\alpha ={\displaystyle \frac{q_\alpha }{\sqrt{c}}}=h_\alpha +{\displaystyle \frac{c}{12}}h_\alpha +c^2(`$ $`\left({\displaystyle \frac{n}{120}}{\displaystyle \frac{11}{1440}}\right)h_\alpha `$
$`{\displaystyle \frac{1}{360}}h_\alpha ^3)+O(c^3).`$ (D.16)
Using (D.1) and (16), this leads to
$`{\displaystyle \frac{2}{L^2}}E(n,L,\gamma )=c{\displaystyle h_\alpha ^2}{\displaystyle \frac{c^2}{6}}{\displaystyle h_\alpha ^2}`$ (D.17)
$`{\displaystyle \frac{c^3}{360}}\left((6n3){\displaystyle h_\alpha ^2}2{\displaystyle h_\alpha ^4}\right)+O(c^4).`$
which coincides with (38) when one uses the properties (C.5, C.6) of the roots $`h_\alpha `$ of the Hermite polynomials.
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# Contents
## 1 Introduction
In its original form Kaluza-Klein reduction was used for the purpose of deriving a four-dimensional theory comprising gravity, a $`U(1)`$ gauge field and a dilatonic scalar, starting from pure gravity in five dimensions. The extra dimension is taken to be a circle, and the five-dimensional metric is then assumed to be independent of the coordinate $`y`$ on the circle. Such a truncation is consistent, and gives rise to an Einstein-Maxwell theory in $`D=4`$, coupled to the dilatonic scalar field. The consistency of the truncation is assured because the reduction ansatz retains all the four-dimensional fields that are independent of $`y`$, while setting all fields that would be associated with $`y`$-dependent harmonics on $`S^1`$ to zero. In a similar vein, Kaluza-Klein reductions involving higher-dimensional theories compactified on tori can also be considered, and again consistent truncations where all fields are taken to be independent of the torus coordinates can be performed.
The situation is much less clear-cut in the case where one performs a reduction on a curved internal manifold, such as a sphere. The new complication in such a case is that the harmonics on the internal space associated with the massless fields in the lower dimension typically now depend on the coordinates of the internal space. This causes no difficulty in a linearised analysis of small fluctuations around a ground-state solution (see, for example, , and references therein), but as soon as one wants to consider the full non-linear structure of the theory it raises the possibility of inconsistencies in a truncation to the massless sector. In fact this is more than a possibility; in general, there will definitely be inconsistencies. This makes it all the more remarkable that there exist certain exceptional cases in which a fully non-linear sphere reduction and truncation is completely and rigorously consistent. Many of the known cases involve special reductions of supergravity theories, notably involving $`S^7`$ or $`S^4`$ reductions of $`D=11`$, the $`S^5`$ reduction of type IIB,<sup>1</sup><sup>1</sup>1The consistency of the $`S^5`$ reduction to five-dimensional maximal gauged supergravity remains conjectural at this time, but strong supporting evidence has been obtained, including various explicit consistent reductions to subsets of the maximal supergravity , and an explicit expression for the complete metric reduction ansatz . and a local $`S^4`$ reduction of the massive type IIA theory . Other exceptional examples of consistent sphere reductions in which all the Yang-Mills gauge fields can be retained have also been found recently, for cases that do not necessarily have any connection with supersymmetry. These comprise the reduction of the low-energy limit of the bosonic string, in an arbitrary dimension $`D`$, on the 3-sphere or the $`(D3)`$-sphere, and the reduction of certain theories of gravity plus a dilaton and a 2-form field strength in $`D`$ dimensions on a 2-sphere . In all these cases, there is no known group-theoretic proof for why the reduction should be consistent.<sup>2</sup><sup>2</sup>2A group-theoretic argument has been used in in order to prove that an $`n`$-sphere reduction of a theory of gravity plus dilaton plus $`n`$-form field strength that retained all the $`SO(n+1)`$ Yang-Mills fields in a massless truncation could not be consistent except in the exceptional cases listed above.
Two approaches to proving the consistency of these supergravity sphere reductions have been pursued in the literature. For the $`S^7`$ and $`S^4`$ reductions from $`D=11`$, the truncations to the maximally supersymmetric gauged $`SO(8)`$ and $`SO(5)`$ supergravities in $`D=4`$ and $`D=7`$ have been argued to be consistent by demonstrating that consistent supersymmetry transformation rules in the lower dimension can be extracted from the original ones in $`D=11`$. A complete and rigorous proof of consistency along these lines would in principle require the analysis of the supersymmetry transformation rules to all orders, including quartic fermion terms, and the difficulties in doing this are considerable. However, it seems reasonable to conclude that the already highly non-trivial success at the quadratic level would persist to all orders. The approach has been used for the $`S^7`$ and $`S^4`$ reductions, and in the latter case has allowed an explicit construction of the exact bosonic reduction ansatz. No analogous complete results have been obtained for the $`S^5`$ reduction of type IIB supergravity, but it seems highly likely to be consistent also.
The alternative approach to proving the consistency of a Kaluza-Klein reduction is a more direct one, in which one explicitly constructs a reduction ansatz which, when substituted into the full set of higher-dimensional equations of motion, gives a consistent embedding provided that the lower-dimensional equations of motion are satisfied. This approach has been used to provide a complete proof of the consistency in several sphere reductions, where further truncations to subsets of the fields of the maximal massless supermultiplet are made. Cases that have been fully proven by this means include $`N=2`$ gauged $`SU(2)`$ supergravity in $`D=7`$ by an $`S^4`$ reduction from $`D=11`$ ; the $`N=4`$ gauged $`SU(2)\times U(1)`$ supergravity in $`D=5`$ by an $`S^5`$ reduction from type IIB in $`D=10`$ ; the $`N=4`$ gauged $`SO(4)`$ supergravity by $`S^7`$ reduction from $`D=11`$ , and the $`N=2`$ gauged $`SU(2)`$ supergravity in $`D=6`$ by a local $`S^4`$ reduction from the massive type IIA theory in $`D=10`$ . (In this last example $`N=2`$ is in fact the largest supersymmetry for gauged supergravity in $`D=6`$, even though ungauged $`N=4`$ supergravity exists.) In addition, the consistency of the truncations of the $`S^4`$, $`S^5`$ and $`S^7`$ reductions to include gravity and all the diagonal scalars of the $`SL(5,R)/SO(5)`$, $`SL(6,R)/SO(6)`$ and $`SL(8,R)/SO(8)`$ submanifolds of the full scalar cosets of the maximal supergravities have been fully demonstrated .
It is significant that all these examples involve reductions on spheres. At the linearised level there is no reason why one should not consider also reductions on internal spaces of other topologies. Examples that have been considered in the past include the Einstein spaces contructed as $`U(1)`$ bundles over $`CP^2\times S^2`$ and $`S^2timesS^2\times S^2`$, as compactifications of eleven-dimensionsal supergravity, and $`U(1)`$ bundles over $`S^2\times S^2`$, as compactifications of type IIB supergravity. The first two examples were first discussed in detail in , where they were constructed as the coset spaces $`M^{pqr}=SU(3)\times SU(2)\times U(1)/(SU(2)\times U(1)\times U(1)`$ and $`Q^{pqr}=SU(2)^3/(U(1)\times U(1))`$ respectively, and the linearised massless spectra were obtained. They were subsequently reconstructed from the viewpoint of $`U(1)`$ bundles over $`CP^2\times S^2`$ and $`S^2\times S^2\times S^2`$ respectively, in , where a stability analysis was also given. The complete massive spectrum was obtained in a linearised analysis in . The five-dimensional example, the $`U(1)`$ bundle over $`S^2\times S^2`$, was discussed from the AdS/CFT viewpoint in , and in a field-theoretic context in . The full Klauza-Klein spectrum was obtained in , and its matching with the conformal operators of the dual CFT was obtained.
Amongst the lower-dimensional massless fields that would result from reductions such as these will be Yang-Mills gauge bosons with gauge group given by the isometry group of the internal space. One may wonder whether a consistent truncation that includes the Yang-Mills gauge fields is possible in these more general reductions too, or whether it is a special feature of the spherical spaces that ensures the consistency.
Some results on certain of these more general reductions were obtained in previous studies. In this paper, we shall address the question in a slightly broader context. The conclusions will be similar to those reached in the previous cases, namely that in general the reductions on non-spherical internal spaces do not allow consistent truncations to the massless sector, even in those exceptional and remarkable theories where consistent sphere reductions are possible. In fact it is much easier to demonstrate the inconsistency of an inconsistent truncation than to prove the consistency of a consistent one. As was discussed in , when there are inconsistencies they tend to show up in relatively easily-studied sectors of the theory, at the level of cubic interaction terms in the Lagrangian. In this paper we shall be considering a specific type of cubic interaction, namely terms that are bilinear in the lower-dimensional gauge fields, and that couple to lower-dimensional linearised spin-2 fields. This sector provides a necessary condition for consistency of a truncation; in general it turns out that the bilinears in gauge fields can act as sources for massive as well as massless spin-2 excitations. If this happens, then setting the massive spin-2 fields to zero is inconsistent with the higher-dimensional equations of motion, and so the reduction is established to be an inconsistent one. We derive this condition in section 2.
As we shall discuss, the absence or presence of these kinds of trilinear couplings is governed by whether or not the Killing vectors on the internal space satisfy a certain quadratic identity. We shall show that although the full set of $`SO(n+1)`$ Killing vectors on the sphere $`S^n`$ do indeed satisfy the identity, implying no inconsistency in this sector, the full sets of Killing vectors in the case of other internal manifolds do not. In particular, we shall show by this means that for the 5-dimensional space $`Q(1,1)`$ (sometimes called $`T^{11}`$), which can be described as a $`U(1)`$ bundle over $`S^2\times S^2`$, only the Killing vector of the $`U(1)`$ factor in its $`U(1)\times SU(2)\times SU(2)`$ isometry group satisfies the consistency condition. Thus in a reduction of type IIB supergravity on the $`Q(1,1)`$ space, only the $`U(1)`$ gauge field of the $`N=2`$ supergravity multiplet can be consistently retained in a massless truncation, whilst the $`SU2)\times SU(2)`$ gauge fields of the matter multiplets must be set to zero.
In order to demonstrate that the $`SU(2)\times SU(2)`$ Killing vectors of the $`Q(1,1)`$ space fail to satisfy the consistency criterion, it is helpful to obtain an explicit construction for them. Motivated by this, we have undertaken a rather more general investigation of the construction of Killing vectors in spaces of this kind. The base space $`S^2\times S^2`$ in the construction of $`Q(1,1)`$ as a $`U(1)`$ bundle is Kähler , and in fact in this specific case it itself is an Einstein space. More generally, one can consider the $`U(1)`$ bundle spaces over any Einstein-Kähler base space, or over a product of Einstein-Kähler spaces. Other relevant examples of this kind are the 7-dimensional $`M(3,2)`$ and $`Q(1,1,1)`$ spaces that have been used in compactifications of $`D=11`$ supergravity . These arise, respectively, as $`U(1)`$ bundles over $`CP^2\times S^2`$ and over $`S^2\times S^2\times S^2`$ . In all the cases, the curvature of the $`U(1)`$ connection is proportional to the sum of the Kähler forms on the factors in the base space.
Intuitively, one expects that if the base space has an isometry group $`G`$, and the curvature of the $`U(1)`$ bundle is invariant under $`G`$, then the isometry group of the bundle space should be at least $`U(1)\times G`$. In section 3 we show how to make this precise, and we obtain explicit formulae that allow one to “lift” the Killing vectors of the base space to Killing vectors in the bundle space. The situation is especially nice if the base space is Einstein-Kähler , or else a product of Einstein-Kähler spaces, and we show how one can then express the Killing vectors of the base, and hence of the bundle space, in terms of certain scalar harmonics on the Einstein-Kähler factors in the base space.
In section 4 we specialise the discussion to the case where the Einstein-Kähler manifolds $`M_i`$ in the product base space $`M=M_1\times M_2\times M_N`$ are taken to be $`M_i=CP^{n_i}`$. We denote the corresponding bundle spaces by $`Q_{n_1n_2\mathrm{}n_N}^{q_1q_2\mathrm{}q_N}`$, where $`q_i`$ is the winding number of the $`U(1)`$ fibre over the factor $`M_i=CP^{n_i}`$ in the product base manifold. The three examples $`Q(1,1)`$, $`M(3,2)`$ and $`Q(1,1,1)`$ described above are special cases within this general class, namely $`Q_{11}^{11}`$, $`Q_{21}^{32}`$ and $`Q_{111}^{111}`$ respectively. We give an explicit construction of the $`SU(n+1)`$ Killing vectors of $`CP^n`$ in terms of certain scalar harmonics. Using this construction and the results from section 3, we are able to lift all the Killing vectors of the product base manifold for $`Q_{n_1n_2\mathrm{}n_N}^{q_1q_2\mathrm{}q_N}`$ into the total bundle space, thereby exhibiting its $`U(1)\times _iSU(n_i+1)`$ isometry group.
We prove also that all the $`U(1)`$ bundle spaces $`Q_{n_1n_2\mathrm{}n_N}^{q_1q_2\mathrm{}q_N}`$ admit Einstein metrics of positive Ricci curvature, for all possible choices of the winding numbers $`q_i`$, provided only that they do not all vanish. In addition we show that when all the $`q_i`$ are given by $`q_i=n_i+1`$, the Einstein metric admits 2 Killing spinors.
In section 5 we make use of some of the general results from sections 3 and 4, to show explicitly that the Killing vectors in the $`SU(n_i+1)`$ factors in spaces such $`Q_{11}^{11}`$, $`Q_{21}^{32}`$ and $`Q_{111}^{111}`$ do not satisfy the consistency criterion for the Kaluza-Klein reductions. These results support the suggestion, made in , that only the massless fields in the supergravity multiplet, as opposed to any massless matter multiplets, can be consistently retained in a Kaluza-Klein reduction using a curved internal space. Thus the reason why spheres work so well in Kaluza-Klein supergravity reductions is because they maximise the number of Killing spinors, and thus their supergravity multiplets are larger than those for any other choice of compactifying space. More generally, in section 5, we analyse the analogue of the consistency condition for bundle spaces of arbitrary dimension, and we show that always the Killing vectors associated with the isometries of the base manifold, when it is a product of two or more complex projective spaces, do not satisfy the consistency condition. Two appendices contain some further general results, including an iterative construction of real metrics on $`CP^n`$, and a detailed analysis of certain bounds on integrals involving the scalar eigenfunctions on $`CP^n`$, which are needed for the results in section 5.
## 2 Consistency conditions on Killing vectors in Kaluza-Klein reductions
In this section, we shall focus principally on the Kaluza-Klein reduction of type IIB supergravity on a 5-dimensional internal space $`_5`$. Analogous results have previously been obtained for reductions of $`D=11`$ supergravity , and we shall mention these briefly at the end of the section.
Since our goal will be to derive a necessary condition for the consistency of the reduction, with a view to showing that the condition is not in fact satisfied except in very special circumstances, it will be sufficient to carry out an analysis that is based on a linearised approximation. Thus we shall consider a situation where $`_5`$ is an Einstein space of positive Ricci curvature, and we shall consider small fluctuations around the AdS$`{}_{5}{}^{}\times _5`$ Freund-Rubin background. In particular, we shall consider the Yang-Mills gauge bosons associated with the isometry group $`G`$ of the internal space $`_5`$.
Although we shall consider only the linearised ansatz for the gauge bosons this will actually enable us to consider the effects of non-linear terms in theory, and in particular to show that bilinears in the gauge fields will in general act as sources for massive spin-2 fields. The reason why we can use a linearised ansatz for this purpose is that gauge invariance ensures that there can be no additional contributions from a full non-linear reduction ansatz that could “help out” and resolve the consistency problems that we shall be able to reveal. Thus, since our goal is only to prove inconsistency, not consistency, the analysis presented here will be sufficient.<sup>3</sup><sup>3</sup>3Note that we shall ignore the contributions of other five-dimensional fields, including the scalar fields, in this discussion. Truncating out these fields, while keeping the Yang-Mills gauge fields, is itself an inconsistent procedure, since the Yang-Mills fields would in principle act as sources for them. The point is, though, that these are quite distinct and separate inconsistencies, which would show up in different sectors of the theory. By focusing, as we shall, on the five-dimensional spacetime components of the ten-dimensional Einstein equation we shall be able to isolate a particular inconsistency that is independent of the neglect of the other fields. In other words, including the other fields in the ansatz would not help to resolve the inconsistency that we shall exhibit.
The fields of the type IIB theory that are relevant for this discussion are the metric tensor $`\widehat{G}_{MN}`$ and the self-dual 5-form field strength $`\widehat{H}_5`$, which we may write as $`\widehat{H}_5=\widehat{G}_5+\widehat{}\widehat{G}_5`$. The type IIB equations of motion for these fields are then
$`\widehat{R}_{MN}`$ $`=`$ $`{\displaystyle \frac{1}{96}}\widehat{H}_{MPQRS}\widehat{H}_N{}_{}{}^{PQRS},`$
$`d\widehat{H}_5`$ $`=`$ $`d\widehat{H}_5=0.`$ (2.1)
The Freund-Rubin AdS$`{}_{5}{}^{}\times _5`$ ground-state solution is then obtained by setting $`\widehat{G}_5=4mϵ_5`$, where $`ϵ_5`$ is the spacetime volume form and $`m`$ is a constant. The equations of motion are then satisfied if the Ricci tensors in the five-dimensional spacetime and the internal space $`_5`$ satisfy
$$R_{\mu \nu }=4m^2g_{\mu \nu },\mathrm{and}R_{mn}=4m^2g_{mn}$$
(2.2)
respectively. Thus we may take the spacetime metric to be AdS<sub>5</sub> with cosmological constant $`4m^2`$, and $`_5`$ to be any 5-dimensional Einstein space with positive cosmological constant<sup>4</sup><sup>4</sup>4We shall adopt the convention throughout this paper of referring to the constant of proportionality $`\mathrm{\Lambda }`$ in the relation $`R_{ab}=\mathrm{\Lambda }g_{ab}`$ on an Einstein space as the cosmological constant. It sometimes differs by a dimension-dependent factor from other terminologies in the literature, but this one has the merit of simplicity. $`4m^2`$.
We may now consider the contributions of the five-dimensional Yang-Mills gauged bosons in the ansätze for the ten-dimensional metric and 5-form field strength, at the leading-order linearised level. For the metric, this will be
$$d\widehat{s}^2=e^\alpha e^\beta \eta _{\alpha \beta }+(e^aK^{Ia}A^I)(e^bK^{Jb}A^J)\delta _{ab},$$
(2.3)
where $`e^\alpha =e^\alpha (x)`$ is the vielbein in the $`d=5`$ spacetime, $`e^a=e^a(y)`$ is the vielbein in the internal space, $`K^{Ia}=K^{Ia}(y)=K^{Im}(y)e_{m}^{}{}_{}{}^{a}(y)`$ are the orthonormal components of the Killing vectors which generate the isometry group $`G`$ of the internal space $`_5`$, and $`A^I=A^I(x)=e^\alpha (x)A_\alpha ^I(x)`$ are the Yang-Mills vector potentials of the Kaluza-Klein reduction. The Killing vectors satisfy
$$[K^I,K^J]=f_{}^{IJ}{}_{K}{}^{}K^K,$$
(2.4)
where $`f_{}^{IJ}{}_{K}{}^{}`$ are the structure constants of $`G`$. The Yang-Mills field strengths $`F^I=\frac{1}{2}F_{\alpha \beta }^Ie^\alpha e^\beta `$ are given by
$$F^I=dA^I+\frac{1}{2}f^{IJK}A^JA^K.$$
(2.5)
In an orthonormal basis $`\widehat{e}^A`$ for $`d\widehat{s}^2`$ we find that the Ricci tensor given by
$`\widehat{R}_{\alpha \beta }`$ $`=`$ $`R_{\alpha \beta }\frac{1}{2}K^{Ia}K_{}^{J}{}_{a}{}^{}F_{}^{I}{}_{\alpha \gamma }{}^{}F_{}^{J}{}_{\beta }{}^{}{}_{}{}^{\gamma },`$
$`\widehat{R}_{ab}`$ $`=`$ $`R_{ab}+{\displaystyle \frac{1}{4}}K_{}^{I}{}_{a}{}^{}K_{}^{J}{}_{b}{}^{}F_{}^{I}{}_{\alpha \beta }{}^{}F^{J\alpha \beta },`$
$`\widehat{R}_{\alpha b}`$ $`=`$ $`\widehat{R}_{b\alpha }=\frac{1}{2}K_{}^{I}{}_{b}{}^{}(D_\beta F_{}^{I}{}_{\alpha }{}^{}{}_{}{}^{\beta }),`$ (2.6)
where $`D_\gamma `$ is the Yang-Mills gauge-covariant derivative,
$$D_\gamma F_{}^{I}{}_{\alpha \beta }{}^{}=_\gamma F_{}^{I}{}_{\alpha \beta }{}^{}+f^{IJK}A_{}^{I}{}_{\gamma }{}^{}F_{}^{K}{}_{\alpha \beta }{}^{}.$$
(2.7)
The curvature scalar is
$$\widehat{R}=R_{(5)}+R_{(M)}\frac{1}{4}K_{}^{I}{}_{a}{}^{}K^{Ja}F_{}^{I}{}_{\alpha \beta }{}^{}F^{J\alpha \beta },$$
(2.8)
where $`R_{(5)}`$ and $`R_{(M)}`$ are the curvature scalars in spacetime and the internal space $`M`$ respectively.
The gauge fields also enter in the linearised ansatz for the 5-form field strength , as follows:
$$\widehat{G}_5=4mϵ_5\frac{1}{m}F^IdK^I.$$
(2.9)
Substituting (2.6), (2.8) and (2.9) into (2.1) we find that the five-dimensional spacetime components of the ten-dimensional Einstein equation in (2.1) are given by
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }(R_{(5)}+R_{(M)})=\frac{1}{2}(F_{}^{I}{}_{\mu \rho }{}^{}F_{}^{J}{}_{\nu }{}^{\rho }\frac{1}{4}g_{\mu \nu }F_{}^{I}{}_{\sigma \rho }{}^{}F_{}^{J}{}_{}{}^{\sigma \rho })Y^{IJ},$$
(2.10)
where
$$Y^{IJ}=Y(K^I,K^J)K_{}^{I}{}_{m}{}^{}K^{Jm}+\frac{1}{2m^2}_mK_{}^{I}{}_{n}{}^{}^mK^{Jn}.$$
(2.11)
The possibility of an inconsistency in the Kaluza-Klein reduction becomes apparent from equation (2.10). The left-hand side is independent of the coordinates $`y`$ on the internal space $`_5`$, while the right-hand side is in general $`y`$-dependent, since the Killing appearing in $`Y^{IJ}`$ are in general $`y`$-dependent. If the right-hand side does have $`y`$-dependence then this is an indication that the assumption that only the massless spin-2 field (the five-dimensional spacetime metric) could be retained in the truncation is an invalid one. One can interpret any $`y`$-dependence on the right-hand side as indicating that there are bilinear terms, built from the Yang-Mills field strengths, that would act as sources for massive spin-2 fields. Thus it would be inconsistent to make a truncation where the massless gauge bosons are retained, while the massive spin-2 fields are set to zero.
By contrast, this inconsistency problem would be evaded if the quantity $`Y_{IJ}`$ defined in (2.11) happened to be independent of $`y`$. In such a case one could, by taking appropriate linear combinations of the Killing vectors, arrange that
$$Y^{IJ}=\beta \delta ^{IJ},$$
(2.12)
where $`\beta `$ is a constant. In this circumstance, (2.10) would become precisely the desired five-dimensional Einstein equation, with the right-hand side being the energy-momentum tensor of the Yang-Mills fields.
Remarkably, all the Killing vectors on the round 5-sphere do satisfy the condition (2.12), thus providing strong evidence for the probable consistency of the $`S^5`$ reduction of type IIB supergravity. On the other hand, it seems that for any other Einstein space $`_5`$ with positive cosmological constant, the Killing vectors do not in general satisfy the condition (2.12), and thus in these cases a consistent massless truncation in which all the Yang-Mills gauge bosons are retained is not possible.
Note that there is an analogous criterion for the consistency of reductions of $`D=11`$ supergravity. This was derived for compactifications on seven-dimensional Einstein spaces in , and takes the identical form (2.11) where the Ricci tensor on the internal seven-dimensional space is given by $`R_{mn}=6m^2g_{mn}`$. Similarly, the analogous consistency condition will arise for reductions of $`D=11`$ supergravity on four-dimensional internal spaces, and indeed for any of the cases where consistent sphere reductions are known to be possible. A detailed enumeration of these cases is given in . In fact in general one can show that for any round sphere $`S^n`$ with $`R_{mn}=(n1)m^2g_{mn}`$, the condition (2.12) is satisfied by all the $`SO(n+1)`$ Killing vectors. (This does not necessarily mean that a consistnet Kaluza-Klein reduction on $`S^n`$ is possible, though.)
It is worth pausing here to emphasise that although we have derived the consistency condition that (2.11) must be constant by means of a consideration only of the linearised ansatz for the Kaluza-Klein reduction of the gauged fields, the result is a completely general one. The reason for this is discussed in detail in ). The crucial point is the following. If (2.11) is $`y`$-dependent, this shows that in a complete Kaluza-Klein reduction in which all the massive as well as massless fields were retained, there would be trilinear couplings involving one power of a heavy spin-2 field, say $`H_{\mu \nu }^{IJ}`$, coupling to the bilinear source term quadratic in gauge fields $`F_{\mu \rho }^I`$ on the right-hand side of (2.10):
$$_{\mathrm{int}}=H^{IJ\mu \nu }(F_{}^{I}{}_{\mu \rho }{}^{}F_{}^{J}{}_{\nu }{}^{\rho }\frac{1}{4}g_{\mu \nu }F_{}^{I}{}_{\sigma \rho }{}^{}F_{}^{J}{}_{}{}^{\sigma \rho }).$$
(2.13)
Now, the masses of all the lower-dimensional massive fields are acquired through a Higgs mechanism, and so it follows that the original gauge invariances must remain unbroken. In consequence, the lower-dimensional massive spin-2 field $`H_{\mu \nu }^{IJ}`$ must have a gauge invariance, implying that the source-current that couples to it must be conserved . Indeed, from an order-by-order analysis it follows that the bilinear current must be conserved by virtue of the free field equations. The bilinear current
$$(F_{}^{I}{}_{\mu \rho }{}^{}F_{}^{J}{}_{\nu }{}^{\rho }\frac{1}{4}g_{\mu \nu }F_{}^{I}{}_{\sigma \rho }{}^{}F_{}^{J}{}_{}{}^{\sigma \rho })$$
(2.14)
appearing in (2.13) is the unique one with this property, and so it is not possible for it to receive any corrections as a result of including higher-order non-linear terms in the reduction ansatz. Thus there is no possibility that the inconsistency we are highlighting could “disappear” in a more complete higher-order analysis. If the quantity (2.11) turns out to be $`y`$-dependent, then no consistent Kaluza-Klein reduction in which the associated gauge fields are retained is possible.
In order to proceed with our discussion, it is now necessary to study in detail the Killing vectors on the internal space. We shall give a rather general discussion, which encompasses many of the compactifications of type IIB supergravity and eleven-dimensional supergravity as special cases. Later, in section 5, we shall apply these results to study the consistency of the Kaluza-Klein reductions.
## 3 Construction of Killing vectors on the internal space
### 3.1 Killing vectors on $`U(1)`$ bundles
Consider a $`D`$-dimensional manifold with a group $`G`$ of isometries. Suppose that there exists a $`U(1)`$ connection on $`M`$ whose curvature is invariant under the isometry group $`G`$. One then expects that the natural metric on the $`(D+1)`$-dimensional bundle space with $`U(1)`$ fibers corresponding to this $`G`$-invariant $`U(1)`$ connection should contain $`G\times U(1)`$ in its isometry group.<sup>5</sup><sup>5</sup>5Generically, this will be the full isometry group of the bundle space, but in special cases it could be a larger group containing $`G\times U(1)`$ as a subgroup. An example is when the base manifold is $`CP^n`$, and the bundle space is the sphere $`S^{2n+1}`$. If the length of the $`U(1)`$ fibres is chosen so as to give the “round” sphere, then the generic $`SU(n+1)\times U(1)`$ isometry group in the bundle enlarges to $`SO(2n+2)`$.
To see that this is indeed the case, suppose that the metric on the base manifold is $`ds^2`$, and that the invariant $`U(1)`$ connection is $`A`$. The natural metric on the $`(n+1)`$-dimensional bundle space is then taken to be
$$d\widehat{s}^2=c^2(dzA)^2+ds^2,$$
(3.1)
where $`c`$ is a constant, and $`z`$ is the coordinate on the $`U(1)`$ fibre. (From this point on, we adopt the convention that quantities with hats refer to the total bundle space $`\widehat{M}`$, while quantities without hats refer to the base space $`M`$. We shall use indices $`m,n,\mathrm{}`$ in the total bundle space $`\widehat{M}`$, and indices $`a,b,\mathrm{}`$ in the base space $`M`$.) In the obvious orthonormal frame, the Riemann tensor for $`d\widehat{s}^2`$ has components
$`\widehat{R}_{abcd}`$ $`=`$ $`R_{abcd}\frac{1}{4}c^2(F_{ac}F_{bd}F_{ad}F_{bc}+2F_{ab}F_{cd}),`$
$`\widehat{R}_{zazb}`$ $`=`$ $`\frac{1}{4}c^2F_a{}_{}{}^{c}F_{bc}^{},`$ (3.2)
$`\widehat{R}_{abcz}`$ $`=`$ $`\frac{1}{2}c_cF_{ab},`$
where $`R_{abcd}`$ is the Riemann tensor of the metric $`ds^2`$ on the base manifold.
From (3.2) it follows that the components of the Ricci tensor for $`d\widehat{s}^2`$ are
$`\widehat{R}_{ab}`$ $`=`$ $`R_{ab}\frac{1}{2}c^2F_{ac}F_b{}_{}{}^{c},`$
$`\widehat{R}_{zz}`$ $`=`$ $`\frac{1}{4}c^2F^{ab}F_{ab},`$ (3.3)
$`\widehat{R}_{az}`$ $`=`$ $`\frac{1}{2}c^bF_{ab},`$
where $`R_{ab}`$ is the Ricci tensor on the base manifold.
We now make the following ansatz in order to lift a Killing vector $`K`$ on the base manifold $`M`$ to a Killing vector $`\widehat{K}`$ on the total bundle space $`\widehat{M}`$:
$$\widehat{K}=K+h_z,$$
(3.4)
where $`h`$ is a function to be determined. By substituting our ansatz (3.4) into the Killing equation of the bundle space:
$$\widehat{}_m\widehat{K}_n+\widehat{}_n\widehat{K}_m=0,$$
(3.5)
we find that $`\widehat{K}`$ is a Killing vector on the bundle space provided that $`h`$ satisfies the following two equations:
$`_ah`$ $`=`$ $`_KA_a,`$ (3.6)
$`_zh`$ $`=`$ $`0,`$ (3.7)
where $`_KA_a`$ is the Lie derivative, defined by
$$_KA_aK^b_bA_a+A_b_aK^b=K^b_bA_a+A_b_aK^b.$$
(3.8)
Equation (3.6) can be rewritten in terms of the field strength $`F=dA`$ as
$$_ah=K^bF_{ba}+_a(K^bA_b).$$
(3.9)
It is easy to see that the two equations (3.9) and (3.7) always admit a solution, provided that $`F`$ is invariant under the action of the Killing symmetry generated by $`K`$. Clearly (3.7) is nothing more than the statement that $`h`$ is independent of the fibre coordinate $`z`$. The integrability condition for solving (3.9) for $`h`$ is that the right-hand side should be expressible as the gradient of a scalar. Since the second term is already a gradient, this means that we must just show that $`_c(K^bF_{ba})_a(K^bF_{bc})=0`$. Calculating this expression, we find
$$_c(K^bF_{ba})_a(K^bF_{bc})=K^b_bF_{ac}F_{ab}_cK^bF_{bc}_aK^b,$$
(3.10)
which is nothing but the Lie derivative $`_KF_{ca}`$. This vanishes precisely by virtue of the assumption that $`F`$ is invariant under the Killing symmetry.
The above argument establishes that every Killing vector on the base manifold lifts to one in the total bundle space. In addition to these Killing vectors of the isometry group $`G`$ of the base manifold, there will also be the $`U(1)`$ Killing vector $`/z`$ on the $`U(1)`$ fibres. Thus the isometry group of the total bundle space will be at least $`G\times U(1)`$.
We are interested in obtaining an explicit construction of the Killing vectors in certain Einstein spaces that can be used for Kaluza-Klein reduction, in order to test the consistency as described in section 2. In all the examples that we shall consider, the Einstein space can be constructed as a $`U(1)`$ bundle over a Kähler base manifold. More specifically, in all cases of interest this Kähler space will itself be a direct product of Einstein-Kähler spaces. The additional structure of the Kähler spaces allows us to obtain more explicit constructions for the Killing vectors in the bundle space.
### 3.2 Killing vectors on Kähler spaces and their $`U(1)`$ bundle spaces
We begin with a review of some basic properties of Killing vectors and Kähler spaces. Consider a compact Kähler manifold $`M`$ equipped with a positive definite metric $`g_{ab}`$, and a Kähler form $`J_{ab}`$. We are interested in the case where $`M`$ has continuous isometries, and hence admits Killing vectors. It follows from the defining equation $`_aK_b+_bK_a=0`$ for a Killing vector that
$$\text{ }\text{ }K_a+R_{ab}K^b=0.$$
(3.11)
Multiplying by $`K^a`$, integrating over $`M`$, and integrating by parts, gives
$$_M(|_aK_b|^2+R_{ab}K^aK^b)=0,$$
(3.12)
where $`|_aK_b|^2`$ means $`(_aK_b)(^aK^b)`$. The metric is positive-definite, and so from (3.12) we deduce that for Killing vectors to exist, there must be appropriate non-negative contributions from the Ricci-tensor term. In fact we are interested in the case where $`R_{ab}`$ is positive definite.
Another consequence that follows from the positivity of the Ricci tensor is that the first Betti number $`b_1`$ of the Kähler space must be zero. This follows from an argument precisely paralleling the one above concerned with the possibility of the existence of Killing vectors. A harmonic 1-form $`H_a`$ satisfies the equation $`\text{ }\text{ }H_a+R_{ab}H^b=0`$. By multiplying by $`H^a`$, integrating over $`M`$, and integrating by parts on the first term, we see that there can be no harmonic 1-forms if the Ricci tensor is positive definite. Since we shall be considering spaces that have strictly positive-definite Ricci tensors, it follows that they will have $`b_1=0`$, and admit no harmonic 1-forms.
Consider now the vector $`V^a`$, constructed from the Killing vector $`K^a`$ as follows:
$$V^aJ^a{}_{b}{}^{}K_{}^{b}.$$
(3.13)
Our first goal will be to show that $`V^a`$ can be written as the gradient of a scalar function. To prove this, define
$$Q_{ab}_aV_b_bV_a.$$
(3.14)
It follows that
$`|Q_{ab}|^2`$ $`=`$ $`(_aV_b_bV_a)(^aV^b^bV^a),`$ (3.15)
$`=`$ $`2(_aK_b)(^aK^b)2J^{ad}J^{cb}(_aK_b)(_cK_d).`$
Integrating this over $`M`$, and integrating by parts on each term, we obtain
$`{\displaystyle _M}|Q_{ab}|^2`$ $`=`$ $`2{\displaystyle _M}K^a\text{ }\text{ }K_a+2{\displaystyle _M}J^{ad}J^{cb}K_b_a_cK_d,`$ (3.16)
$`=`$ $`2{\displaystyle _M}K^a\text{ }\text{ }K_a+2{\displaystyle _M}J^{ad}J^{cb}K_bR^e{}_{acd}{}^{}K_{e}^{},`$
$`=`$ $`2{\displaystyle _M}K^a(\text{ }\text{ }K_a+R_{ab}K^b),`$
$`=`$ $`0,`$
where we have used the standard Killing-vector identity $`_a_cK_d=R^e{}_{acd}{}^{}K_{e}^{}`$ in reaching the second line, and the standard Kähler identity $`R_{abcd}=J_c{}_{}{}^{e}J_{d}^{}{}_{}{}^{f}R_{abef}^{}`$ in reaching the third line. The final result follows from using (3.11). Thus we conclude that $`Q_{ab}=0`$, and hence that $`V_a`$, viewed as a 1-form, is closed; $`dV=0`$. Locally, therefore, we can write $`V=d\psi `$. As we discussed previously, we shall be interested in Kähler spaces with positive-definite Ricci tensor, and such spaces have vanishing first Betti number. Since there are no harmonic 1-forms in such spaces, it follows that $`dV=0`$ can be solved globally by writing $`V=d\psi `$. In other words, we have the result that on a Kähler space with vanishing first Betti number, any Killing vector can be written as
$$K^a=J^{ab}_b\psi ,$$
(3.17)
for some scalar $`\psi `$.
This scalar $`\psi `$ has a clear interpretation if we impose that our Kähler space is also an Einstein space, $`R_{ab}=\mathrm{\Lambda }g_{ab}`$, where $`\mathrm{\Lambda }`$ is the “cosmological constant” on $`M`$. Then (3.11) reduces to
$$\text{ }\text{ }K_a+\mathrm{\Lambda }K_a=0.$$
(3.18)
It is now straightforward to see, by substituting (3.17) into (3.18), that this scalar field is actually an eigenfunction of the Laplacian on $`M`$, $`\text{ }\text{ }\psi =\lambda \psi `$, with
$$\text{ }\text{ }\psi +2\mathrm{\Lambda }\psi =0.$$
(3.19)
Moreover, the implication goes in the other direction as well. In other words, if $`\psi `$ is an eigenfunction of the scalar Laplacian, satisfying (3.19), then $`K^a`$ defined by (3.17) is a Killing vector. To see this, we define $`P_{ab}_aK_b+_bK_a`$. Substituting (3.17) into this, writing $`|P|^2`$, and then performing appropriate integrations by parts, we find that
$$_M|P_{ab}|^2=2(\lambda 2\mathrm{\Lambda })_M|_a\psi |^2.$$
(3.20)
(Again, standard results from Kähler geometry are needed in intermediate steps.) Thus if the scalar eigenfunction $`\psi `$ has eigenvalue $`\lambda =2\mathrm{\Lambda }`$, it follows that $`P_{ab}=0`$ and hence that $`K_a`$ constructed as in (3.17) is a Killing vector.
Thus we see that there is a one-to-one correspondence between Killing vectors, and scalar eigenfunctions with eigenvalue $`\lambda =2\mathrm{\Lambda }`$, where $`\mathrm{\Lambda }`$ is the cosmological constant of the Einstein-Kähler space.
Using (3.17), we can now obtain explicit expressions for the Killing vectors on the space of the $`U(1)`$ bundle over an Einstein-Kähler base space, where the curvature of the $`U(1)`$ connection is taken to be proportional to the Kähler form. Taking the Einstein-Kähler base metric to have cosmological constant $`\mathrm{\Lambda }`$ as above, and taking the field strength of the connection $`A`$ on the $`U(1)`$ bundle to be $`F=\alpha J`$, where $`\alpha `$ is a constant, it follows from (3.3) that the Ricci tensor on the bundle space will be given by
$$\widehat{R}_{ab}=(\mathrm{\Lambda }\frac{1}{2}c^2\alpha ^2)\delta _{ab},\widehat{R}_{zz}=\frac{1}{4}c^2\alpha ^2D,\widehat{R}_{az}=0,$$
(3.21)
where $`D`$ is the dimension of the base manifold. In particular, the metric on the $`U(1)`$ bundle becomes Einstein if $`a`$ is chosen such that
$$\mathrm{\Lambda }=\frac{1}{4}c^2\alpha ^2(D+2).$$
(3.22)
Substituting (3.17) into (3.9), we now obtain the result that
$$_ah=_a(\alpha \psi +K^bA_b),$$
(3.23)
which can be integrated to give $`h=\alpha \psi +K^bA_b`$. Thus for each Killing vector $`K`$ on the Einstein-Kähler base space, with its associated scalar $`\psi `$ as given in (3.17), the corresponding Killing vector in the $`U(1)`$ bundle space is
$$\widehat{K}=K+(\alpha \psi +K^bA_b)\frac{}{z}.$$
(3.24)
It is worth noting at this stage that there is an elegant expression for the Killing vector $`\widehat{K}`$, viewed as a 1-form by lowering its vector index using the metric (3.1) on the bundle space. After doing this, we find that as a 1-form we have
$$\widehat{K}=\mathrm{i}(\overline{})\psi +\alpha c^2\psi (dzA),$$
(3.25)
where $``$ and $`\overline{}`$ are the holomorphic and anti-holomorphic exterior derivatives; $`d=+\overline{}=d\zeta _\alpha _\alpha +d\overline{\zeta }^{\overline{\alpha }}_{\overline{\alpha }}`$.
Note that in the case of an Einstein-Kähler base space we can easily express $`\psi `$ in terms of the Killing vector $`K^a`$, since from (3.17) we have $`J^{ab}_aK_b=\text{ }\text{ }\psi `$, and hence from (3.19) we shall have
$$\psi =\frac{1}{2\mathrm{\Lambda }}J^{ab}_aK_b.$$
(3.26)
### 3.3 Killing vectors on a product of Kähler spaces and their $`U(1)`$ bundles
In subsequent sections, we shall be interested in constructing Killing vectors on $`U(1)`$ bundles over products of 2-spheres and more generally complex projective spaces $`CP^n`$. These are are particular examples of Einstein-Kähler spaces. In this section, we shall give results for the construction of Killing vectors on $`U(1)`$ bundles over the direct product of $`N`$ Einstein-Kähler spaces $`M_i`$, of real dimensions $`d_i`$, i.e. $`M=M_1\times M_2\times \mathrm{}\times M_N`$, with total real dimension
$$D=\underset{i=1}{\overset{N}{}}d_i.$$
(3.27)
The metric on the bundle space will be given by
$$d\widehat{s}^2=c^2(dzA)^2+\underset{i=1}{\overset{N}{}}ds_i^2,$$
(3.28)
where $`ds_i^2`$ is the Einstein-Kähler metric on the factor $`M_i`$ in the base space, with cosmological constant $`\mathrm{\Lambda }_i`$. The total connection $`A`$ is equal to the sum of contributions from each factor, $`A=_iA^{(i)}`$.
Since the base space is a direct product, we can choose the natural block-diagonal basis for its Killing vectors, where there is no mixing between the isometries of each factor in the product. Thus if $`K^{(i)}`$ is a Killing vector on $`M_i`$ then it is also a Killing vector on $`M`$, and vice versa. If we use $`a_i`$ to denote a coordinate index on $`M_i`$, then this result follows by combining the Killing equation $`_{a_i}K_{b_i}+_{b_i}K_{a_i}=0`$ on $`M_i`$, with the fact that the $`K^{(i)}`$ are covariantly constant with respect to $`_{b_j}`$ for $`ji`$.
This fact allows us to use the results of the previous section to express any of these Killing vectors $`K^{(i)}`$ on $`M`$ as
$$K^{a_i}=J^{a_ib_i}_{b_i}\psi ^{(i)},,$$
(3.29)
where $`\psi ^{(i)}`$ is the corresponding scalar eigenfunction of the Laplacian on $`M_i`$ with eigenvalue $`2\mathrm{\Lambda }_i`$, and $`J_{a_ib_i}`$ are the components of the Kähler form on $`M_i`$. From the results in the previous section, it then follows that the corresponding Killing vector in the bundle space will be
$$\widehat{K}^{(i)}=K^{(i)}+(\alpha _i\psi ^{(i)}+K^{b_i}A_{b_i})\frac{}{z},$$
(3.30)
where $`A_{a_i}`$ is the contribution to $`A`$ from the factor $`M_i`$ in the base space, $`A=_iA^{(i)}`$, and we are taking
$$F=dA=\underset{i}{}\alpha _iJ^{(i)},$$
(3.31)
where $`J^{(i)}`$ is the Kähler form on $`M_i`$.
We may again obtain an elegant expression for the Killing vector viewed as a 1-form, generalising (3.25):
$$\widehat{K}^{(i)}=\mathrm{i}(\overline{})\psi ^{(i)}+\alpha _ic^2\psi ^{(i)}(dzA).$$
(3.32)
The period $`\mathrm{\Delta }z`$ of the fibre coordinate $`z`$ must be compatible with the integrals of $`F`$ over all 2-cycles in the base manifold. Specifically, we must have
$$\mathrm{\Delta }z=\frac{1}{q_k}_{\mathrm{\Sigma }_k}F,$$
(3.33)
where $`q_k`$ is an integer and $`\mathrm{\Sigma }_k`$ is any 2-cycle in the base space. We are taking each factor $`M_i`$ in the base space to be Einstein-Kähler , with cosmological constant $`\mathrm{\Lambda }_i`$, and so it follows that the Ricci form $`P^{(i)}`$ in $`M_i`$ is given by $`P^{(i)}=\mathrm{\Lambda }_iJ^{(i)}`$. Since $`1/(2\pi )P^{(i)}`$ defines the first Chern class of $`M_i`$, it follows that $`1/(2\pi )P^{(i)}=`$integer, where the integral is taken over any 2-cycle in $`M_i`$, whilst the integral will be zero for any 2-cycle in $`M_j`$ with $`ji`$. If we define $`k_i`$ to be the greatest common divisor of the integers obtained by integrating $`P^{(i)}`$ over all possible 2-cycles in $`M_i`$, then it follows from (3.33) that $`z`$ must have a period such that
$$\mathrm{\Delta }z=\frac{2\pi \alpha _ik_i}{\mathrm{\Lambda }_iq_i},$$
(3.34)
for all $`i`$, where the $`q_i`$ are integers. Thus we must have
$$\alpha _i=\frac{b\mathrm{\Lambda }_iq_i}{k_i},$$
(3.35)
where $`b`$ is related to the period $`\mathrm{\Delta }z`$ by $`b=\mathrm{\Delta }z/(2\pi )`$, and it is a constant independent of $`i`$. Since we have also included the constant $`c`$ in (3.1), we are free to choose $`b`$ at will, to give $`z`$ a convenient period. The integers $`q_i`$ can be thought of as the winding numbers of the $`U(1)`$ bundle over each factor $`M_i`$ in the product base manifold.
Note that it one wants the total $`U(1)`$ bundle space to be Einstein, with cosmological constant $`\widehat{\mathrm{\Lambda }}`$, then it follows from (3.3) that we must have
$`\widehat{\mathrm{\Lambda }}`$ $`=`$ $`\mathrm{\Lambda }_i\frac{1}{2}c^2\alpha _i^2,\mathrm{for}\mathrm{all}i,`$ (3.36)
$`\widehat{\mathrm{\Lambda }}`$ $`=`$ $`\frac{1}{4}c^2{\displaystyle \underset{i}{}}d_i\alpha _i^2,`$ (3.37)
where $`d_i`$ is the dimension of the manifold $`M_i`$, and $`\alpha _i`$ is given by (3.35). To solve these equations, one can view $`\widehat{\mathrm{\Lambda }}`$ and the winding numbers $`q_i`$ as freely-specifiable quantities, with the $`N`$ equations (3.36) then being solved for the individual cosmological constants $`\mathrm{\Lambda }_i`$ of the factors in the base space, and (3.37) being solved for the scale factor $`c`$ in fibre direction of the metric (3.28) on the bundle space. As we shall now show, one can always solve these equations for the $`\mathrm{\Lambda }_i`$ and $`c`$, for any choice of the integers $`q_i`$, provided that they are not all zero.<sup>6</sup><sup>6</sup>6If some of the $`q_i`$ are zero, we can just separate off the corresponding Einstein spaces $`M_i`$ in the base space, and prove the existence of an Einstein metric on the bundle over the remaining base-space factors for which all the $`q_i`$ are non-zero. The product of this bundle space with the Einstein spaces associated with the $`q_i=0`$ factors can clearly be made Einstein, by appropriate choice of the $`\mathrm{\Lambda }_i`$. If all the $`q_i`$ were zero the $`U(1)`$ bundle would be trivial and the total $`(D+1)`$-dimensional space would be $`S^1\times M`$, which clearly cannot be Einstein since the factors in the base space $`M`$ are assumed to have strictly-positive cosmological constants.
To see this, we substitute (3.35) into (3.36), and note that for each $`i`$ the equation allows a real solution for $`\mathrm{\Lambda }_i`$ only if
$$c^2\alpha _i^2\frac{\mathrm{\Lambda }_i^2}{2\widehat{\mathrm{\Lambda }}}.$$
(3.38)
Summing over $`i`$, and using (3.37) then gives
$$\widehat{\mathrm{\Lambda }}^2\frac{1}{8}\underset{i}{}d_i\mathrm{\Lambda }_i^2.$$
(3.39)
On the other hand, combining (3.36) and (3.37) we have
$$\widehat{\mathrm{\Lambda }}=\frac{1}{D+2}\underset{i}{}d_i\mathrm{\Lambda }_i.$$
(3.40)
Combining (3.39) and (3.40) then gives the result
$$(D+2)^2\underset{i}{}d_i\mathrm{\Lambda }_i^28\underset{i,j}{}d_id_j\mathrm{\Lambda }_i\mathrm{\Lambda }_j0.$$
(3.41)
This is the criterion for the existence of a real Einstein space. Since it is just a quadratic form in $`\mathrm{\Lambda }_i`$, it can be expressed as the condition that the $`N\times N`$ matrix $`M_{ij}`$, defined by
$$M_{ij}=(D+2)d_i\delta _{ij}8d_id_j,$$
(3.42)
must have non-negative eigenvalues.
To show this, we first note that the matrix $`M_{ij}`$ has determinant given by
$$det(M_{ij})=(D2)^2(D+2)^{2N2}\underset{i}{}d_i,$$
(3.43)
which is strictly positive, since we may always assume $`D>2`$.<sup>7</sup><sup>7</sup>7The case where the total dimension $`D`$ of the base space is equal to 2 can easily be disposed of in a separate discussion. The only possibility would be for the base space to be $`S^2`$, and we already know that the $`U(1)`$ bundle over this is $`S^3`$, which admits an Einstein metric.
Secondly, we note that if the dimensions $`d_i`$ are all taken to be equal, $`d_i=d`$, then the eigenvalues of $`M_{ij}`$ are $`(D+2)^2d`$ (occurring $`N1`$ times) and $`(D2)^2d`$ (occurring once). Thus in this special case all the eigenvalues of $`M_{ij}`$ are strictly positive. If $`M_{ij}`$ were to have any negative values for any valid choice of the $`d_i`$, it would have to be the case that $`det(M_{ij})`$ passed through 0 as the parameters $`d_i`$ were adjusted from $`d_i=d`$ to these putative values of $`d_i`$. However, we saw from (3.43) that the determinant is strictly positive, and so it follows that $`M_{ij}`$ cannot have negative eigenvalues for any valid choice of $`d_i`$. Thus it is guaranteed that the inequality (3.41) is satisfied, and so a real solution to the conditions (3.36) and (3.37) always exists.
Although we have given an existence proof for an Einstein metric on the bundle spaces for any choice of the winding numbers $`q_i`$, it is not in general easy to solve explicitly for the cosmological constants $`\mathrm{\Lambda }_i`$ of the individual factors in the base space. (In general, one has to solve high-order polynomial equations.) However, a simple solution of (3.36) and (3.37) can always be explicitly obtained in the special case where we choose the winding numbers $`q_i`$ to be such that $`q_i=k_i/\mathrm{}`$, where $`\mathrm{}=\mathrm{gcd}(k_i)`$ is the greatest common divisor of the $`k_i`$. In this case, from (3.36) we see that this set of $`N`$ equations, labelled by $`i`$, all become equivalent. Therefore, defining $`\mathrm{\Lambda }\mathrm{\Lambda }_i`$ and $`\alpha \alpha _i`$, we have
$$\mathrm{\Lambda }=\frac{D+2}{D}\widehat{\mathrm{\Lambda }},$$
(3.44)
and
$$c^2\alpha ^2=\frac{4}{D+2}\mathrm{\Lambda },$$
(3.45)
where $`D`$ is the total dimension of the base manifold. Combining (3.35) and (3.45) it follows that the parameters of the metric satisfy the relation
$$\mathrm{\Lambda }b^2c^2=\frac{4\mathrm{}^2}{D+2}.$$
(3.46)
Note that since in this special case we have all the $`\mathrm{\Lambda }_i`$ equal, the product of Einstein-Kähler base spaces is itself an Einstein space. This situation with $`q_i=k_i/\mathrm{}`$ will be seen to be of particular significance in the next section, when we take the factors in the product base space all to be complex projective spaces. It turns out that the Einstein spaces with $`q_i=k_i/\mathrm{}`$ then all admit 2 Killing spinors.
## 4 Products of $`CP^n`$ spaces, and their $`U(1)`$ bundles
### 4.1 Geometry of $`CP^n`$, and its Killing vectors
We begin by reviewing the Fubini-Study construction of the Einstein-Kähler metric on $`CP^n`$. Let $`Z^A`$ be complex coordinates on $`C^{n+1}`$, with the flat metric
$$ds_{2n+2}^2=dZ^Ad\overline{Z}_A.$$
(4.1)
We shall split the index $`A`$ into $`A=(0,\alpha )`$, where $`1\alpha n`$, and introduce inhomogeneous coordinates $`\zeta ^\alpha =Z^\alpha /Z^0`$, in the patch where $`Z^00`$. We make the further definitions
$$Z^0=e^{\mathrm{i}\tau }|Z^0|,r=\sqrt{Z^A\overline{Z}_A},f=1+\zeta ^\alpha \overline{\zeta }^{\overline{\alpha }}.$$
(4.2)
Substituting into (4.1), we find that the flat metric on $`C^{n+1}`$ becomes
$$ds_{2n+2}^2=dr^2+r^2d\mathrm{\Omega }_{2n+1}^2,$$
(4.3)
where $`d\mathrm{\Omega }_{2n+1}^2`$ is the metric on the unit sphere $`S^{2n+1}`$, given by
$$d\mathrm{\Omega }_{2n+1}^2=(d\tau +B)^2+f^1d\zeta ^\alpha d\overline{\zeta }^{\overline{\alpha }}f^2\overline{\zeta }^{\overline{\alpha }}\zeta ^\beta d\zeta ^\alpha d\overline{\zeta }^{\overline{\beta }},$$
(4.4)
where
$$B=\frac{1}{2}\mathrm{i}f^1(\zeta ^\alpha d\overline{\zeta }^{\overline{\alpha }}\overline{\zeta }^{\overline{\alpha }}d\zeta ^\alpha ).$$
(4.5)
The metric (4.4) is the unit $`S^{2n+1}`$ described as a $`U(1)`$ bundle over $`CP^n`$, and the last two terms are precisely the Fubini-Study metric $`d\mathrm{\Sigma }_n^2`$ on $`CP^n`$:
$$d\mathrm{\Sigma }_n^2=f^1d\zeta ^\alpha d\overline{\zeta }^{\overline{\alpha }}f^2\overline{\zeta }^{\overline{\alpha }}\zeta ^\beta d\zeta ^\alpha d\overline{\zeta }^{\overline{\beta }},$$
(4.6)
and so
$$d\mathrm{\Omega }_{2n+1}^2=(d\tau +B)^2+d\mathrm{\Sigma }_n^2.$$
(4.7)
The quantity $`B`$ defined in (4.5) is a potential for the Kähler form, with
$$J=dB=\mathrm{i}f^1d\zeta ^\alpha d\overline{\zeta }^{\overline{\alpha }}+\mathrm{i}f^2\overline{\zeta }^{\overline{\alpha }}\zeta ^\beta d\zeta ^\alpha d\overline{\zeta }^{\overline{\beta }},$$
(4.8)
which is the Kähler form. This can be written as $`J=\mathrm{i}g_{\alpha \overline{\beta }}d\zeta ^\alpha d\overline{\zeta }^{\overline{\beta }}`$, where the metric $`g_{\alpha \overline{\beta }}`$ and its inverse $`g^{\alpha \overline{\beta }}`$ are given by
$$g_{\alpha \overline{\beta }}=\frac{1}{2}f^1\delta _{\alpha \overline{\beta }}\frac{1}{2}f^2\overline{\zeta }^{\overline{\alpha }}\zeta ^\beta ,g^{\alpha \overline{\beta }}=2f\delta ^{\alpha \overline{\beta }}+2f\zeta ^\alpha \overline{\zeta }^{\overline{\beta }}.$$
(4.9)
The Fubini-Study metric (4.6) is Einstein, with cosmological constant
$$\mathrm{\Lambda }=2(n+1).$$
(4.10)
We shall refer to the Fubini-Study metric (4.6) with this specific normalisation for the cosmological constant as the “unit $`CP^n`$ metric,” since it is the one that corresponds to the Hopf fibration of the unit $`(2n+1)`$-sphere. Note that $`CP^n`$ has the isometry group $`SU(n+1)`$, which can be seen from the fact that the metric (4.1) and the coordinate $`r`$ are both invariant under $`SU(n+1)`$, acting by matrix multiplication on the column vector $`Z^A`$.
Since we eventually want to be able to construct Killing vectors on $`U(1)`$ bundles over products of $`CP^n`$ spaces, we need to find the eigenfunctions of the scalar Laplacian on $`CP^n`$ with eigenvalue $`2\mathrm{\Lambda }`$, as discussed in section (3.2). In fact the construction of all scalar eigenfunctions on $`CP^n`$ is very simple. Let $`T_{A_1\mathrm{}A_p}^{B_1\mathrm{}B_q}`$ be a constant Hermitean $`SU(n+1)`$ tensor, which is symmetric in the index set $`\{A_1,\mathrm{},A_p\}`$ and the index set $`\{B_1,\mathrm{},B_q\}`$, and traceless in any contraction between an $`A`$ and a $`B`$ index. This defines the $`(p,q)`$ representation of $`SU(n+1)`$. Clearly the scalar function
$$\mathrm{\Phi }=T_{A_1\mathrm{}A_p}{}_{}{}^{B_1\mathrm{}B_q}Z_{}^{A_1}\mathrm{}Z^{A_p}\overline{Z}_{B_1}\mathrm{}\overline{Z}_{B_q}$$
(4.11)
is a zero mode of the Laplacian on $`C^{n+1}`$:
$$\text{ }\text{ }_{C^{n+1}}\mathrm{\Phi }=\frac{^2}{Z^A\overline{Z}_A}\mathrm{\Phi }=0,$$
(4.12)
where we can write this Laplacian in terms of $`r`$ and the Laplacian on the unit $`S^{2n+1}`$ as
$$0=\text{ }\text{ }_{C^{n+1}}\mathrm{\Phi }=\frac{1}{r^{2n+1}}\frac{}{r}\left(r^{2n+1}\frac{\mathrm{\Phi }}{r}\right)+\frac{1}{r^2}\text{ }\text{ }_{S^{2n+1}}\mathrm{\Phi }.$$
(4.13)
Note that $`\mathrm{\Phi }`$ can be written as
$$\mathrm{\Phi }=r^{p+q}e^{\mathrm{i}(pq)\tau }\mathrm{\Psi },$$
(4.14)
where $`\mathrm{\Psi }`$ depends only on the inhomogeneous $`CP^n`$ coordinates $`\zeta ^\alpha `$.
It is straightforward to show from (4.7) that the components of the sphere metric $`\widehat{g}_{AB}`$ and the $`CP^n`$ metric $`g_{ab}`$ are related by
$`\widehat{g}_{ab}=g_{ab}+B_aB_b,\widehat{g}_{a\tau }=B_a,\widehat{g}_{\tau \tau }=1,`$
$`\widehat{g}^{ab}=g^{ab},\widehat{g}^{a\tau }=B^a,\widehat{g}^{\tau \tau }=1+B_aB^a,`$ (4.15)
where $`B^ag^{ab}B_b`$. From this it is easily seen that the scalar Laplacian on $`S^{2n+1}`$ is given by
$$\text{ }\text{ }_{S^{2n+1}}=\left(_aB_a\frac{}{\tau }\right)^2+\frac{^2}{\tau ^2}.$$
(4.16)
Substituting (4.14) into (4.13) and (4.16), we therefore find that $`\mathrm{\Psi }`$ is an eigenfunction on $`CP^n`$, satisfying
$$𝒟_a𝒟^a\mathrm{\Psi }=2[2pq+n(p+q)]\mathrm{\Psi },$$
(4.17)
where $`𝒟_a=_a\mathrm{i}(pq)B_a`$. This is the Laplacian for scalar fields of charge $`(pq)`$, in the $`(p,0,0,\mathrm{},0,q)`$ representation of $`SU(n+1)`$. The uncharged scalars therefore occur in the $`(p,0,0,\mathrm{},0,p)`$ representations, with eigenvalues $`\lambda =4p(p+n)`$.
In section (3.2) the Killing vectors on an Einstein-Kähler space were constructed in terms of uncharged scalar eigenfunctions with eigenvalue $`2\mathrm{\Lambda }`$. On $`CP^n`$, the appropriate eigenfunctions are the ones with $`(p,q)=(1,1)`$, since, as can be seen from (4.17), they have eigenvalue $`4(n+1)`$, which, from (4.10), is $`2\mathrm{\Lambda }`$. They are indeed in the adjoint representation of $`SU(n+1)`$, as should be since they are supposed to be in one-to-one correspondence with the Killing vectors of $`CP^n`$.
Thus we see that the scalars $`\psi `$ that generate the Killing vectors on $`CP^n`$ are given by
$$\psi =\frac{1}{r^2}T_A{}_{}{}^{B}Z_{}^{A}\overline{Z}_B,$$
(4.18)
where $`T_A^B`$ is an arbitrary Hermitean traceless tensor. From the previous definitions, it has the following expression in terms of the inhomogeneous coordinates on $`CP^n`$:
$$\psi =f^1(T_0{}_{}{}^{0}+T_0{}_{}{}^{\alpha }\overline{\zeta }_{}^{\overline{\alpha }}+T_\alpha {}_{}{}^{0}\zeta _{}^{\alpha }+T_\alpha {}_{}{}^{\beta }\zeta _{}^{\alpha }\overline{\zeta }^{\overline{\beta }}).$$
(4.19)
Note that since $`T_A^B`$ is traceless, we can write $`T_0{}_{}{}^{0}=T_\alpha ^\alpha `$, and thus we can regard the unconstrained constant tensors $`T_0^\alpha `$, $`T_\alpha ^0`$ and $`T_\alpha ^\beta `$ as parameterising the set of scalars $`\psi `$ corresponding to the full set of $`n(n+2)`$ Killing vectors of $`CP^n`$.
From the scalars $`\psi `$, we can readily construct the Killing vectors using (3.17). From (4.9) we therefore find that the complex components of the Killing vector associated with $`\psi `$ are given by
$$K^\alpha =\mathrm{i}g^{\alpha \overline{\beta }}_{\overline{\beta }}\psi =\frac{\mathrm{i}}{2}(T_0{}_{}{}^{\alpha }+T_\beta {}_{}{}^{\alpha }\zeta _{}^{\beta }T_0{}_{}{}^{0}\zeta _{}^{\alpha }T_\beta {}_{}{}^{0}\zeta _{}^{\beta }\zeta ^\alpha ),$$
(4.20)
with $`K^{\overline{\alpha }}`$ being the complex conjugate of $`K^\alpha `$.
As a check on this construction of the Killing vectors from the scalar eigenfunctions $`\psi `$, we may also construct them directly, using the fact that they must correspond to infinitesimal $`SU(n+1)`$ transformations of the form $`\delta Z^A=\mathrm{i}ϵT_B{}_{}{}^{A}Z_{}^{B}`$ on the homogeneous coordinates, where $`T_B^A`$ is again an arbitrary Hermitean traceless tensor. This translates into $`\delta \zeta ^\alpha =\delta Z^\alpha /Z^0Z^\alpha /(Z^0)^2\delta Z^0`$, giving
$$\delta \zeta ^\alpha =\mathrm{i}ϵ(T_0{}_{}{}^{\alpha }+T_\beta {}_{}{}^{\alpha }\zeta _{}^{\beta }T_0{}_{}{}^{0}\zeta _{}^{\alpha }T_\beta {}_{}{}^{0}\zeta _{}^{\beta }\zeta ^\alpha ),$$
(4.21)
which is in precise agreement with (4.20), since Killing vectors generate the coordinate transformations $`\delta \zeta ^\alpha =2ϵK^\alpha `$. Of course we also need to know the explicit scalar functions $`\psi `$, for the purpose of lifting the Killing vectors to the $`U(1)`$ bundle space.
Note that $`CP^n`$ is a space of constant holomorphic sectional curvature, and in fact in terms of a real index notation the orthonormal components of the Riemann tensor of the unit $`CP^n`$ with metric (4.6) are given by
$$R_{abcd}=\delta _{ac}\delta _{bd}\delta _{ad}\delta _{bc}+J_{ac}J_{bd}J_{ad}J_{bc}+2J_{ab}J_{cd}.$$
(4.22)
It is sometimes useful to work with an explicit real metric for $`CP^n`$. In Appendix A, we obtain an iterative construction for a real metric on $`CP^n`$, in terms of a metric on $`CP^{n1}`$.
It is now straightforward to follow the procedure described in sections (3.2) and (3.3), to construct the $`U(1)`$ bundle space over an arbitrary product of $`CP^n`$ metrics. Specifically, we take the base manifold to be $`M=M_1\times M_2\times \mathrm{}\times M_N`$, where $`M_i`$ is the complex projective space $`CP^{n_i}`$, with real dimension $`d_i=2n_i`$. We shall denote the total bundle spaces by
$$Q_{n_1n_2\mathrm{}n_N}^{q_1q_2\mathrm{}q_N},$$
(4.23)
where the integers $`q_i`$ are the winding numbers of the $`U(1)`$ bundle over the factors $`CP^{n_i}`$ in the base manifold.
### 4.2 Killing spinors on $`Q_{n_1\mathrm{}n_N}^{q_1\mathrm{}q_N}`$ spaces
As we discussed in section 3.3, one can always find a solution to the conditions (3.36) and (3.37) for any choice of the $`q_i`$. A particularly simple case is when $`q_i=k_i`$. In fact in $`CP^{n_i}`$ there is only one 2-cycle, and the integer $`k_i`$ is therefore simply the result from integrating the first Chern class $`P_i/(2\pi )`$ over this cycle, which turns out to give
$$k_i=n_i+1.$$
(4.24)
In fact the Einstein spaces $`Q_{n_1n_2\mathrm{}n_N}^{q_1q_2\mathrm{}q_N}`$ with $`q_i=(n_i+1)/\mathrm{}`$ where $`\mathrm{}`$ is the greatest common divisor of the $`(n_i+1)`$ have a further nice feature, namely that they all admit Killing spinors. To show this, we note that the Killing spinor equation
$$D_A\eta \frac{\mathrm{i}}{2}\sqrt{\frac{\widehat{\mathrm{\Lambda }}}{D}}\mathrm{\Gamma }_A\eta =0$$
(4.25)
has the integrability condition
$$\frac{1}{4}\widehat{R}_{ABCD}\mathrm{\Gamma }^{CD}\eta \frac{\widehat{\mathrm{\Lambda }}}{2D}\mathrm{\Gamma }_{AB}\eta =0,$$
(4.26)
which is obtained by taking a commutator of the generalised derivatives appearing in (4.25). From (4.26) one can easily deduce that the metric on the total bundle space must be Einstein, and furthermore that
$$\widehat{C}_{ABCD}\mathrm{\Gamma }^{CD}\eta =0,$$
(4.27)
where $`\widehat{C}_{ABCD}`$ is the Weyl tensor on the total space.
If for every space $`CP^{n_i}`$ we take $`q_i=k_i/\mathrm{}`$, where $`\mathrm{}=\text{gcd}(k_i)`$, then we can use (3.44) and (3.45) to express the non-zero orthonormal components of the Riemann tensor on the $`U(1)`$ bundle space as:
$`\widehat{R}_{a_ib_ic_id_i}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{d_i+2}}\left(\delta _{a_ic_i}\delta _{b_id_i}\delta _{a_id_i}\delta _{b_ic_i}\right)`$
$`+\mathrm{\Lambda }\left[{\displaystyle \frac{1}{d_i+2}}{\displaystyle \frac{1}{D+2}}\right](J_{a_ic_i}J_{b_id_i}J_{a_id_i}J_{b_ic_i}+2J_{a_ib_i}J_{c_id_i}),`$
$`\widehat{R}_{a_ib_ia_jb_j}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Lambda }}{D+2}}J_{a_ib_i}J_{a_jb_j},`$ (4.28)
$`\widehat{R}_{a_ia_jb_ib_j}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{D+2}}J_{a_ib_i}J_{a_jb_j},`$
$`\widehat{R}_{\mathrm{\hspace{0.17em}0}a_i\mathrm{\hspace{0.17em}0}b_i}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{(D+2)}}\delta _{a_ib_i}`$
where $`D=_id_i=2n_i`$ is the total dimension of the base space, $`\mathrm{\Lambda }`$ is the (universal) cosmological constant of the $`CP^{n_i}`$, and the indices $`a_i`$ label the coordinates on $`CP^{n_i}`$. (We are using the expression (4.22) for the Riemann tensor of $`CP^n`$, appropriately rescaled so that the cosmological constant is $`\mathrm{\Lambda }`$.)
From (4.28) it follows that the non-zero components of the Weyl tensor are
$`\widehat{C}_{a_ib_ic_id_i}`$ $`=`$ $`\mathrm{\Lambda }[{\displaystyle \frac{1}{d_i+2}}{\displaystyle \frac{1}{D+2}}](\delta _{a_ic_i}\delta _{b_id_i}\delta _{a_id_i}\delta _{b_ic_i}+`$
$`J_{a_ic_i}J_{b_id_i}J_{a_id_i}J_{b_ic_i}+2J_{a_ib_i}J_{c_id_i}),`$
$`\widehat{C}_{a_ib_ia_jb_j}`$ $`=`$ $`{\displaystyle \frac{2\mathrm{\Lambda }}{D+2}}J_{a_ib_i}J_{a_jb_j},`$ (4.29)
$`\widehat{C}_{a_ia_jb_ib_j}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Lambda }}{D+2}}(\delta _{a_ib_i}\delta _{a_jb_j}+J_{a_ib_i}J_{a_jb_j}).`$
The integrability conditions (4.26) for the existence of Killing spinors therefore become
$`\mathrm{\Gamma }_{a_ib_i}\eta +J_{a_ic_i}J_{b_jd_j}\mathrm{\Gamma }_{c_id_j}\eta =0,`$ (4.30)
$`(Dd_i)(\mathrm{\Gamma }_{a_ib_i}+J_{a_ic_i}J_{b_id_i}\mathrm{\Gamma }_{c_id_i}+J_{a_ib_i}J_{c_id_i}\mathrm{\Gamma }_{c_id_i})\eta `$
$`(d_i+2)J_{a_ib_i}{\displaystyle \underset{ji}{}}J_{c_jd_j}\mathrm{\Gamma }_{c_jd_j}\eta =0.`$ (4.31)
One can show that (4.31) is implied by (4.30), and in fact the full set of independent conditions can be summarised succinctly as follows. Without loss of generality we can choose a basis for the $`CP^{n_i}`$ spaces in which the orthonormal components of the Kähler forms are:
$$J_{12}=J_{34}=J_{56}=\mathrm{}=+1,$$
(4.32)
with all other components being either zero, or implied by antisymmetry from the given ones. The conditions (4.30) and (4.31) can then be shown to be precisely equivalent to the conditions
$$\mathrm{\Gamma }_{12}\eta =\mathrm{\Gamma }_{34}\eta =\mathrm{\Gamma }_{56}\eta =\mathrm{}\mathrm{\Gamma }_{D1,D}\eta .$$
(4.33)
Since $`D`$ is even, and the total bundle space has dimension $`D+1`$, it follows that the spinors have $`2^{D/2}`$ components. There are $`\frac{1}{2}D1`$ equations in (4.33), each of which implies a halving of the original number of components, and so the final conclusion is that there are always 2 Killing spinors in these bundle spaces (real or complex, according to whether the spinors are Majorana or not). Special cases of this result that have appeared previously in the literature include the $`U(1)`$ bundles over $`S^2\times S^2`$, $`S^2\times S^2\times S^2`$ and $`CP^2\times S^2`$. We shall in general refer to all the $`q_i=k_i/\mathrm{}`$ Einstein spaces as “supersymmetric”spaces, although of course their Killing spinors are really only associated with supersymmetric compactifications in certain low-dimensional examples.
## 5 Consistency condition for Kaluza-Klein reductions
We saw in section 2 that in the cases of interest in supergravity reductions, a criterion for the consistency of the Kaluza-Klein reduction, and truncation to the massless gauge-boson sector, is that the Killing vectors $`\widehat{K}^I`$ associated with any gauge bosons that are to be retained must satisfy the condition that
$$Y(\widehat{K}^I,\widehat{K}^J)=\widehat{K}^{Im}\widehat{K}_m^J+\frac{1}{2m^2}(\widehat{}^m\widehat{K}^{In})(\widehat{}_m\widehat{K}_n^J)Y^{IJ}$$
(5.1)
should be constant, independent of the coordinates $`y`$ of the internal space. Here, $`m`$ is the related to the cosmological constant $`\widehat{\mathrm{\Lambda }}`$ of the internal Einstein space by $`\widehat{\mathrm{\Lambda }}=Dm^2`$, where the dimension of the internal space is $`D+1`$.
We begin by noting that the second term in (5.1) can be re-expressed more simply by using the following identity:
$`\widehat{\text{ }\text{ }}(\widehat{K}^m\widehat{L}_m)`$ $`=`$ $`\widehat{K}^m\widehat{\text{ }\text{ }}\widehat{L}_m+\widehat{L}^m\widehat{\text{ }\text{ }}\widehat{K}_m+2(\widehat{}^m\widehat{K}^n)(\widehat{}_m\widehat{L}_n)`$ (5.2)
$`=`$ $`2\widehat{\mathrm{\Lambda }}\widehat{K}^m\widehat{L}_m+2(\widehat{}^m\widehat{K}^n)(\widehat{}_m\widehat{L}_n),`$
for any pair of Killing vectors $`\widehat{K}^m`$ and $`\widehat{L}^m`$, where we have made use of the fact that Killing vectors on an Einstein space with cosmological constant $`\widehat{\mathrm{\Lambda }}`$ satisfy the equation $`\widehat{\text{ }\text{ }}\widehat{K}^m+\widehat{\mathrm{\Lambda }}\widehat{K}^m=0`$. This allows us to express the second term in (5.1) in terms of $`\widehat{K}^m\widehat{L}_m`$:
$$(\widehat{}^m\widehat{K}^n)(\widehat{}_m\widehat{L}_n)=\frac{1}{2}\text{ }\text{ }(\widehat{K}^m\widehat{L}_m)+\widehat{\mathrm{\Lambda }}\widehat{K}^m\widehat{L}_m.$$
(5.3)
Note that we just need the Laplacian $`\text{ }\text{ }`$ on the base space here, since it is equal to the Laplacian $`\widehat{\text{ }\text{ }}`$ in the bundle space when acting on scalars that are independent of the fibre coordinate $`z`$. The quantity $`Y(\widehat{K},\widehat{L})`$ defined in (5.1), whose constancy is need for consistency, is therefore expressible as
$$Y(\widehat{K},\widehat{L})=\frac{1}{2}(D+2)\widehat{K}^m\widehat{L}_m+\frac{D}{4\widehat{\mathrm{\Lambda }}}\text{ }\text{ }(\widehat{K}^m\widehat{L}_m).$$
(5.4)
We shall refer to the criterion that $`Y(\widehat{K},\widehat{L})`$ in (5.1) be constant as “The Consistency Condition” for short.
With our results from the previous sections we are now able to test the consistency condition in general, for any Einstein space $`\widehat{M}`$ that is constructed as a $`U(1)`$ bundle over a product of complex projective base spaces. Before doing so, we shall show that for any sphere $`S^n`$, with its standard round metric, all the $`SO(n+1)`$ Killing vectors satisfy the consistency condition. This is an important point not only for the discussion of Kaluza-Klein reductions on spheres themselves, but also we shall need to make use of this fact later in the section, when we examine the consistency condition in more general cases.
One way to prove that the full set of $`SO(n+1)`$ Killing vectors on the sphere $`S^n`$ satisfy the consistency condition is by using the fact that there are always Killing spinors on the sphere, equal in number to the dimension of the spinors, that satisfy
$$\widehat{}_m\eta ^A\frac{\mathrm{i}}{2}m\mathrm{\Gamma }_m\eta ^A=0.$$
(5.5)
From any pair of these, one can construct vectors $`\widehat{K}_m^{AB}=\overline{\eta }^A\mathrm{\Gamma }_m\eta ^B`$, which can easily be seen to satisfy the Killing vector equation. One can also show that all the Killing vectors of $`SO(n+1)`$ are obtained by this means. Furthermore, it follows from (5.5) that $`\widehat{}_m\widehat{K}_n^{AB}=\mathrm{i}m\overline{\eta }^A\mathrm{\Gamma }_{mn}\eta ^B`$. It is now relatively straightforward to show, using Fierz rearrangements, that the Killing vectors do indeed satisfy the consistency condition.
There is another way of showing that the full set of Killing vectors on the sphere satisfy the consistency condition, which is, perhaps, a little more geometrically appealing. We can describe the unit sphere $`S^n`$ as the surface $`x^Ax^A=1`$ in $`\text{I}\mathrm{R}^{n+1}`$, where $`x^A`$ are Cartesian coordinates in $`\text{I}\mathrm{R}^{n+1}`$. The Killing vectors on $`S^n`$ are then given by
$$K_{AB}=x^A_Bx^B_A.$$
(5.6)
If we write $`x^A=ru^A`$, where the $`u^A`$ satisfy $`u^Au^A=1`$ and are coordinates on the unit $`S^n`$, and $`r^2=x^Ax^A`$, then the metric on $`\text{I}\mathrm{R}^{n+1}`$ is given by
$$ds^2(\text{I}\mathrm{R}^{n+1})=dr^2+r^2du^Adu^A,$$
(5.7)
where $`du^Adu^A`$ is the metric on the unit $`S^n`$. If we denote by $`g_{AB}`$ the metric on the unit $`S^n`$, it is clear that it is related to the flat metric $`\delta _{AB}`$ on $`\text{I}\mathrm{R}^{n+1}`$ by
$$g_{AB}=\frac{1}{r^2}\left(\delta _{AB}\frac{x^Ax^B}{r^2}\right),$$
(5.8)
since this gives $`g_{AB}dx^Adx^B=du^Adu^A`$. An elementary calculation then shows that the inner product between Killing vectors $`K_{AB}`$ and $`K_{CD}`$ given in (5.6), with respect to the metric $`g_{AB}`$, is
$$(K_{AB}K_{CD})=\delta _{AC}u_Bu_D+\delta _{BD}u_Au_C\delta _{AD}u_Bu_C\delta _{BC}u_Au_D.$$
(5.9)
Now, the Laplacian on $`\text{I}\mathrm{R}^{n+1}`$ is related to the Laplacian on the unit $`S^n`$ by
$$\text{ }\text{ }_{R^{n+1}}=\frac{1}{r^n}\frac{}{r}\left(r^n\frac{}{r}\right)+\frac{1}{r^2}\text{ }\text{ }_{S^n}.$$
(5.10)
From (5.9), and $`x^A=ru^A`$, we shall have
$$\text{ }\text{ }_{R^{n+1}}\left(r^2(K_{AB}K_{CD})\right)=4(\delta _{AC}\delta _{BD}\delta _{AD}\delta _{BC}),$$
(5.11)
and hence using (5.10) we obtain
$$\text{ }\text{ }_{S^n}(K_{AB}K_{CD})+2(n+1)(K_{AB}K_{CD})=4(\delta _{AC}\delta _{BD}\delta _{AD}\delta _{BC}).$$
(5.12)
Since the unit $`S^n`$ has cosmological constant $`(n+1)`$, which corresponds to $`m^2=1`$ in (5.1), we finally arrive at the result that on the unit $`S^n`$
$$Y(K_{AB},K_{BC})=\delta _{AC}\delta _{BD}\delta _{AD}\delta _{BC}.$$
(5.13)
This shows that indeed all the $`SO(n+1)`$ Killing vectors on the sphere $`S^n`$ satisfy the consistency condition.
We now turn to the case where the internal manifold is a general Einstein space that can be constructed as a $`U(1)`$ bundle over a product of complex projective spaces, of the kind we have discussed in the previous sections. In section 3, we derived the expression (3.30) for a Killing vector on the bundle space, and (3.32) for its expression as a 1-form. It is now straightforward to calculate the inner product between any two Killing vectors, which we shall need for testing the consistency condition. Let us first establish the notation that we shall write the $`U(1)`$ Killing vector that generates translations along the fibres as
$$U\frac{}{z}.$$
(5.14)
It is easily seen that written as a 1-form, this is
$$U=c^2(dzA).$$
(5.15)
We shall use $`\widehat{K}_i`$ to denote a Killing vector lifted from the factor $`M_i`$ in the base manifold. There are four different sectors to consider in the consistency condition, namely $`Y(U,U)`$, $`Y(U,\widehat{K}_i)`$, $`Y(\widehat{K}_i,\widehat{K}_j)`$ (with $`ij)`$ and $`Y(\widehat{K}_i,\widehat{L}_i)`$ (where $`\widehat{K}_i`$ and $`\widehat{L}_i`$ are two Killing vectors in the same factor $`M_i`$ in the base manifold).
Taking $`Y(U,U)`$ first we see from (5.14) and (5.15) that $`U^mU_m=c^2=`$constant, and hence from (5.4) we shall have $`Y(U,U)=`$constant. So the $`U(1)`$ Killing vector by itself always satisfies the consistency condition.
Next, consider $`Y(U,\widehat{K}_i)`$. From (5.14) and (3.32) we have
$$U^m\widehat{K}_m^{(i)}=\alpha _ic^2\psi ^{(i)},$$
(5.16)
and so from (5.4) we obtain
$$Y(U,\widehat{K}_i)=\frac{\alpha _ic^2}{2\widehat{\mathrm{\Lambda }}}\left[(D+2)\widehat{\mathrm{\Lambda }}D\mathrm{\Lambda }_i\right]\psi ^{(i)}.$$
(5.17)
Since $`\psi ^{(i)}`$ is never constant (it satisfies $`\text{ }\text{ }\psi ^{(i)}=2\mathrm{\Lambda }_i\psi ^{(i)}`$), it follows that for a Killing vector $`\widehat{K}^{(i)}`$ coming from the base to be included in a consistent truncation as well as the $`U(1)`$ Killing vector $`U`$, the quantity in square brackets would have to vanish, i.e.
$$C_i(D+2)\widehat{\mathrm{\Lambda }}D\mathrm{\Lambda }_i=0.$$
(5.18)
We shall not analyse this condition extensively at this stage, since as we shall see later, more severe inconsistency problems generally occur in other sectors. We just note, however, that in view of the relation (3.40), consistency in this sector would require
$$\underset{j}{}d_j\mathrm{\Lambda }_jD\mathrm{\Lambda }_i=0.$$
(5.19)
In particular, this would be satisfied if all the $`\mathrm{\Lambda }_j`$ were equal, $`\mathrm{\Lambda }_j=\mathrm{\Lambda }`$, since $`_jd_j=D`$ (this is the case for all the spaces with $`q_i=k_i/\mathrm{}`$, i.e. the ones that admit 2 Killing spinors). However, we shall see below that the Killing vector $`\widehat{K}_i`$ will still run into other consistency problems in this case.
Moving on to the $`Y(\widehat{K}_i,\widehat{K}_j)`$ sector, where $`\widehat{K}_i`$ and $`\widehat{K}_j`$ come from different factors $`M_i`$ and $`M_j`$ in the base space, we find from (3.30) and (3.32) that the inner product for two such Killing vectors is
$$\widehat{K}_i^m\widehat{K}_{mj}=\alpha _i\alpha _jc^2\psi ^{(i)}\psi ^{(j)}.$$
(5.20)
Since the two functions $`\psi ^{(i)}`$ and $`\psi ^{(j)}`$ are assumed to live in two different factors in the base space here, it follows that $`^a\psi ^{(i)}_a\psi ^{(j)}=0`$, and hence, substituting into (5.4), we find
$$Y(\widehat{K}_i,\widehat{K}_j)=\frac{\alpha _i\alpha _jc^2}{2\widehat{\mathrm{\Lambda }}}\left[(D+2)\widehat{\mathrm{\Lambda }}D(\mathrm{\Lambda }_i+\mathrm{\Lambda }_j)\right]\psi ^{(i)}\psi ^{(j)}.$$
(5.21)
Again, since the $`\psi ^{(i)}`$ and $`\psi ^{(j)}`$ functions are always non-constant, the only way for $`Y(\widehat{K}_i,\widehat{K}_j)`$ to be constant would be if the quantity in square brackets vanished, namely
$$C_{ij}(D+2)\widehat{\mathrm{\Lambda }}D(\mathrm{\Lambda }_i+\mathrm{\Lambda }_j)=0.$$
(5.22)
Again, without fully analysing this condition here we may note that in the cases of principal interest with $`\mathrm{\Lambda }_k=\mathrm{\Lambda }`$ for all $`k`$ (the “supersymmetric” cases where there are 2 Killing spinors), equation (3.40) now allows us to deduce that
$$C_{ij}=D\mathrm{\Lambda },$$
(5.23)
and so the consistency condition is not satisfied. Thus we already see that we could not include Killing vectors from both of two factors $`M_i`$ and $`M_j`$ in the base space, at least in the supersymmetric cases where all the $`\mathrm{\Lambda }_k`$ are equal.
The fourth sector to consider is when two Killing vectors $`\widehat{K}`$ and $`\widehat{L}`$ come from the same factor $`M_i`$ in the base space. In order to avoid an unnecessary profusion of indices, we shall now suppress the “$`i`$” index that labels the particular factor in the product base manifold where the two Killing vectors are living. Thus the quantities $`\widehat{K}`$, $`\widehat{L}`$, $`\psi `$, $`\stackrel{~}{\psi }`$, $`d`$, $`\alpha `$, $`\mathrm{\Lambda }`$ in the following discussion all refer to this specific factor in the base space.
Now, the calculation of the inner product of the gives the result
$$\widehat{K}^m\widehat{L}_m=\alpha ^2c^2\psi \stackrel{~}{\psi }+^a\psi _a\stackrel{~}{\psi },$$
(5.24)
where $`K^a=J^{ab}_b\psi `$ and $`L^a=J^{ab}_b\stackrel{~}{\psi }`$. Substituting into $`Y`$ defined in (5.4), we now find
$`Y(\widehat{K},\widehat{L})=`$
$`{\displaystyle \frac{1}{2\widehat{\mathrm{\Lambda }}}}\{\alpha ^2c^2[(D+2)\widehat{\mathrm{\Lambda }}2D\mathrm{\Lambda }]\psi \stackrel{~}{\psi }+[(D+2)\widehat{\mathrm{\Lambda }}+D(\alpha ^2c^2\mathrm{\Lambda })]^a\psi _a\stackrel{~}{\psi }`$
$`+D(^a^b\psi )(_a_b\stackrel{~}{\psi })\}.`$ (5.25)
This equation can be simplified considerably, as follows. We may invoke the fact that if we consider the case where the base manifold has just a single factor $`M_i=CP^{n_i}`$, then the corresponding bundle space, with its Einstein metric, is the standard round metric on the sphere $`S^{2n_i+1}`$. Furthermore, we know that in this case all the Killing vectors on $`S^{2n_i+1}`$ satisfy the consistency condition, as we discussed earlier. This, therefore, allows us to deduce that the scalars $`\psi `$ and $`\stackrel{~}{\psi }`$ must satisfy equations such that (5.25) is constant when we take just the single factor $`M_i`$ in the base space. In this case we shall have $`D=d`$ (the dimension of the single space $`M_i`$). Substituting into (5.25), we then learn that
$$\frac{4\mathrm{\Lambda }^2}{d+2}\psi \stackrel{~}{\psi }+\frac{4\mathrm{\Lambda }}{d+2}^a\psi _a\stackrel{~}{\psi }+(^a^b\psi )(_a_b\stackrel{~}{\psi })$$
(5.26)
must be a constant, for any choice of $`\psi `$ and $`\stackrel{~}{\psi }`$ on $`M_i`$. This result<sup>8</sup><sup>8</sup>8One can also prove this result directly, as follows. We know that any Killing vector $`K^a`$ satisfies $`_a_bK_c=R^d{}_{abc}{}^{}K_{d}^{}`$. Since we have $`K^a=J^{ab}_b\psi `$ here, and furthermore the Riemann tensor on $`CP^n`$ is given by (4.22), we can conclude, after rescaling to cosmological constant $`\mathrm{\Lambda }`$ on $`CP^n`$, that
$$_a_b_c\psi =\frac{\mathrm{\Lambda }}{d+2}\left[J_{ab}J_{cd}^d\psi +J_{ac}J_{bd}^d\psi g_{ab}_c\psi g_{ac}_b\psi 2g_{bc}_a\psi \right],$$
(5.27) where $`d=2n`$. After some simple further manipulations, the constancy of (5.26) follows. for the eigenfunctions $`\psi `$ on $`CP^n`$ that they satisfy the condition that (5.26) is constant for any pair of such eigenfunctions. can now be fed back into (5.25) in the cases that really interest us, namely when there is more than one factor in the product base manifold. Specifically, we can use (5.26) in order to eliminate the $`(^a^b\psi )(_a_b\stackrel{~}{\psi })`$ terms in (5.25). Thus, we can deduce that consistency in this sector will be achieved only if
$$Q^a\psi _a\stackrel{~}{\psi }\beta \psi \stackrel{~}{\psi }$$
(5.28)
is constant, where the constant $`\beta `$ is given by
$$\beta =\frac{\frac{4D\mathrm{\Lambda }^2}{d+2}\alpha ^2c^2[2D\mathrm{\Lambda }(D+2)\widehat{\mathrm{\Lambda }}]}{\frac{4D\mathrm{\Lambda }}{d+2}(D+2)\widehat{\mathrm{\Lambda }}D(\alpha ^2c^2\mathrm{\Lambda })}.$$
(5.29)
Using (3.36) and (3.37), this can be rewritten as
$$\beta =\frac{4D\mathrm{\Lambda }^22(\mathrm{\Lambda }\widehat{\mathrm{\Lambda }})[2D\mathrm{\Lambda }(D+2)\widehat{\mathrm{\Lambda }}](d+2)}{4D\mathrm{\Lambda }(d+2)[(D+2)\widehat{\mathrm{\Lambda }}+D(\mathrm{\Lambda }2\widehat{\mathrm{\Lambda }})]}.$$
(5.30)
It is easiest to analyse this condition in the case where the Killing vector $`\widehat{L}`$ is taken to be the same as $`\widehat{K}`$, since if we can show that $`Y(\widehat{K},\widehat{K})`$ is not a constant, then that will show that no Killing vector from the base space can be retained in a consistent truncation. Let us therefore just consider one scalar eigenfunction $`\psi `$, with $`\stackrel{~}{\psi }=\psi `$. Thus we wish to study whether the quantity
$$Q^a\psi _a\psi \beta \psi ^2$$
(5.31)
can be constant. If $`Q`$ is constant then $`_aQ`$ will be zero, and so we can follow the familiar strategy of integrating $`(_aQ)^2`$ over the factor $`M_i=CP^{n_i}`$ in the product base manifold, where the scalar eigenfunction $`\psi `$ resides. If we can show that this integral is positive, then it will establish that $`Q`$ is not constant, and hence that the gauge boson associated to the corresponding Killing vector cannot be retained in a consistent Kaluza-Klein reduction.
Using integrations by parts, and the equation $`\text{ }\text{ }\psi =2\mathrm{\Lambda }\psi `$, repeatedly, we can establish the following results:
$`{\displaystyle \psi ^2|\psi |^2}`$ $`=`$ $`\frac{2}{3}\mathrm{\Lambda }{\displaystyle \psi ^4},`$
$`{\displaystyle _a|\psi |^2^a(\psi ^2)}`$ $`=`$ $`\frac{8}{3}\mathrm{\Lambda }^2{\displaystyle \psi ^4}2{\displaystyle |\psi |^4},`$
$`{\displaystyle _a|\psi |^2^a|\psi |^2}`$ $`=`$ $`\frac{8}{3}\mathrm{\Lambda }^3{\displaystyle \psi ^4}{\displaystyle \frac{2(d2)}{d+2}}\mathrm{\Lambda }{\displaystyle |\psi |^4},`$ (5.32)
where $`|\psi |^2_a\psi ^a\psi `$ and $`|\psi |^4(|\psi |^2)^2`$. (We have used the relation (5.27) in obtaining the last of these three equations.) Using these results, we find that
$$|Q|^2=\frac{8}{3}\mathrm{\Lambda }(\beta \mathrm{\Lambda })^2\psi ^4+2\left(2\beta \frac{d2}{d+2}\mathrm{\Lambda }\right)|\psi |^4.$$
(5.33)
Using this, it is possible to show that, except for “trivial” cases that we shall discuss below, the quantity $`Q`$ can never be constant for any of the eigenfunctions $`\psi `$ associated with the Killing vectors of the $`SU(n_i+1)`$ factors in the isometry group of the bundle space. We shall first discuss the “supersymmetric” cases, where the winding numbers $`q_i`$ satisfy $`q_i=k_i/\mathrm{}`$, since the proof is very simple in these cases, and furthermore they are the examples of principal physical interest. After that, we shall present a complete analysis for all possible choices of winding numbers.
As we saw in section 3, when $`q_i=k_i/\mathrm{}`$ the cosmological constants of all the $`CP^{n_i}`$ factors in the base space are equal, as are the constants $`\alpha _i`$; they are given by (3.44) and (3.45). Substituting these into (5.29) we find $`\beta =\mathrm{\Lambda }`$, and so (5.33) gives
$$|Q|^2=\frac{2\mathrm{\Lambda }(d+6)}{(d+2)}|\psi |^4.$$
(5.34)
The right-hand side is manifestly positive, and so the result that $`Q`$ cannot be constant follows.
For the general (non-supersymmetric) case with arbitrary winding numbers $`q_i`$, consider first the situation when the factor $`M_i`$ in the base space where $`\psi `$ resides is $`CP^1`$. In this particular case, because $`CP^1`$ is the sphere $`S^2`$, it follows that the three eigenfunctions $`\psi `$ that generate the $`SO(3)`$ Killing vectors actually satisfy the equation $`_a_b\psi =\mathrm{\Lambda }g_{ab}\psi `$, and from this it follows that on $`CP^1`$ we have
$$|\psi |^4=\frac{8}{3}\mathrm{\Lambda }^2\psi ^4.$$
(5.35)
Substituting this into (5.33) gives
$$|Q|^2=\frac{8}{3}\mathrm{\Lambda }(\beta +\mathrm{\Lambda })^2\psi ^4.$$
(5.36)
Thus we see that in this case it must be that $`Q`$ is constant if and only if $`\beta =\mathrm{\Lambda }`$. It is easy to see from the equations (3.35), (3.36), (3.37) and (5.29) that this can happen only in the extreme case where the fibres in the $`U(1)`$ bundle have a non-zero winding number only over the $`S^2`$ factor in the base space where $`\psi `$ resides. But in this extreme case the total space is simply the direct product of $`S^3`$ times the remaining $`CP^{n_i}`$ factors in the base. Not surprisingly, since $`S^3`$ is a group manifold, it has Killing vectors for which the associated quantity $`Q`$ will be constant. (Since any given Killing vector is associated with a left-translation or right-translation under $`SU(2)`$.) Aside from this extreme case, which is certainly not the one of interest to us in this paper, we see that $`Q`$ can never be constant.
Next, consider the case where the eigenfunction $`\psi `$ lives in a $`CP^2`$ factor in the base space. It is necessary, again, to determine the relation between $`|\psi |^4`$ and $`\psi ^4`$. Clearly this will be of the form
$$|\psi |^4=c\mathrm{\Lambda }^2\psi ^4,$$
(5.37)
where $`c`$ is a pure (dimensionless) number. It is evident from the expressions(4.18) or (4.19) for $`\psi `$ that the two integrals on $`CP^2`$ must be expressible in terms of $`SU(3)`$-invariant quartic polynomials built from the traceless Hermitean tensor $`T_A^B`$. Since there is no independent fourth-order Casimir for $`SU(3)`$, it must be that both integrals in (5.37) for $`CP^2`$ are pure numbers times $`(T_A{}_{}{}^{B}T_{B}^{}{}_{}{}^{A})^2`$, the numbers being independent of the choice of $`T_A^B`$. Thus the constant $`c`$ can be determined by evaluating the two sides of (5.37) for any convenient choice of eigenfunction $`\psi `$. From (4.19), a simple choice is to take the $`\psi `$ corresponding to $`T_\alpha {}_{}{}^{\beta }=\delta _\alpha ^\beta `$, which implies $`T_0{}_{}{}^{0}=2`$, with all other components of $`T_A^B`$ zero. This gives
$$\psi =13f^1,$$
(5.38)
where $`f`$ is given in (4.2). It is easy to substitute this into (5.37), leading to the result that
$$c=2.$$
(5.39)
Finally, using this result in (5.33), with $`n=\frac{1}{2}d`$, we arrive at the following:
$$|Q|^2=\frac{8}{3}\mathrm{\Lambda }\left[(\beta +\frac{1}{2}\mathrm{\Lambda })^2+\frac{1}{4}\mathrm{\Lambda }^2\right]\psi ^4,$$
(5.40)
which shows that $`Q`$ can never be constant in this case.
Finally, we can consider the general case where $`\psi `$ lives in a $`CP^n`$ factor in the base space. Now the calculation is a little more involved, since the ratio of $`\psi |^4`$ to $`\psi ^4`$ depends on the specific choice of eigenfunction $`\psi `$, when $`n3`$. In order to achieve the best chance of having $`Q|^2`$ be zero, one wants the ratio of $`|\psi |^4`$ to $`|\psi |^4`$ to be as large as possible, since then the (possibly negative) second term on the right-hand side of (5.33) has the best chance to outweigh the always-positive contribution from the first term on the right-hand side. In the Appendix we present some calculations that provide a determination of the largest value of this ratio; see (B.17) and (B.18). Thus from (5.33) we find that when $`n`$ is odd, we shall have
$$|Q|^2\frac{8}{3}\mathrm{\Lambda }\left(\beta +\frac{2\mathrm{\Lambda }}{n+1}\right)^2\psi ^4,$$
(5.41)
whilst when $`n`$ is even we instead find
$$|Q|^2\frac{8}{3}\mathrm{\Lambda }\left[\left(\beta +\frac{2n\mathrm{\Lambda }}{n^2+n+2}\right)^2+\frac{4(n+2)}{(n^2+n+2)^2}\right]\psi ^4,$$
(5.42)
From these results we see that $`Q`$ can never be constant when $`n`$ is even. When $`n`$ is odd instead, we see that $`Q`$ can be constant if and only if
$$\beta =\frac{2\mathrm{\Lambda }}{n+1}.$$
(5.43)
Now from (3.36) and (3.37) it immediately follows that if we define $`x\widehat{\mathrm{\Lambda }}/\mathrm{\Lambda }`$, then
$$\frac{d}{d+2}x1.$$
(5.44)
The lower limit is saturated if the $`U(1)`$ fibres wind only over the chosen $`CP^n`$ factor in the base space, whilst the upper limit is saturated if instead the $`U(1)`$ fibres have zero winding number over the chosen $`CP^n`$ factor. Using (5.30), we find that
$$\beta +\frac{4\mathrm{\Lambda }}{d+2}=\mathrm{\Lambda }\frac{2(D+2)(d+2)^2x^2+2(d+2)(3Dd+8D+2d)x4dD(d+4)}{(d+2)[(D2)(d+2)xD(d2)]}.$$
(5.45)
The denominator is positive for all $`x`$ in the interval (5.44), and the numerator has no extremum in this interval. It then follows that we shall have
$$\beta \frac{4\mathrm{\Lambda }}{d+2},$$
(5.46)
with equality being achieved only if $`x=d/(d+2)`$. Since $`d=2n`$, we conclude from this and (5.43) that $`Q`$ can be constant only in the extreme case where the $`U(1)`$ fibres wind purely over the $`CP^n`$ factor in the base space in which the eigenfunction $`\psi `$ resides.
The reason for the occurrence of these exceptional cases where $`Q`$ can be constant is the following. When $`n`$ is odd, say $`n=2q+1`$, and the fibres wind only over the $`CP^{2q+1}`$ factor in the base manifold, the total space is the direct product of $`S^{4q+3}`$ with the other $`CP^{n_i}`$ factors in the base space. Now the sphere $`S^{4q+3}`$ can be described as an $`SU(2)`$ bundle over the quaternionic projective space $`HP^q`$. Consequently, an $`SU(2)`$ subgroup of the $`SO(4q+4)`$ isometry group of the sphere corresponds to left translations by $`SU(2)`$ on the $`SU(2)`$ fibres, and therefore the associated $`SU(2)`$ Killing vectors $`K^I`$ will necessarily have the property that $`K^IK^J=`$constant, and so they will be associated with eigenfunctions $`\psi `$ on $`CP^{2q+1}`$ that satisfy the condition $`Q=`$constant. It is these Killing vectors that are being “detected” by the saturation of the bound (5.41).
These exceptional cases are higher-dimensional generalisations of the exception arising for $`n=1`$, with the fibres winding only over the $`CP^1`$ factor to give $`S^3`$, which we discussed previously. Again they are “trivial,” from the point of view of our analysis of compactifications, since we are not particularly interested in cases where the internal space is a direct product of a sphere $`S^{4q+3}`$ with a Kähler space. Nonetheless, it is reassuring to find that our rather intricate general analysis has indeed, as it should, detected these slightly obscure exceptions to the general rule.
With these results, we have proved that the non-abelian Killing vectors on the $`U(1)`$ bundle spaces over any product of $`CP^{n_i}`$ factors in the base space will never satisfy the consistency requirement that $`Y^{IJ}`$ in (5.1) is a constant.<sup>9</sup><sup>9</sup>9Except in the previously-discussed trivial cases of $`SU(2)`$ Killing vectors in the $`S^{4q+3}`$ factors in a bundle space where the fibres wind only over a $`CP^{2q+1}`$ base-space factor. This means that the associated Kaluza-Klein Yang-Mills fields associated with the non-abelian part of the symmetry group cannot be consistently retained in a massless truncation. In particular, this proves that of the $`U(1)\times SU(2)\times SU(2)`$ Yang-Mills fields in the $`Q(1,1)`$ compactification of the type IIB theory to $`D=5`$, the $`SU(2)\times SU(2)`$ fields cannot be retained in a consistent massless truncation.
## 6 Conclusions
In this paper, we have studied a necessary condition for the occurrence of a consistent Kaluza-Klein reduction on an internal Einstein manifold, in which all the Yang-Mills fields associated with the isometry group of the compactifying space are retained in a massless truncation. This condition, that the quantitiy $`Y^{IJ}`$ defined in (5.1) should be constant, is of rather general relevance in all the known non-trivial consistent Kaluza-Klein reductions. In particular, this consistency criterion is satsified by all the Killing vectors on a sphere, of arbitrary dimension. Our principal goal in this paper has been to show that the consistency criterion is never satisfied by the non-abelian $`SU(n_i+1)`$ Killing vectors in the isometry groups of the spaces $`Q_{n_1\mathrm{}n_N}^{q_1\mathrm{}q_N}`$, which are defined as $`U(1)`$ bundles over the product $`_iCP^{n_i}`$ of complex-projective spaces $`CP^{n_i}`$, with winding numbers $`q_i`$. In particular, this shows that space $`Q_{11}^{11}`$ (sometimes known as $`T^{11}`$), the $`U(1)`$ bundle over $`S^2\times S^2`$, does not allow a consistent Kaluza-Klein reduction of type IIB supergravity in which the non-abelian Yang-Mills fields of its $`SU(2)\times SU(2)\times U(1)`$ isometry group are retained in a massless truncation. Likewise, the compactifications of $`D=11`$ supergravity on the $`U(1)`$ bundles over $`S^2\times S^2\times S^2`$ and over $`CP^2\times S^2`$ do not allow the retention of the corresponding non-abelian Yang-Mills fields in massless truncations. These facts will be of relevance in the AdS/CFT correspondence , where it should turn out that certain correlation functions involving products of single massive operators with massless ones will correspondingly be non-zero in these cases (see, for example, ).
We have set our proof of the inconsistency of the full massless truncations in these cases in a more general context, in which we show in general that the non-abelian Killing vectors on the bundle spaces $`Q_{n_1\mathrm{}n_N}^{q_1\mathrm{}q_N}`$ do not satisfy the consistency criterion that all Killing vectors on all spheres satisfy. In order to show this, we have made a detailed analysis that should also be of more general utility. In particular, we studied the lifting of Killing vectors from an arbitrary base manifold to a $`U(1)`$ bundle over the base, and then we specialised to the case where the base is a Kähler -Einstein space, or a product of Kähler -Einstein spaces. In such cases, more complete results can be obtained, based on the fact that any Killing vector in the base can be expressed in terms of a certain eigenfunction of the scalar Laplacian.
We then turned to the cases of principal interest, where the base space is the product of complex projective spaces $`CP^{n_i}`$. We made a study of the Fubini-Study metrics, and in an appendix we obtained a rather useful iterative construction for real metrics on $`CP^n`$. We showed that all of the bundle spaces $`Q_{n_1\mathrm{}n_N}^{q_1\mathrm{}q_N}`$ can be given Einstein metrics, provided only that all the winding numbers $`q_i`$ do not vanish. We also showed that in the special case where $`q_i=(n_i+1)/\mathrm{}`$, where $`\mathrm{}`$ is the greatest common divisor of the $`(n_i+1)`$, the Einstein spaces all admit two Killing spinors. These cases, for $`Q_{11}^{11}`$, $`Q_{111}^{111}`$ and $`Q_{21}^{32}`$, are the ones of principal interest in the context of supergravity compactifications, since they imply the existence of unbroken supersymmetries.
We showed also that the question of whether the non-abelian Killing vectors of the $`U(1)\times _iSU(n_1+1)`$ isometry group of $`Q_{n_1\mathrm{}n_N}^{q_1\mathrm{}q_N}`$ satisfy the consistency criterion in (5.1) can be reduced to the question of whether the scalar eigenfunctions on $`CP^n`$ that are related to its Killing vectors satisfy certain integral bounds. We studied these bounds in detail, and used these to obtain our proofs of the inconsistency of the Kaluza-Klein reductions.
## Acknowledgements
We are grateful to Mirjam Cvetič, Hong Lü, Arta Sadrzadeh, Kelly Stelle and Tuan Tran for helpful discussions. C.N.P. thanks the Caltech-USC Center for Theoretical Physics for hospitality during the completion of this work.
Appendices
## Appendix A An iterative construction of $`CP^n`$
On occasion, it is helpful to have a real expression for the Fubini-Study metric on $`CP^n`$ available. This is easily done for low-dimensional examples by making specific adapted coordinate choices (see, for example, for an explicit real metric on $`CP^2`$). In general, we can give an elegant iterative construction for the metric on $`CP^n`$ in terms of the metric on $`CP^{n1}`$.
We take as our starting point the standard Fubini-Study metric (4.6) on $`CP^n`$, and write the inhomogeneous coordinates $`\zeta ^\alpha `$ as
$$\zeta ^\alpha =\mathrm{tan}\xi u^\alpha ,\text{with}u^\alpha \overline{u}^{\overline{\alpha }}=1.$$
(A.1)
With this coordinate redefinition the $`CP^n`$ metric (4.6) becomes
$$d\mathrm{\Sigma }_n^2=d\xi ^2+\mathrm{sin}^2\xi du^\alpha d\overline{u}^{\overline{\alpha }}\mathrm{sin}^4\xi |\overline{u}^{\overline{\alpha }}du^\alpha |^2.$$
(A.2)
Noting that the $`n`$ quantities $`u^\alpha `$ are themselves complex coordinates on $`C^n`$, subject to the constraint $`u^\alpha \overline{u}^{\overline{\alpha }}=1`$, we can follow the same strategy as in the original $`CP^n`$ construction, by introducing $`(n1)`$ inhomogeneous coordinates $`v^i`$, with $`1in1`$, defined by
$$v^i=\frac{u^i}{u^n},$$
(A.3)
where $`u^n`$ here denotes the $`n`$’th of the coordinates $`u^\alpha `$. In addition, we define
$$u^n=|u^n|e^{\mathrm{i}\stackrel{~}{\tau }}.$$
(A.4)
After a little calculation, we see that the metric (4.6) on $`CP^n`$ now takes the form
$$d\mathrm{\Sigma }_n^2=d\xi ^2+\mathrm{sin}^2\xi \mathrm{cos}^2\xi (d\stackrel{~}{\tau }+\stackrel{~}{B})^2+\mathrm{sin}^2\xi d\mathrm{\Sigma }_{n1}^2,$$
(A.5)
where $`d\mathrm{\Sigma }_{n1}^2`$ is the Fubini-Study metric on the unit $`CP^{n1}`$, and $`\stackrel{~}{B}`$ is a potential for the Kähler form of $`CP^{n1}`$:
$`d\mathrm{\Sigma }_{n1}^2`$ $`=`$ $`\stackrel{~}{f}^1dv^id\overline{v}^{\overline{i}}\stackrel{~}{f}^2|\overline{v}^{\overline{i}}dv^i|^2,\stackrel{~}{f}=1+v^i\overline{v}^{\overline{i}},`$
$`\stackrel{~}{B}`$ $`=`$ $`\frac{1}{2}\mathrm{i}\stackrel{~}{f}^1(v^id\overline{v}^{\overline{i}}\overline{v}^{\overline{i}}dv^i).`$ (A.6)
Thus (A.5) gives us an iterative construction of the Fubini-Study metric on the unit $`CP^n`$ in terms of the Fubini-Study metric on the unit $`CP^{n1}`$. (In fact the metric in $`CP^2`$ obtained in is precisely of this form, with the metric on $`CP^1`$ being the standard metric on the 2-sphere.) Note that the potential $`B`$ for $`CP^n`$, appearing in (4.8), is given in terms of the analogous potential $`\stackrel{~}{B}`$ for $`CP^{n1}`$ by
$$B=\mathrm{sin}^2\xi (d\stackrel{~}{\tau }+\stackrel{~}{B}).$$
(A.7)
The function $`f`$ appearing in (4.2) is given by
$$f=\mathrm{sec}^2\xi .$$
(A.8)
## Appendix B Inequalities on $`CP^n`$
In section 5, we show that the gauge boson associated with any Killing vector on a factor $`CP^n`$ in the base manifold whose associated scalar harmonic $`\psi `$ has a $`Q`$, as defined in (5.31), that is non-constant, cannot be retained in a consistent massless Kaluza-Klein reduction. In this appendix we derive some inequalities involving the integrals $`\psi ^4`$ and $`|\psi |^4`$ appearing in (5.33), which are used in the calculations in section 5.
In terms of the construction (4.18) or (4.19) for the eigenfunctions $`\psi `$, it is clear that the integrals $`\psi ^4`$ and $`|\psi |^4`$ must necessarily give rise to quartic $`SU(n+1)`$ invariants built from the traceless Hermitean tensor $`T_A^B`$. If we define
$$I_2T_A{}_{}{}^{B}T_{B}^{}{}_{}{}^{A},I_4T_A{}_{}{}^{B}T_{B}^{}{}_{}{}^{C}T_{C}^{}{}_{}{}^{D}T_{D}^{}{}_{}{}^{A},$$
(B.1)
it follows therefore that on $`CP^n`$ we must have
$$\psi ^4=a(I_2)^2+bI_4,|\psi |^4=\stackrel{~}{a}(I_2)^2+\stackrel{~}{b}I_4,$$
(B.2)
for pure numbers $`a`$, $`b`$, $`\stackrel{~}{a}`$ and $`\stackrel{~}{b}`$ that are dependent only on the value of $`n`$. In order to determine these constants, it suffices to consider just two special cases of eigenfunctions $`\psi `$ that have different values for the ratio $`I_4/(I_2)^2`$.
A convenient choice for the two eigenfunctions $`\psi _1`$ and $`\psi _2`$ is as follows. For $`\psi _1`$, we take $`T_\alpha {}_{}{}^{\beta }=\delta _\alpha ^\beta `$, $`T_0{}_{}{}^{0}=n`$, with all other components of $`T_A_B`$ vanishing. For $`\psi _2`$, we take instead $`T_n{}_{}{}^{0}=T_0{}_{}{}^{n}=\frac{1}{2}`$, with all other components vanishing. (Here “$`n`$” indicates that $`\alpha `$ takes the value $`\alpha =n`$.) For these two special cases the invariants $`I_2`$ and $`I_4`$ are given by:
$`\psi _1:`$ $`I_2=n(n+1),I_4=n(n^3+1),`$
$`\psi _2:`$ $`I_2=\frac{1}{2},I_4=\frac{1}{8}.`$ (B.3)
Thus when $`n2`$, we see that $`I_4/(I_2)^2`$ is different for the two eigenfunctions, and so by evaluating the integrals in (B.2), we shall be able to determine $`a`$, $`b`$, $`\stackrel{~}{a}`$ and $`\stackrel{~}{b}`$.
In order to evaluate the integrals, it is convenient to make use of the iterative construction of $`CP^n`$ metrics that we obtained in Appendix A. Specifically, we iterate twice, to give
$$d\mathrm{\Sigma }_n^2=d\xi ^2+\mathrm{sin}^2\xi \mathrm{cos}^2\xi (d\stackrel{~}{\tau }+\stackrel{~}{B})^2+\mathrm{sin}^2\xi \left(d\lambda ^2+\mathrm{sin}^2\lambda \mathrm{cos}^2\lambda (dz+C)^2+\mathrm{sin}^2\lambda d\mathrm{\Sigma }_{n2}^2\right).$$
(B.4)
(Our notation should be self-evident, by comparing with the construction in Appendix A.) The two eigenfunctions $`\psi _1`$ and $`\psi _2`$ can then be seen to be given by
$$\psi _1=1(n+1)\mathrm{cos}^2\xi ,\psi _2=\mathrm{sin}\xi \mathrm{cos}\xi \mathrm{cos}\lambda \mathrm{cos}\stackrel{~}{\tau }.$$
(B.5)
Other relevant points are that the determinant of the metric (B.4) is given by
$$\sqrt{g_n}=(\mathrm{sin}\xi )^{2n1}\mathrm{cos}\xi (\mathrm{sin}\lambda )^{2n3}\mathrm{cos}\lambda \sqrt{g_{n2}},$$
(B.6)
where $`g_{n2}`$ is the determinant of the metric $`d\mathrm{\Sigma }_{n2}^2`$ on $`CP^{n2}`$. Furthermore, the relevant components of the inverse metric are given by
$$g^{\xi \xi }=1,g^{\lambda \lambda }=\frac{1}{\mathrm{sin}^2\xi },g^{\stackrel{~}{\tau }\stackrel{~}{\tau }}=\frac{\mathrm{sec}^2\xi +\mathrm{tan}^2\lambda }{\mathrm{sin}^2\xi }.$$
(B.7)
For functions $`\varphi `$ of $`\xi `$, $`\lambda `$ and $`\stackrel{~}{\tau }`$ only, we have
$$|\varphi |^2=\left(\frac{\varphi }{\xi }\right)^2+\frac{1}{\mathrm{sin}^2\xi }\left(\frac{\varphi }{\lambda }\right)^2+\frac{\mathrm{sec}^2\xi +\mathrm{tan}^2\lambda }{\mathrm{sin}^2\xi }\left(\frac{\varphi }{\stackrel{~}{\tau }}\right)^2.$$
(B.8)
As a final preliminary, we note from (4.4) that since the unit sphere $`S^{2n+1}`$ has volume $`\mathrm{\Omega }_{2n+1}=2\pi ^{n+1}/\mathrm{\Gamma }(n+1)`$, it follows that the unit $`CP^n`$ has volume $`\mathrm{\Sigma }_n`$ given by
$$\mathrm{\Sigma }_n=\frac{\pi ^n}{\mathrm{\Gamma }(n+1)}.$$
(B.9)
It is now straightforward to evaluate all the necessary integrals, and thus to determine the constants $`a`$, $`b`$, $`\stackrel{~}{a}`$ and $`\stackrel{~}{b}`$ appearing in (B.2). We find that for the unit $`CP^n`$, with $`n2`$, we shall have
$`{\displaystyle \psi ^4}`$ $`=`$ $`{\displaystyle \frac{3\pi ^{n1}}{2\mathrm{\Gamma }(n+5)}}\left[(I_2)^2+2I_4\right],`$
$`{\displaystyle |\psi |^4}`$ $`=`$ $`{\displaystyle \frac{8\pi ^{n1}}{\mathrm{\Gamma }(n+5)}}\left[(n^2+5n+7)(I_2)^2+(n+1)(n+2)I_4\right].`$ (B.10)
For the discussion in section 5, it turns out that we need to know the largest possible value that the ratio $`(|\psi |^4)/(\psi ^4)`$ can attain. It is easy to see from (B.10) that this will occur for a tensor $`T_A^B`$ that gives the smallest possible value of $`I_4/(I_2)^2`$. To determine this value, let the (real) eigenvalues of $`T_A^B`$ be $`\lambda _A`$. Tracelessness implies that $`_A\lambda _A=0`$. If we solve for the eigenvalue $`\lambda _0`$ in terms of the $`\lambda _\alpha `$ for $`1\alpha n`$, we shall therefore have
$$I_2=\underset{\alpha }{}\lambda _\alpha ^2+(\underset{\alpha }{}\lambda _\alpha )^2,I_4=\underset{\alpha }{}\lambda _\alpha ^4+(\underset{\alpha }{}\lambda _\alpha )^4,$$
(B.11)
and so the ratio $`RI_4/(I_2)^2`$ is extremised when the $`\lambda _\alpha `$ satisfy
$$\lambda _\alpha ^3I_2\lambda _\alpha I_4+(\underset{\beta }{}\lambda _\beta )^3I_2\underset{\beta }{}\lambda _\beta I_4=0.$$
(B.12)
If we suppose that two of the extremising eigenvalues, say $`\lambda _\alpha `$ and $`\lambda _\beta `$, are unequal, then by subtracting their equations (B.12) we find that
$$\lambda _\alpha ^2+\lambda _\alpha \lambda _\beta +\lambda _\beta ^2=\frac{I_4}{I_2}.$$
(B.13)
If a third eigenvalue, say $`\lambda _\gamma `$, is unequal to both $`\lambda _\alpha `$ and $`\lambda _\beta `$, then by subtractions we can see that
$$\lambda _\alpha +\lambda _\beta +\lambda _\gamma =0.$$
(B.14)
Finally, if we suppose that a fourth eigenvalue $`\lambda _\delta `$ is unequal to all of the previous three, then by subtractions we arrive at the contradiction that $`\lambda _\delta =\lambda _\alpha `$. Therefore any set of $`\lambda _\alpha `$ that extremises the ratio $`R`$ can involve at most three different values.
It now becomes rather straightforward to find the extrema, and in particular, to identify the global minima. There are two distinct cases, depending upon whether $`n`$ is even or odd. We find that the minimum is achieved for
$`n=2q+1:`$ $`\lambda _0=\lambda _1=\mathrm{}=\lambda _q=\lambda ,\lambda _{q+1}=\lambda _{q+2}=\mathrm{}=\lambda _{2q+1}=\lambda ,`$ (B.15)
$`n=2q:`$ $`\lambda _0=\lambda _1=\mathrm{}=\lambda _q=\lambda ,\lambda _{q+1}=\lambda _{q+2}=\mathrm{}=\lambda _{2q}={\displaystyle \frac{n+2}{n}}\lambda .`$
(Of course in each case there are symmetry-related minima corresponding to permuting the eigenvalues. The minimisation occurs when the set of eigenvalues $`\lambda _i`$ divides into two subsets that are as nearly as possible equal in size, within each of which all eigenvalues are equal. This 50/50 partitioning is exact only if $`n`$ is odd, since the total number of eigenvalues $`\lambda _i`$ is then even.) Thus we find that the following inequalities hold:
$`n=2q+1:`$ $`{\displaystyle \frac{I_4}{I_2^2}}{\displaystyle \frac{1}{n+1}},`$
$`n=2q:`$ $`{\displaystyle \frac{I_4}{I_2^2}}{\displaystyle \frac{n^2+2n+4}{n(n+1)(n+2)}}.`$ (B.16)
Substituting into (B.10), and reinstating the cosmological constant $`\mathrm{\Lambda }`$ by the appropriate constant rescaling, we then find that
$$|\psi |^4c_{\mathrm{max}}\psi ^4,$$
(B.17)
with
$`n=2q+1:`$ $`c_{\mathrm{max}}={\displaystyle \frac{4(n+3)\mathrm{\Lambda }^2}{3(n+1)}},`$
$`n=2q:`$ $`c_{\mathrm{max}}={\displaystyle \frac{4(n+1)(n+2)\mathrm{\Lambda }^2}{3(n^2+n+2)}}.`$ (B.18)
These inequalities are used in section 5.
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# Space-borne global astrometric surveys: the hunt for extra-solar planets.
## 1 Introduction
At the very end of 1995, the discovery \[Mayor & Queloz 1995\] of the first Jupiter-mass ($`M_J`$) planet orbiting a normal star other than the Sun heralded the beginning of a new era of extraordinary discoveries in the realm of extra-solar planets, bringing with them the hope for a better understanding of the formation and frequency of planetary systems, and perhaps of bringing us closer to the ultimate goal of discovering extraterrestrial life.
After four years since that discovery, spectroscopic programs have been able to reveal some twenty extra-solar planets, i.e., objects with a lower mass limit below the 13-$`M_J`$ cut-off which has been adopted by Oppenheimer and Kulkarni \[Oppenheimer & Kulkarni 1999\] to differentiate giant planets from brown dwarfs.
However, these discoveries have raised new and troubling questions in our understanding of the properties of planetary systems. The fundamental tenets upon which present theories are based include nearly circular orbits and giant planets formed several AU from the central star, in contrast with the very short orbital periods \[Mayor & Queloz 1995\] and high eccentricities \[Latham et al. 1989, Cochran et al. 1996, Mazeh et al. 1996\] found for several of the new discoveries. Their interpretation as bona-fide planets rests on our understanding of correlations shown by their orbital and physical parameters, as recently discussed by Black \[Black 1997\] and earlier by Duquennoy & Mayor \[Duquennoy & Mayor 1991\] in their work on solar-type binary stars.
New models, which employ specific physical and dynamical mechanisms like in-situ formation \[Bodenheimer et al. 1999\] or orbital migration \[Lin et al. 1996, Trilling et al. 1998, Murray et al. 1998\], have been proposed to justify the presence of hot jupiters around normal stars, demonstrating that the interplay between additional theoretical work and more observational data will be necessary for a continued improvement in our theoretical understanding of how planets form and evolve, and where Earth-like planets could eventually be found.
However, simply adding a few tens of additional detections of giant extra-solar planets is not enough. A better understanding of the conditions under which planetary systems form and of their general properties requires large, complete samples of planets, with useful upper limits on Jupiter-mass planets at several AU from the central star.
Ongoing and planned radial velocity surveys \[Mayor & Queloz 1995, Cochran et al. 1996, Noyes et al. 1997, Marcy & Butler 1992, Marcy & Butler 1998\] have started filling significant portions of the relevant parameter space. Searches based on relative astrometry from the ground and in space \[Gatewood 1987, Colavita et al. 1999, Colavita et al. 1998, Mariotti et al. 1998, Pravdo & Shaklan 1996\] will be an important complement to the spectroscopic work and, probably, the preferred means for establishing the existence of planets around young stars and that of low mass planets down to a few Earth masses, as will be the case for SIM \[Boden et al. 1997, Unwin 1999\].
A HIPPARCOS-like, space-borne global astrometric mission, which can survey the whole sky to faint magnitudes and with high astrometric accuracy, will enable the monitoring of large ($`>10^5`$) samples of stars, with well understood completeness properties. This, depending on actual values of planetary frequencies \[Marcy et al. 1999\], could yield the possibility of making firm measurements of statistical properties of planetary systems. For, different correlations among orbital parameters (eccentricity, period or semi-major axis) and measurable differences in planetary frequency are likely to be generated by diverse planetary formation scenarios (core accretion and disk instability are the two known to date) and evolution mechanisms, as well as different formation and evolution processes of the parent star (binarity, spectral type, metallicity, age). An astrometric mission such as GAIA appears well poised for such a systematic census of planetary systems within $``$ 200 pc from the Sun.
The GAIA concept was originally proposed by Lindegren and Perryman \[Lindegren & Perryman 1996\] as a possible Cornerstone–class mission within the Horizon 2000+ program of scientific satellites of the European Space Agency. This satellite is designed to chart more than one billion objects (stars, extra-galactic objects, and solar system objects) on the sky down to the limiting magnitude of $`I=20`$. The targeted final accuracy is $``$ 10 $`\mu `$as on positions and parallaxes, and $``$ 10 $`\mu \mathrm{as}`$/year on proper motions at the reference magnitude of $`V=15`$ for a G2V star \[Gilmore et al. 1998\], and for a mission life time of 5 years.
In the following sections we show and discuss some relevant results derived from detailed end-to-end simulations of the data acquisition and analysis process for GAIA, which, as we will see, appears capable of discovering Jupiter-mass planets around $``$ 3$`\times 10^5`$ candidate stars (including dwarfs earlier than K5).
## 2 Data simulation
The simulation code is an adaptation of that used by Galligani et al. \[Galligani et al. 1989\] for the assessment of the astrometric accuracy of the sphere reconstruction in the HIPPARCOS mission. We generate catalogs of single stars randomly distributed on the sky; each run produces a sphere of 200 stars. As the satellite observing strategy (or scanning law) is most easily described in ecliptic coordinates, positions, proper motions and parallax ($`\lambda `$, $`\beta `$, $`\mu _\lambda `$, $`\mu _\beta `$, and $`\pi `$, respectively) are also given in the same coordinate system. For the moment, $`\mu _\lambda `$, $`\mu _\beta `$ and $`\pi `$, as well as magnitudes and colors, are drawn from simple distributions which do not represent any specific Galaxy model; in particular, in each run the 200 stars simulated have the same parallax, total proper motions, magnitudes, and colors.
The satellite is made to sweep the sky in such a way that the spin axis precesses around the Sun at a rate of about 6.5 rev/year and with a constant angle of $`43\mathrm{deg}`$. Stars that at any given time are “seen” within a strip $`1\mathrm{deg}`$ wide along the great circle (GC) being scanned are considered observed; a GC is completed in about 2.5 hours. Basic observations, which in principle can easily be derived from the fringe pattern measurements above, are the abscissae along a GC, as GAIA, like HIPPARCOS, makes very precise measurements in one dimension only. The mission lifetime is set to 5 years and the scanning law is such that the number of basic observations per star is a function of ecliptic latitude: each given object on the sky is re-observed approximately every month, for a total of $`60`$ one-dimensional position measurements, on average. The minimum number of observations is $``$ 30 and occurs for stars at the ecliptic equator. The position of a star at a given time, as described by the combined effects of parallax and proper motions, is called here barycentric location and it has been described in Euclidean space (flat Lorentzian). General relativistic effects, which will have to be considered in the future (especially for the case of Earth-like planets), are not taken into account.
Finally, gravitational perturbations (Keplerian motions) induced by a (single) nearby orbiting mass are added to the barycentric location resulting in the “true” geometric location of a target.
Simulated observations are generated by adding the appropriate astrometric noise to the true locations. The error sources considered in the simulation are discussed below.
As in HIPPARCOS, GAIA will have two viewing directions separated by a large angle named Base Angle (BA). It is by connecting directions far apart on the sky that the principle of space-borne global astrometry demonstrated by HIPPARCOS is implemented. Therefore, it is not surprising that the accuracy with which the BA is known throughout the mission is probably the most important single item within the GAIA concept. At the moment there are two mature optical designs for the BA. One configuration feature two telescopes mounted on the same optical bench and pointing along the two different line-of-sights. The other design uses a beam combiner (an adaptation of that used on HIPPARCOS) which physically materializes the BA and multiplexes the two viewing directions into a single telescope unit; beam combiner and telescope are bolted to the same bench. The twin telescopes designed for the first configuration feature an off-axis monolithic primary (with no central obscuration), while the collecting telescope of the latter design is a Fizeau interferometer with a baseline of 2.45m.
The details of the GAIA optical configuration (see e.g. Gilmore et al. \[Gilmore et al. 1998\], and references therein) are not critical to our discussion of photon-driven astrometric errors. The major difference with our earlier work \[Casertano et al. 1995\] is that both designs feature a significantly increased telescope aperture. The monolithic configuration has a rectangular primary of 1.7m by 0.7m, and each circular aperture of the interferometric option is 0.65m in diameter. High-accuracy measurements of the position (phase) of the PSFs are made directly on the focal plane using CCD detectors (see next section).
### 2.1 Payload configurations
The payload design has greatly advanced since the idea sketched in Lindegren & Perryman \[Lindegren & Perryman 1996\] and further developed in Loiseau & Shaklan \[Loiseau & Shaklan 1996\]. The two industrial studies commissioned by ESA have looked into the feasibility of two different options for the GAIA payload \[Gilmore et al. 1998\]: an all-passive configuration with two identical telescopes (with rectangular-shape monolithic primary mirrors) to be operated in L2, and an all-active configuration with an interferometric (diluted) beam combiner and a Fizeau interferometer as light collector behind it to be operated in geosynchronous orbit. Stability of the BA is passively maintained in the case of the monolithic configuration by utilizing a sophisticated active thermal control system (which must operate at the $`\mu `$K level) and a silicon carbide structure for the optical bench. On the other hand, the interferometric configuration is all-active, with the stability achieved by real-time monitoring implemented via high-accuracy laser metrology \[Gai et al. 1997\] and control of all critical degrees of freedom.
As we are more familiar with the interferometric option for GAIA \[Lattanzi et al. 1997\], we will be referring to that in the discussion below. However, precision and accuracy requirements and estimates are quite similar for both configurations, especially at the bright magnitudes we are concerned with in this paper.
### 2.2 Photon noise and point spread function measurement
The ability to measure accurately the position of a star ultimately depends on the width and shape of the point spread function (PSF) of the imaging system and on the number $`N`$ of photons detected. For a well-sampled PSF, the theoretical limit is shown by \[Lindegren 1978\] to be $`ϵ=\lambda /(4\pi x_{\mathrm{rms}}\sqrt{N})`$, where $`x_{\mathrm{rms}}`$ is the rms size of the aperture in the measurement direction. For two circular apertures of diameter $`D`$ and with a central separation $`B`$, we have $`x_{\mathrm{rms}}=\sqrt{(B/2)^2+(D/4)^2}`$; for the baseline GAIA parameters ($`B=2.45`$ m, $`D=0.65`$ m), $`x_{\mathrm{rms}}=1.24`$ m \[Lindegren & Perryman 1995\]. For $`\lambda =750`$ nm (the baselined operational wavelength), this translates into a theoretical monochromatic measurement accuracy of $`10\mathrm{mas}/\sqrt{N}`$.
The measurement accuracy is degraded for non-monochromatic measurements, by about a factor 2 for a Gaussian filter centered at 7500 Å and 2000 Å wide \[Gai et al. 1998\]. In addition, the requirement of optimal sampling may be difficult to achieve, since the central fringe is only about $`0\stackrel{}{.}04`$ wide and the field of view is $`1\mathrm{deg}`$. In practice, this will probably cause a loss of accuracy of about 20–40 per cent \[Gai et al. 1995, Gai et al. 1998\]. In the following, we assume a “best reasonable” single-measurement photon noise error of $`24\mathrm{mas}/\sqrt{N}`$.
Since scans overlap partially, each “observation” of a star will consist of about 8 consecutive scans with 20 sec of integration time allocated per scan, a total integration time of 160 s. Assuming a total collecting area of 0.664 $`\mathrm{m}^2`$ (2 apertures of 0.65 m diameter) and a total system efficiency of 20 per cent, a star with $`V=15`$ would generate about $`10^6`$ photons per observation, corresponding to a photon-limited measurement accuracy of the fringe position of 24 $`\mu \mathrm{as}`$. The accuracy scales with the inverse square root of the flux. This approximate calculation agrees with the values in Table 1, obtained for the current version of the interferometric option. The photon noise values (second column in Table 1) for $`10V18`$ stars were obtained with 3.1 pixels per fringe period, a RON of 3 electrons/pixels, and a DQE of 0.6. Objects brighter than $`V`$ = 10 mag will be also observed by GAIA. However, limitations on the dynamic range achievable with CCDs, requires to limit the exposure time for the brighter stars. Thus is practice, the precision starts to level off brighter than $`V`$ 10.
This accuracy is based on the photon statistics only, and does not take into account possible systematic effects, such as distortions in the optical system or in the detector, imperfection in the fringes, and so on. Many such systematic effects can be calibrated using closure methods, others will require enhancements in system design.
### 2.3 Accuracy of the Base Angle
The systematic effects mentioned above will lead to residual (systematic) errors in the knowledge of the BA which do not scale with magnitude. Ultimately, these will be the errors which will limit the maximum accuracy of GAIA for bright sources. We combine such residual errors in what is hereafter called residual bias, or simply bias. In the present error budget we assume that the bias can be described as a stationary stochastic process; therefore the bias contributed at the single-observation level to each object (Table 1) scales, like photon noise precision, with the average number of observations collected over the mission lifetime.
The BA can be measured and monitored accurately over time intervals longer than a full revolution ($``$ 3 hours) by making use of the $`2\pi `$ closure properties of the consecutive scans. This is an especially important feature of GAIA whose potential, at the $`mas`$ level, has been beautifully proven by HIPPARCOS. However, the bias over shorter time scales need to be controlled by ensuring that the relative positions and shapes of all optical elements of the beam combiner do not vary throughout the observation. The necessary sub-nm accuracy will probably be achieved by a combination of passive control and of laser metrology. For simplicity, we summarize the overall accuracy (bias) with which the (relative) positions of the mirrors of the beam combiner must be known by a single number, the “baseline error” $`\sigma _\mathrm{b}`$. This helps visualize the complex interplay of the beam combiner mirrors (which materialize the two baselines of the interferometric design) with a more familiar bias of the type $`B\times \delta `$, where B is the baseline and $`\delta `$ the angular uncertainty<sup>1</sup><sup>1</sup>1While the main contribution to the baseline error will probably come from the relative position of the beam combiner mirrors, the motions and/or distortions of other optical elements (such as the primary apertures of the Fizeau interferometer) can also generate an effective baseline error.. This also helps understand the bias contribution in Table 1. For example, if the photon error is subtracted (in quadrature) from the total error for the case with $`\sigma _\mathrm{b}=30`$ pm, the resulting value, $`10`$ $`\mu \mathrm{as}`$, represents the angular bias corresponding to that linear “baseline error”. Therefore, the requirement on the metrology is $``$ 3.8 times more stringent that one would have guessed from the simple calculation $`B\times \delta `$, with B=2.45 m and $`\delta `$=10 $`\mu \mathrm{as}`$. This example shows that the baseline bias is driven by the dimension of the single mirrors (D=0.65 m) forming the baseline (B/D$``$3.8) within the beam combination system.
Laser metrology in laboratory settings has already achieved extremely high performances, with relative measurements at the picometer level \[Gürsel 1993, Noecker et al. 1993, Noecker 1995, Reasenberg et al. 1995\]. Such precision has been reached over short (few wavelengths) variations in path length, which are appropriate to the GAIA design if a good active thermal control is included.
However, the few pm error quoted refers to the precision and stability of a one-dimensional laser gauge measurement of a single optical path. Maintaining the accuracy of the interferometer baseline is much more complex, first, because the three-dimensional positions of several optical elements may need to be monitored simultaneously, and second, because of the possible differences between the optical path of starlight and of the laser gauge beams. Noecker \[Noecker 1995\] lists a number of possible systematic errors for the POINTS mission concept. An experiment to demonstrate picometer laser metrology for a three-dimensional system on the GAIA scale is underway \[Gai et al. 1997\]. For the moment, it remains difficult to give firm figures for the baseline accuracy that will eventually be achieved. In Table 1, we consider two cases which probably bracket realistic expectations: a more “conservative” accuracy $`\sigma _\mathrm{b}=30`$ pm, which already appears within reach from the preliminary results available, and an “optimistic” accuracy $`\sigma _\mathrm{b}=10`$ pm. Notice that these bias values are intended at the level of what is called here the single-observation error. The accuracy levels at the end of the mission will be $`3`$ $`\mu \mathrm{as}`$ and 1 $`\mu \mathrm{as}`$ for $`\sigma _\mathrm{b}=30`$ pm and 10 pm, respectively, if errors in successive passes are uncorrelated, as discussed previously. Notice that, for $`\sigma _\mathrm{b}=30`$ pm, the astrometric precision begins to level off brighter than $`V=12`$.
Finally, it is important to mention that the numbers quoted for the linear bias apply to a 3-dimensional monitoring of the relevant structure. The laser beams will have to monitor the corresponding one-dimensional degrees of freedom with significantly better resolution—a factor $`<2`$ for the current design.
## 3 DETECTION AND ORBIT DETERMINATION METHODS
The magnitude of the gravitational perturbations induced by a planet on the parent star, as seen by an astrometric mission, can be quantified through its astrometric signature $`\alpha `$ defined as:
$$\alpha =\frac{M_p}{M_s}\frac{a_p}{d}$$
(1)
where $`M_p`$, $`M_s`$ are the masses of the planet and star respectively, $`a_p`$ the semi-major axis of the planetary orbit, $`d`$ the distance of the system from us. If $`a_p`$ is in AU and $`d`$ in parsec, then $`\alpha `$ is expressed in arcsec.
GAIA’s sensitivity to this signal depends of course on the errors $`\sigma _\psi `$ of each measurement, with theoretical values listed in Table 1. We have verified through tests with different values of $`\sigma _\psi `$ (see section 4.1) that, as could be expected, the detection probability depends on $`\sigma _\psi `$ and $`\alpha `$ only through their ratio, the “astrometric” signal-to-noise ratio
$$S/N=\alpha /\sigma _\psi ,$$
(2)
so that the results obtained can be straightforwardly rescaled to different measurement errors.
For simplicity, we have thus kept the single-observation error fixed at $`\sigma _\psi =10\mu \mathrm{as}`$ throughout our simulations. This is the value expected for relatively bright stars ($`V`$ 12 mag, corresponding to the Sun at 200 pc), with a conservative baseline error of 30 pm (see Table 1). Our simulations should thus give realistic estimates of: a) GAIA’s detectability horizon of planetary mass companions to solar-type stars in the vicinity of our Sun, and b) limits on distance for accurate orbital parameters determination.
The starting point for our two-level statistical investigation is the computation of detection probabilities, which will in principle depend upon 1) mission parameters, 2) noise sources, and 3) orbital elements. The contributions to points 1) and 2) are listed in Section 2. As for point 3), we will express the detection probabilities as function of the period $`P`$ and the signature $`\alpha `$, which we expect to be the major contributors, and generally average over the expected distribution of the other orbital parameters.
### 3.1 Detection
Our first detection method applies a standard $`\chi ^2`$ test (with the confidence level set to 95 per cent) to the residuals $`\psi \psi _r`$, where the $`\psi `$ are the actual measurements, and the $`\psi _r`$ are the GC abscissae recomputed after fitting the observations of each star with a single-star model, i.e., solving only for the five astrometric parameters appropriate for a single star (position, parallax, and proper motion). Being $`\chi _o^2`$ the value provided by the null model (no planet), the test fails when Pr$`(\chi ^2\chi _o^2)0.95`$. In this case the planet is considered detected. Note that this method only measures deviations from the single-star model, it makes no assumptions on the nature of the residuals nor does it provide initial guesses for the computation of the planet’s mass and orbital parameters.
We employed the $`\chi ^2`$ test to analyze 160 000 stars uniformly distributed on the sky, perturbed by dark companions inducing astrometric signatures $`\alpha `$ ranging from 10 to 100 $`\mu \mathrm{as}`$, with periods P between 0.5 and 20 years. The remaining orbital elements were distributed randomly in the ranges: $`0^{}i90^{}`$, $`0e1`$, $`0\mathrm{\Omega }2\pi `$, $`0\omega 2\pi `$, $`0\tau P`$. The single observation error was set to $`\sigma _\psi `$ = 10 $`\mu `$as, thus the signal-to-noise ratio varied in the range $`1S/N10`$. Finally, we repeated the same simulations without planets, to verify the correct behaviour of the test and the choice of the confidence level. As expected the number of false detections was $``$ 5 per cent.
### 3.2 Orbital parameters
Once a planet is “detected”, the goal is to derive reliable estimates for of its orbital elements and mass. For a complete reconstruction of the orbital geometry, we implemented an analytic model in which the observation residuals $`\psi \psi _r`$ are evaluated with a recomputed abscissa of the form:
$$\psi _r=\psi _r(\lambda ,\beta ,\mu _\lambda ,\mu _\beta ,\pi ,T_1,T_2,T_3,T_4,e,P,\tau )$$
(3)
The 1-d measurement along the scanned GC is thus expressed as function of both the standard astrometric parameters ($`\lambda `$, $`\beta `$, $`\mu _\lambda `$, $`\mu _\beta `$, $`\pi `$) and the parameters describing the stellar relative orbit with respect to the barycenter of the planetary system: period $`P`$, periastron epoch $`\tau `$, eccentricity $`e`$ and the 4 Thiele-Innes elements $`T_i`$, functions of semi-major axis $`a`$, inclination $`i`$, argument of periastron $`\omega `$, position angle of the ascending node $`\mathrm{\Omega }`$.
The solution of the non linear system of observation equations is obtained employing an iterative linear Least Squares procedure by means of which the entire set of orbital elements can be simultaneously estimated within a well defined accuracy level. The details of this procedure are given in the following sections.
## 4 RESULTS
In this section we account for what is our present understanding of GAIA’s sensitivity to astrometric perturbations induced by Jupiter-mass planets orbiting around nearby stars. Given the present evidence of a number of extra-solar planets, we also present results on how GAIA would perform on a selection of three such systems.
### 4.1 Detection probabilities
Figure 1 gives the percentage of failure of the $`\chi ^2`$ test on the single-star hypothesis in the case of $`\sigma _\psi =10`$ $`\mu \mathrm{as}`$. We note that at relatively low $`S/N`$ ratios the detection probability is dominated by sampling of the orbital period. At higher $`S/N`$ values orbital sampling is less critical and long period planets (up to about twice the mission duration) are detectable: as a matter of fact, when $`S/N`$ 10 detection probability reaches about 100 per cent.
The dip at $`P=1`$ year was somewhat expected, as both orbital motion and parallactic factor have the same period. However, the decrease is small, indicating that the coupling is less critical than might have been anticipated. This is probably due to the fact that the parallactic motion has fixed phase and aspect ratio, and a relatively small mismatch in any of the orbital parameters—phase, inclination, eccentricity—is sufficient to separate the two signals.
Figure 1 also shows that the $`\chi ^2`$ test quickly loses its sensitivity for $`S/N`$ approaching unity; thus planets for which the error in individual observations is comparable to the apparent semi-major axis are essentially undetectable with this technique.
As mentioned before, the results of Figure 1 apply to measurement errors other than the assumed $`\sigma _\psi =10`$ $`\mu \mathrm{as}`$, as long as the $`S/N`$ value remains the same. For example, for a single observation error $`\sigma _\psi =1`$ $`\mu \mathrm{as}`$, the detection probability is exactly the same as shown in Figure 1, but for amplitudes ten time smaller.
We can get important physical informations from the statistical results of Figure 1, simply changing the interpretative perspective from which we are leading the discussion. The solid lines in Figure 2 are the empirical relations, derived from Figure 1, that express the amplitude of the perturbation as a function of orbital period needed for detection probabilities of 25, 50, and 95 per cent, respectively. If the orbital period is shorter than the assumed mission life-time of 5 years, the detection probability is about 50 per cent for $`S/N1`$. For longer periods, the sampling of the orbit is worse, probability drops significantly and a much higher signal is required for the planet’s signature to be detected. This is in qualitative agreement with the results of Babcock \[Babcock 1994\], who studied the detection and convergence probability of a complete orbital model for simulated planetary systems as observed by the mission POINTS. Babcock did find a slightly larger sensitivity on planet’s period, manifested in an earlier turn-up and steeper slope at long periods of the 50 per cent probability curve; this most likely depends on the fact that the determination of reliable orbital elements is more challenging than detection only.
Let us now discuss the other elements of Figure 2. The overplotted dashed and dotted lines represent the signal expected for a Jupiter-mass planet and solar-mass star at various distances and orbital periods, obtained by substituting Kepler’s third law in the expression for $`\alpha `$:
$$\alpha \frac{M_p}{M_s^{2/3}}\frac{P^{2/3}}{d}$$
(4)
The vertical solid line at $`P`$ = 11.8 years indicates the locus of the actual Jupiter-Sun system at different distances. Jupiter-Sun systems appear detectable with probability $`50`$ per cent up to a distance of 100 pc, while Jupiter-mass planets with shorter periods can be detected to larger distances: to 150 pc for periods between 2.5 and 8 years. The relation can be rescaled to lower planet masses by reducing the distance of the system in proportion to the planet mass, since it is the ratio $`M_p/d`$ that enters in the astrometric signature $`\alpha `$.
Periodic photospheric activity across the disk of the parent star (e.g. star-spots cycles) can induce displacements in the position of the photocenter, thus adding astrometric signal which could be difficult to disentangle from the planet signature. However, the magnitude of such intrinsic astrometric noise appears to be at most of a few $`\mu `$as \[Woolf & Angel 1999\] for a solar-type star at a distance of 10 pc, i.e., significantly smaller than the expected single-measurement errors.
Finally, although the single-observation error significantly deteriorates with magnitude (Table 1), jovian planets around late type stars can reliably be detected at relatively large distances.
### 4.2 Orbital parameters estimation
The next step beyond the simple detection of candidate planets is the task of estimating the orbital parameters of each system and the mass of the unseen companion.
Our method, described in more detail below, consists of applying an iterative non-linear least-squares fitting procedure to each of the simulated orbits; the ’known’ orbital parameters are utilized as initial guesses to start the fit. Convergence of the non-linear fitting method and quality of the orbital solutions can both be significantly affected by the choice of the starting guesses. Therefore, the use of the true values of the orbital parameters to initialize the fit leaves open some important questions related to how and to what extent effective starting values, i.e., leading to successful orbital solutions, will be identified from the data as function of actual performances of the satellite, uncertainties in the error model, and properties intrinsic to planetary systems. For this, realistic global search strategies must be implemented and double-blind test campaigns conducted. Work on these issues is in progress and will be reported in the future.
Instead, this work focuses on the important goal of gauging GAIA’s ultimate ability in measuring planets (properly, single giant planets orbiting single solar-type stars). An efficient way of achieving this is by assuming perfect knowledge of GAIA’s characteristics (mainly, error model and satellite attitude) and, indeed, by using the known values of each orbital parameter as suitable initial guesses for the least-squares solutions <sup>2</sup><sup>2</sup>2Note that, as explained later in this section, the fitting program does not know that the initial values provided for the parameters are also their true values. Then, the post-fit differences to the true values of the parameters should be a reliable measure of the accuracy in orbit reconstruction that can ultimately be achieved by GAIA with the given measurement errors.
Similarly to what we did in section 4.1, we look at the accuracy of the results primarily as function of distance and orbital period. We consider systems with Sun-Jupiter masses and three values of the period: 0.5 years, to test GAIA’s ability to cope with poorly-sampled motion; 5 years, a near-ideal case where the orbital period is as long as the mission; and 11.8 years, the true period of the Jupiter-Sun system, which stretches GAIA’s ability to solve long-period orbits. For each case, we generate 200 systems randomly placed on the sky, with randomly distributed values of the other orbital parameters. These 200 systems are then placed all at the same distance, and the simulation is repeated for distance values ranging from from 10 to 100 pc.
A fit to the simulated observations is made directly at the GC level, taking into consideration a theoretical model in which GAIA’s unidimensional measurements are expressed as functions of the 5 astrometric parameters plus the 7 orbital elements of the keplerian orbit: for the latter we adopted a slightly linearized analytic form where semi-major axis, inclination, periastron longitude and position angle of the ascending node are combined in the four Thiele-Innes elements, while only orbital period, eccentricity and periastron epoch are left unchanged.
The first solution of the linearized system is thus found with respect to the initial guess. Then, the linearization at step $`k`$ is updated with respect to the solution obtained at step $`k1`$ and the process repeated until the differential corrections $`\delta a_{i,k}`$ to each of the 12 parameters satisfy the relation: $`\left|\delta a_{i,k}/a_{i,k1}\right|<10^6`$, where $`i=1,2,\mathrm{},12`$ and $`a_{i,k1}`$ are the parameters adjusted at the previous step. We then record what fraction of the final values for each parameter falls within a certain fractional error of the true value: this we call the “convergence probability”, which is a function of the distance of the system, the period, and the desired fractional error.
We evaluate the convergence probability for the parameters which are most likely to affect the efficiency of the reconstruction of the orbit and of the mass determination, namely semimajor axis $`a`$, period $`P`$, inclination $`i`$, and eccentricity $`e`$. For each parameter, we consider fractional errors of 10, 30, and 50 per cent.
The results are shown in Figures 3 to 5. Points of different shapes correspond to different periods: triangles for 0.5 years, diamonds for 5 years, and crosses for 11.8 years. Figure 3 shows the probability of convergence within 10 per cent of the true values; Figure 4 to within 30 per cent; and Figure 5 to within 50 per cent. In each Figure, different panels correspond to different quantities.
As expected, the 5 year-period case is the best of the three. The 11.8 year period, although it corresponds to a larger signal (as the semimajor axis, and thus $`\alpha `$, increases with period), suffers from the incomplete sampling of the orbit during the mission life-time, while the short-period case suffers from both the smaller signal amplitude and a (generally) non-optimal timing of the observations.
The periodicity of the signal is the characteristic which can be evaluated with the best accuracy: $`P`$ is the only element which is always estimated with high accuracy (better than 10 per cent) throughout the ranges covered by our simulations. We also note that the semimajor axis $`a`$ and the inclination $`i`$ behave similarly to each other. This may be in part a consequence of the use of the Thiele-Innes representation, since both parameters are obtained from combination of the Thiele-Innes elements.
Somewhat surprisingly, orbital eccentricity—which in the previous simulation is assumed to be distributed uniformly from 0 to 1—has a very significant effect on the quality of the estimated orbital parameters, including the mass of the planet. We illustrate this in Figures 6 through 8, where the different symbols now refer to different orbital eccentricities, and the convergence probability is given for a fixed fractional error of 20 per cent. Different figures correspond to different orbital periods. Generally, high eccentricities deteriorate convergence percentages for all orbital parameters. The reason is that the regularly-spaced observations of a survey satellite, such as GAIA, may be ill-suited to sample an orbital motion with large velocity variations, as happens for high eccentricities. The effect is especially prominent for the long-period case, $`P=11.8`$ years, for such orbits will often never be observed during the periastron.
Our main results can thus be summarized as follows:
* For a uniform eccentricity distribution in the range $`0e1`$, then more than a half of all existing Jupiter-mass planets orbiting around solar-type stars with periods comparable to the mission lifetime can be detected, and their masses and orbital parameters can be estimated to an accuracy of 10 per cent up to distances of about 100 pc. From the completeness limit of the HIPPARCOS catalog ($`V7.58`$ mag) and considering spectral types no later than G5, we estimate 20 000, 65 000, and 150 000 stars to 100, 150, and 200 pc respectively. In this estimate the contribution from early type stars and giants is negligible. Note that these numbers increase significantly when considering spectral types earlier than K5: calculations based on current Galaxy models predict $`\stackrel{>}{}`$ 300 000 stars within 200 pc from the Sun \[Lattanzi et al. 1999\].
* For a fractional error no worse than 20 per cent, the distance limit is reduced to 50 pc for short-period and long-period orbits (0.5 and 11.8 years);
* Low-eccentricity systems are easier to detect and solve: this could indicate a potential limit in GAIA’s capability of reconstructing the precise behaviour of the mass function of low mass companions to solar-type stars in the neighborhood of our Sun, as to our knowledge first observational data seem to indicate than brown dwarfs are more likely to revolve around the parent stars on significantly eccentric orbits \[Black 1997\]. On the other hand, low orbital eccentricities may be prevalent for longer periods, thus enhancing GAIA’s capabilities for such systems.
### 4.3 47 Uma, 70 Vir, and 51 Peg
We now consider more specifically how the known planetary systems fall within GAIA’s capabilities. Figure 9 shows period and astrometric signature for the planets detected to date by radial velocity measurements as compared to GAIA’s 50 per cent iso-probability curve. This is the minimum signature, corresponding to an orbital inclination of 90; the radial velocity method cannot generally remove the degeneracy between planetary mass and inclination. Although all systems are very close to the Sun ($`d\stackrel{<}{}40`$ pc), several planets have very small signatures due to their short orbital period (and thus small semi-major axis, by Kepler’s third law). Some of the shortest-period objects, i.e, the 51 Peg-class planets, may prove difficult to detect and measure with GAIA. On the other hand, planets with signature $`100\mu \mathrm{as}`$ will be easily detected, and their orbital elements can be found with high accuracy by GAIA.
The considerations and results obtained in the previous sections from a statistical viewpoint are obviously preliminary to the much more difficult task of the development of self-consistent detection and orbital parameters estimation algorithms devoted to the investigation of individual objects, in which a more complete and realistic error model can be taken into account: to this end detailed system studies are in progress and will be presented elsewhere. It is nevertheless of interest a first glance to how well GAIA will behave once at work: thus, to test further GAIA’s capabilities of detecting periodic oscillations and signatures due to planetary companions around stars in the neighborhood of our solar system, we have concentrated our attention on the results obtained with radial velocity techniques, for what concerns three of the star-planet systems known to date: 47 UMa, 70 Vir and 51 Peg. All these stars are nearby and very similar to our Sun.
According to the spectroscopic measurements \[Butler & Marcy 1996, Marcy & Butler 1996, Mayor & Queloz 1995\], around these three stars orbit planetary bodies of minimum mass, respectively, $``$2.46 $`M_\mathrm{J}`$, $``$6.50 $`M_\mathrm{J}`$, and $``$0.5 $`M_\mathrm{J}`$, with orbital periods of about 3 years, 4 months, 4 days. This translates in the following astrometric signatures:
$`\alpha `$ (47 UMa) $``$ 361.88 $`\mu `$as
$`\alpha `$ (70 Vir) $``$ 167.93 $`\mu `$as
$`\alpha `$ (51 Peg) $``$ 1.66 $`\mu `$as
Because of its extremely short period, the astrometric perturbation in the case of the 51 Peg system is so small that reliable estimates of its astrometric orbit with GAIA will be highly difficult. On the other hand, the detection of the signatures induced on 47 UMa and 70 Vir will be extremely easy for a GAIA-like satellite.
We have generated 100 simulated planetary systems on the celestial sphere, respectively identical to 47 UMa and 70 Vir, assuming perfectly circular orbits and stellar masses $`M`$ = $`M_{}`$. The unknown orbital inclination was chosen to be $`i=45^{}`$.
First of all, we applied the $`\chi ^2`$ test to the simulated observational data, with a conservative single-observation error $`\sigma _\psi =10`$ $`\mu \mathrm{as}`$. The detection probabilities derived are, as expected, always about 100 per cent, given the signal-to-noise ratio to be always in the regime $`S/N1`$.
In order to obtain a visual indication of the quality of the reconstruction, we then compared the apparent motion of the star with the motion obtained with a single star fit and that obtained by a fit with the fitted parameters for the planet. Figure 10 shows the differences in 47 UMa’s path on the sphere, when calculated starting from: a) the true values of $`\mu _\lambda `$, $`\mu _\beta `$, $`\pi `$, adopted in the simulation (asterisks, solid line); b) the values of $`\mu _\lambda `$, $`\mu _\beta `$, $`\pi `$, obtained after fitting the observations with a single star model (crosses, dotted line); and c) the values of $`\mu _\lambda `$, $`\mu _\beta `$, $`\pi `$ and of the orbital parameters derived after fitting the observations with a 10 parameter model reproducing the standard motion plus the Keplerian circular motion induced on the observed star by the presence of the planet (diamonds, dashed line).
The annual proper motion of 47 UMa is very large, compared to the magnitude of the astrometric signature produced by the gravitational influence of its planetary companion: in ecliptic coordinates, $`\mu _\lambda `$ = -0.283 arcsec/year, $`\mu _\beta `$ = 0.152 arcsec/year, while $`\alpha 0.5`$ mas for an orbital inclination of $`45^{}`$. The top left panel shows the full range of motion of 47 UMa over five years; the looping motion is the parallactic eclipse, and the effect of the planet is essentially undetectable on this scale. In order to see the difference in the residuals, we zoom first on region A marked by a small square in the top left panel, then we identify in the upper right panel two 1-mas<sup>2</sup> regions (B and C), representing each a 1000x enlargement of a small fraction of the motion (bottom panels). The true and reconstructed motions are very close, consistent with the $`10\mu `$as single-measurement error of the observations used to derive the orbital parameters, while on this scale the orbit obtained with the single-star assumptions (dotted) is clearly very different from the actual motion.
Third, we derive the uncertainties in the fitted values of the important physical parameters of the two systems, and especially the companion masses. We derive the mass from the fitted orbital parameters via the mass function:
$$\frac{M_p^3}{(M_s+M_p)^2}=\frac{a_s^3}{\pi ^3}\frac{1}{P^2}$$
(5)
where again $`M_p`$ and $`M_s`$ are the planetary and stellar masses in solar masses, $`\pi `$ the parallax, $`P`$ the orbital period in years, and $`a_s`$ the semi-major axis of the parent star, expressed in arcsec.
Assuming $`M_pM_s`$, we obtain:
$$M_p\left(\frac{a_s^3}{\pi ^3}\frac{M_s^2}{P^2}\right)^{1/3}$$
(6)
Fitting the simulated observations with a 10 parameter model we can derive estimates of $`\pi `$ and $`P`$ directly, while $`a_s`$ can be obtained as a combination of the approximated four Thiele-Innes elements. Provided that one gets informations about the mass of the primary, for example thanks to spectroscopy, it is possible to calculate $`M_p`$.
In our case we have, for simplicity, supposed to know $`M_{47UMa}`$ and $`M_{70Vir}`$ to be exactly equal to $`M_{}`$.
Tables 2 and 3 show the approximated values of $`M_p`$ in the two cases, starting from the estimated values of $`P`$, $`\pi `$ and $`a_s`$ obtained during the fit, as function of orbital inclination. The Tables stop at $`i=20^{}`$ for 47 UMa and $`i=10^{}`$ for 70 Vir, taking into consideration the upper limits on their masses found by Perryman et al \[Perryman et al. 1996\] using HIPPARCOS data. The HIPPARCOS measurements, as a matter of fact, agreed with the stars to be single, as no evident signature greater than the nominal error ($``$1 mas) was revealed. This implies the following upper limits: $`M_p`$$``$ 7 $`M_\mathrm{J}`$, for 47 UMa, $`M_p`$$``$ 38 $`M_\mathrm{J}`$, for 70 Vir, which again means, considering the results obtained spectroscopically, choosing a minimum orbital inclination as in the tables. In both cases, it is clearly evident the high accuracy with which the relevant orbital parameters and the planet’s mass are recovered: thanks to the fact that $`S/N1`$, the fit to GAIA’s simulated observations of 47 UMa and 70 Vir is very satisfactory; the RMS errors between fitted and nominal values of the various parameters are always well within 10 per cent.
Such results confirm that GAIA could reveal itself a very powerful instrument to investigate many of the known candidate planetary systems, with the exception of very short-period systems, such as 51 Peg: provided the signal-to-noise ratio is sufficiently high and the period not too short, our simulations show that such systems will be easily detected, and their orbital parameters accurately determined by GAIA, at least if such systems are effectively simple, as the first spectroscopic measurements would suggest. In case the signatures produced cannot be interpreted as due to the presence of only one planet, this would introduce many interpretative complications in the signals observed: for planetary systems resembling our own, detection should be little affected and residuals analysis might reveal hints of the influence of the outer planets ($`\alpha `$ for a Saturn-like companion is $``$ 60 $`\mu \mathrm{as}`$ at 10 pc but its period is $``$ 6 times the mission lifetime). However, reliable orbital fitting would probably be restricted to the main component. On the other hand, multiple orbital fitting to good accuracy should be possible for systems composed of giant planets of comparable mass, orbiting with periods within the interval $``$0.5-5 yr, and yielding astrometric S/N $`\stackrel{>}{}`$ 10 , like it would be the case for the two outermost planets of the recently discovered system $`\upsilon `$ And \[Butler et al. 1999\].
The capability of detecting and measuring multiple planets with GAIA will become matter of future simulations and system studies.
## 5 Summary and conclusions
In this work we have given the first quantitative evaluation of the detectability horizon of single extra-solar giant planets around single normal stars in the neighborhood of our solar system for the global astrometry mission GAIA. Complete simulations, comprehensive of observations of star-planet systems and successive statistical analysis of the simulated data, have yielded the following results:
* it will be possible to detect more than 50 per cent of all Jupiter-like planets (orbital period $`P=11.8`$ years) orbiting solar-type stars within 100 pc; Jupiter-size planets, with shorter orbital periods, will be detectable up to 200 pc, with similar probabilities;
* for true Sun-Jupiter systems it will be possible to determine the full set of orbital parameters and to derive accurate estimates of the masses up to distances of order of 50 pc, value which doubles if we consider the range of periods in the vicinity of the mission lifetime;
* simulated observations of a selection of the actually known extra-solar planets, discovered by means of spectroscopic measurements, provide a meaningful estimate of the uncertainty with which masses and orbital elements can be determined for the known star-planet systems and for a substantial fraction of those that will be found within the context of such a global astrometry mission. Although preliminary, our results indicate that these systems will be easy to discover and their orbital parameters will be accurately determined with GAIA, except very short-period systems such as 51 Peg.
Hence, our results indicate that: $`1)`$ GAIA would monitor all of the hundreds of thousand F-G-K stars (i.e., whose masses are within a factor $``$ 1.5 that of the Sun) up to a distance of $``$ 200 pc from the Sun, in search for astrometric signatures due to the presence of giant planets ($`MM_J`$) with orbital periods up to Jupiter’s; $`2)`$ a significant fraction of the detected planets would have the main orbital parameters (semi-major axis, period, eccentricity, inclination) measured to better than 30 per cent accuracy.
Therefore, the GAIA survey would uniquely complement the expectations from other ongoing and planned spectroscopic and astrometric planet searches, both from ground and in space, thus helping with the creation of the fundamental testing ground on which to measure the validity of actual theoretical models of formation and evolution. GAIA’s discovery potential might have significant impact on our knowledge of the distribution laws of the most relevant orbital parameters, and it would contribute to determine the frequency of planetary systems themselves in the solar neighborhood and, by extrapolation, in the whole Galaxy. A vast all-sky astrometric survey would help understand peculiar characteristics of these systems, e.g., whether giant planets lying in the outer regions are common: such planetary scenarios may be worth further investigation, as, according to present theoretical models, this could indicate presence of low mass planets in the inner regions, possibly in the parent stars’ habitable zones. The monitoring of hundreds of thousands stars directly implies the chance to investigate objects belonging to a wide range of spectral types, thus providing the important observational material \[Boss 1998\] to decide whether giant planets are more likely to form by means of gravitational instability in disks (once they are found more often around young stars \[Kuiper 1951, Cameron 1978, Bodenheimer et al. 1980\]), or by means of accretion of planetesimals (once they are found more often around old stars \[Pollack 1984, Lissauer 1987, Pollack et al. 1996\]). The high precision global astrometric measurements will estimate the inclination $`i`$ of the orbital planes for the majority of the presently known planetary systems and for a large fraction of those that will be eventually discovered: it will then be possible to provide unambiguous mass estimations of such dark companions, reducing significantly the uncertainty on the mass range in the transition region from brown dwarfs to giant planets.
## Acknowledgments
We wish to offer our special thanks to M.A.C. Perryman for lending initial impetus and continuing support to this investigation. Over the course of this work, we have benefited from discussions with numerous colleagues, and especially Gerry Gilmore, David Latham, Lennart Lindegren, Robert Reasenberg and Stuart Shaklan. Also, we wish to thank the referee for her/his careful comments which helped us improve the original manuscript. All four authors gratefully acknowledge partial financial support from the Italian Space Agency, under contract ASI/ARS-96-77.
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# Edinburgh 2000/12LTH 455hep-lat/xxx Instability in the Molecular Dynamics Step of Hybrid Monte Carlo in Dynamical Fermion Lattice QCD Simulations
## I Introduction
Hybrid Monte Carlo (HMC) remains the most widely used algorithm for lattice QCD computations with dynamical fermions. In such computations, trial configurations are produced by integrating the Hamiltonian equations of motion from an initial configuration for some fictitious molecular dynamics (MD) time $`\tau `$. Configurations are then accepted or rejected by subjecting the energy change $`\delta H`$ along a trajectory to a Metropolis accept/reject step.
It has been observed that the equations of motion in the MD evolution of such an algorithm are chaotic in the case of QCD. This implies that rounding errors induced by the use of finite precision in a digital computer may grow exponentially. Such growth can be characterised in terms of the leading Liapunov exponent of the system. Furthermore, it has been shown that the most commonly used MD integration scheme – the leapfrog method – has the potential to become unstable. Instability is a problem for lattice QCD simulations since it results in large energy changes along MD trajectories and hence negligible acceptance rates in the HMC algorithm.
The instability in the leapfrog method has been illustrated in for the case of free field theory where a mechanism has been proposed which could explain the onset of such an instability in lattice QCD. Numerical studies of the latter were carried out on small lattices at a variety of couplings and quark masses. The onset of instability was found to be at smaller step sizes for lighter quark masses.
Edwards, Horváth and Kennedy also investigated an optimisation strategy in which reduced work (and hence accuracy) in the MD calculation was balanced against the resulting reduced acceptance in the Metropolis step. Each MD step requires the iterative solution of a system of linear equations. Since dynamical fermion HMC codes spend a substantial fraction of their execution time performing such solves it it clearly important to investigate whether substantial efficiency gains can be made without introducing undesirable effects such as loss of reversibility in the MD. The investigation was quite preliminary and the errors quoted were quite large. This issue was also investigated on small lattices in . The present paper investigates many of the issues raised in and extends the numerical studies to production-scale lattices.
The paper is organised as follows. In section II we summarise the Hybrid Monte Carlo formalism and give details of the algorithms used. Section III contains a discussion of the effects of numerical roundoff errors on reversibility. In section IV we present results and discussion of our analysis of instability in the MD step. In section V we present the results of an optimisation analysis involving reduced accuracy in the MD step.
Finally, in section VI we summarise our results and conclusions.
## II Hybrid Monte Carlo and lattice QCD
### A HMC algorithm
Consider a system with canonical coordinates $`q`$ and action $`S(q)`$. One wishes to generate configurations $`q`$ with an equilibrium probability distribution in which the statistical weight of configuration $`q`$ is proportional to $`e^{S(q)}`$.
In Hybrid Monte Carlo, we introduce fictitious momenta $`p`$ conjugate to $`q`$ and define a Hamiltonian function $`H(q,p)=\frac{p^2}{2}+S(q)`$.
One may then generate configurations $`(q,p)`$ distributed according to
$$P(q,p)dqdp=\frac{1}{Z}e^{H(q,p)}dqdp\text{where}Z=𝑑q𝑑pe^{H(q,p)}.$$
(1)
After the integration over the momenta, we obtain the desired distribution for the coordinates.
Given an initial configuration $`(q,p)`$, a sequence of configurations is generated by repeated iteration of the following steps:
1. Momentum refreshment: Draw new fictitious momenta $`p`$ from a Gaussian distribution with zero mean and unit variance.
2. Molecular dynamics: Integrate the Hamiltonian equations of motion for some fictitious time trajectory of length $`\tau `$, from the initial configuration $`(q(0),p(0))=(q,p)`$ to obtain the trial configuration $`(q(\tau ),p(\tau ))=(q^{},p^{})`$.
3. Accept/reject step: The trial configuration $`(q^{},p^{})`$ is accepted with probability
$$P_{\mathrm{acc}}(q^{},p^{}q,p)=\mathrm{min}(1,e^{\delta H})$$
(2)
where
$$\delta H=H(q^{},p^{})H(q,p).$$
(3)
If the trial configuration is rejected the new configuration is $`(q,p)`$.
### B Leap–frog integration
For the HMC algorithm to satisfy detailed balance, the MD is required to be reversible and measure preserving. This can be achieved through the use of symmetric symplectic integration schemes, such as the leapfrog algorithm. In this algorithm, one constructs an approximation $`𝒰_3(\delta \tau )`$ to the time evolution operator $`𝒰(\delta \tau )`$ for advancing a phase space vector $`(q,p)`$ through a step of length $`\delta \tau `$ in molecular dynamics time. The approximate operator $`𝒰_3(\delta \tau )`$ is itself composed of a symmetric combination of the symplectic partial coordinate and momentum update operators $`𝒰_\mathrm{q}(\delta \tau )`$ and $`𝒰_\mathrm{p}(\delta \tau )`$ respectively, for example as
$$𝒰_3(\delta \tau )=𝒰_\mathrm{p}\left(\frac{\delta \tau }{2}\right)𝒰_\mathrm{q}(\delta \tau )𝒰_\mathrm{p}\left(\frac{\delta \tau }{2}\right).$$
(4)
The partial update operators are themselves defined as
$`𝒰_\mathrm{q}(\delta \tau )(q,p)`$ $`=`$ $`(q+p\delta \tau ,p)`$ (5)
$`𝒰_\mathrm{p}(\delta \tau )(q,p)`$ $`=`$ $`(q,p+F\delta \tau )`$ (6)
where $`F=\frac{S}{q}`$ is the MD force. Due to its symmetric construction, $`𝒰_3(\delta \tau )`$ is reversible and, due to the symplectic nature of its component updates, it is area preserving. The process of iteratively acting on an initial phase space vector with $`𝒰_3(\delta \tau )`$ is called leapfrog integration. The method is accurate to $`O(\delta \tau ^3)`$ per time step.
### C Higher order integration schemes
The construction of higher order integration schemes (see for example ) is recursive, proceeding from the leapfrog scheme. Given an approximate time evolution operator $`U_{n+1}(\delta \tau )`$ accurate to $`O(\delta \tau ^{n+1})`$ for some even $`n`$, one can construct the operator
$$𝒰_{n+3}(\delta \tau )=𝒰_{n+1}(\delta \tau _1)^i𝒰_{n+1}(\delta \tau _2)𝒰_{n+1}(\delta \tau _1)^i$$
(7)
with
$`\delta \tau _1={\displaystyle \frac{\delta \tau }{2is}}`$ (8)
$`\delta \tau _2={\displaystyle \frac{\delta \tau }{1\frac{2i}{s}}}`$ (9)
where $`i`$ is an arbitrary positive integer and $`s=(2i)^{\frac{1}{n+2}}`$. The step sizes $`\delta \tau _1`$ and $`\delta \tau _2`$ are chosen to cancel truncation orders of $`O(\delta \tau ^{n+1})`$ and symmetry with respect to time ensures that there are no truncation errors of $`O(\delta \tau ^{n+2})`$. Hence such a scheme is accurate to $`O(\delta \tau ^{n+3})`$.
Sexton and Weingarten have considered the general case where the action $`S`$ can be split into two parts as $`S(q)=S_1(q)+S_2(q)`$ and constructed an $`O(\delta \tau ^3)`$ algorithm in which the coefficient of leading order truncation error term may be decreased. The method is advantageous if evaluating the force corresponding to $`S_1`$ is computationally much cheaper than the force associated with $`S_2`$ (or vice versa). For example, one may take $`S_1`$ to be the gauge action and $`S_2`$ to be some computationally expensive fermion action. The coefficient of the leading error term could then be decreased by performing more gauge update steps than momentum updates.
### D Formulation of MD for lattice QCD
The canonical coordinate variables for lattice QCD are the $`SU(3)`$ link matrices $`U_\mu (x)`$ associated with the link emanating from site $`x`$ of the lattice and ending on neighbouring site $`x+\widehat{\mu }`$, where $`\widehat{\mu }`$ is a unit vector in one of the Euclidean space–time directions. The conjugate momentum fields $`\pi _\mu (x)`$ are members of the Lie algebra $`su(3)`$.
In general, one can write the fictitious Hamiltonian for a lattice QCD system with two degenerate flavours of Sheikholeslami-Wohlert (clover) improved fermions as
$$\stackrel{~}{H}=\frac{1}{2}\underset{x,\mu }{}\pi _\mu ^2+S_\mathrm{g}(\beta ;U)+\varphi ^{}\stackrel{~}{Q}^1(\kappa ,c;U)\varphi ,$$
(10)
where
$$\stackrel{~}{Q}(\kappa ,c;U)=M^{}(\kappa ,c;U)M(\kappa ,c;U).$$
(11)
Here $`M(\kappa ,c;U)`$ is the clover improved fermion matrix with improvement coefficient $`c`$, $`\varphi `$ are pseudofermions and $`S_\mathrm{g}(\beta ;U)`$ is the standard Wilson gauge action
$$S_\mathrm{g}(\beta ;U)=\frac{\beta }{6}\underset{\mathrm{}}{}\mathrm{Re}\mathrm{Tr}U_{\mathrm{}}.$$
(12)
In (12) the sum is over all elementary plaquettes $`U_{\mathrm{}}`$ on the lattice and $`\beta =\frac{6}{g^2}`$ where $`g`$ is the bare gauge coupling constant.
In our computations we have employed the technique of even–odd preconditioning which changes the form of $`\stackrel{~}{Q}`$ and $`\stackrel{~}{H}`$ somewhat. Each lattice site is labelled with a parity $`p`$ which is either even or odd so that any one lattice site has an opposite parity from all of its neighbours. This allows the fermion matrix to be block diagonalised and the Hamiltonian to be re–written as:
$$H=\frac{1}{2}\underset{x,\mu }{}\pi ^2+S_\mathrm{g}(\beta ;U)2\mathrm{Tr}\mathrm{ln}A_e+\varphi _o^{}Q^1(\kappa ,c;U)\varphi _o.$$
(13)
Here, $`A`$ is the so called clover term summed over sites of one parity (even in the equation above) and $`Q`$ is the preconditioned fermion matrix coupling lattice sites of the opposite parity (odd in the equation above) only. Thus $`Q`$ has half the rank of $`\stackrel{~}{Q}`$. This leads to some memory saving at the additional expense of having to evaluate $`\mathrm{Tr}\mathrm{ln}A`$ directly on sites of one parity. The precise formulation of the preconditioned matrices can be found in .
We do not expect that splitting the Hamiltonian in this way will change conclusions regarding reversibility and related issues in any significant way. Although there is an extra force term to be computed to integrate the equations of motion, the logarithm of the clover term is computed directly and is independent of the parameters used for the solution of the system of linear equations. Likewise, for the inversion of the clover term, we use a direct method that is not controlled by algorithmic parameters such as a target relative residue. Hence we regard the effects of preconditioning as a minor technicality and shall disregard them for the rest of this paper.
The leapfrog partial update steps for the gauge fields and the momenta are
$`𝒰_q(\delta \tau )(U_\mu (x),\pi _\mu (x))`$ $`=`$ $`(\mathrm{exp}\{i\delta \tau \pi _\mu (x)\}U_\mu (x),\pi _\mu (x))`$ (14)
$`𝒰_p(\delta \tau )(U_\mu (x),\pi _\mu (x))`$ $`=`$ $`(U_\mu (x),\pi _\mu (x)+\delta \tau F_\mu (x))`$ (15)
where
$$F_\mu (x)=F_\mu ^\mathrm{g}(x)+F_\mu ^\mathrm{f}(x)$$
(16)
and $`F^\mathrm{g}`$, $`F^\mathrm{f}`$ are the respective gauge and fermionic force contributions,
$`F_\mu ^\mathrm{g}(x)`$ $`=`$ $`{\displaystyle \frac{S_\mathrm{g}(U)}{U_\mu (x)}}`$ (17)
$`F_\mu ^\mathrm{f}(x)`$ $`=`$ $`\left[Q^1\varphi \right]^{}{\displaystyle \frac{Q}{U_\mu (x)}}\left[Q^1\varphi \right].`$ (18)
### E Solution of the linear system
Computation of the fermion force requires the quantity
$$X=Q^1\varphi $$
(19)
which is obtained via the solution of the linear system
$$QX=\varphi .$$
(20)
This is normally carried out with a Krylov subspace solver such as the Conjugate Gradients (CG) or the Stabilized BiConjugate Gradients (BiCGStab) algorithm. With the BiCGStab solver, the solution consists of two solves:
$`M^{}(\kappa ,c)Y`$ $`=`$ $`\varphi `$ (21)
$`M(\kappa ,c)X`$ $`=`$ $`Y`$ (22)
whereas with CG, one can solve (20) directly. When using CG with a Hermitean positive definite matrix such as $`Q`$, the solution is guaranteed to converge monotonically. With BiCGStab, one has no such guarantee. Since the condition number of $`Q`$ is the square of the condition numbers of either $`M`$ or $`M^{}`$, we expect the two stage solution using BiCGStab to be faster on the whole than using one CG solve. As the convergence of BiCGStab can be erratic, it is prudent to restart the solution process for $`X`$ with CG using, as an initial guess, the solution for $`X`$ from the previous BiCGStab solve.
The solver residual $`r_i`$ at the $`i`$-th iteration of a CG solve is defined as
$$r_i^{\mathrm{Real}}=\varphi QX_i$$
(23)
where $`X_i`$ is the approximate solution at iteration $`i`$. The relative residual at the $`i`$-th iteration is then defined as
$$\rho _i^{\mathrm{Real}}=\frac{r_i^{\mathrm{Real}}}{\varphi }.$$
(24)
In solver algorithms, $`r_i`$ is not usually computed using (23). Instead, $`r_i`$ is generally defined through some three term or coupled two term recurrence relation. We will refer to this latter definition of the residual as $`r^{\mathrm{Acc}}`$, the accumulated residual. The corresponding definition of the relative residual is
$$\rho _i^{\mathrm{Acc}}=\frac{r_i^{\mathrm{Acc}}}{\varphi }.$$
(25)
These two definitions are equivalent in exact arithmetic. However, computation of the accumulated residual needs only vector additions and scalar multiplications whereas computation of the real residual needs a matrix multiplication and so can differ in finite arithmetic. In our computations we use the accumulated residual. We will denote by $`r`$ our target relative residual. Hence the iterative process terminates when $`\rho _i^{\mathrm{Acc}}<r`$. In the remainder of this paper we refer to $`r`$ as the solver target residual, or just simply the solver residual.
## III Reversibility
Reversibility and area preservation of the Molecular Dynamics step are required for a correct HMC algorithm. The leapfrog algorithm described in section 2, is reversible and area preserving in exact arithmetic. Computations are of necessity carried out in finite precision and exact reversibility is lost. It is therefore important to verify that implementation of the fundamental steps of the algorithm are as close to reversible as it is possible to make them.
Ideally, one would like to establish the least level of precision required such that the accumulation of rounding errors does not introduce a significant bias into the end results of a calculation. At present, it is not possible to give a fully quantitative answer to this question. The accumulation of rounding errors is a complex phenomenon and, since the underlying equations of motion are known to be chaotic, the potential for introducing large uncontrolled errors is great . The best one can do is to ensure that the implementation of each algorithmic component is as close to reversible as practical and that the accumulation of errors grow in the expected way and so remain under control.
We study the reversibility of gauge and momentum update components separately.
### A Gauge update
The gauge update involves the process of exponentiating the conjugate momenta on all lattice links . One wishes to verify here that
* the exponentiation of the momenta does produces a suitable unitary matrix;
* the exponentiation of the momenta is reversible in the sense that
$$\mathrm{exp}(i\pi _\mu (x)\delta \tau )=\mathrm{exp}(i\pi _\mu (x)\delta \tau )^{}.$$
(26)
To check these properties, we studied
$`\mathrm{\Delta }\mathrm{Unit}`$ $`=`$ $`\underset{x,\mu ,a,b}{\mathrm{max}}\left|(\mathrm{exp}(i\pi _\mu (x)\delta \tau )\mathrm{exp}(i\pi _\mu (x)\delta \tau )^{}1)_{ab}\right|`$ (27)
$`\mathrm{\Delta }\mathrm{Rev}`$ $`=`$ $`\underset{x,\mu ,a,b}{\mathrm{max}}\left|(\mathrm{exp}(i\pi _\mu (x)\delta \tau )\mathrm{exp}(i\pi _\mu (x)\delta \tau )^{})_{ab}\right|,`$ (28)
where $`x`$, $`\mu `$, $`a`$ and $`b`$ are site, direction and colour indices respectively. These observables measure the maximum violations of unitarity and hermiticity on a given lattice.
In tests of the gauge field update reversibility, we used quenched lattices with $`V=4^4`$ sites at $`\beta =5.4`$. For the MD evolution we used $`\tau =1`$ and $`\delta \tau =\frac{1}{10}`$. The maximum values of both $`\mathrm{\Delta }\mathrm{Unit}`$ and $`\mathrm{\Delta }\mathrm{Rev}`$ along a molecular dynamics trajectory were found to be
$$\underset{\mathrm{traj}}{\mathrm{max}}\mathrm{\Delta }\mathrm{Unit}=\underset{\mathrm{traj}}{\mathrm{max}}\mathrm{\Delta }\mathrm{Rev}=0.59604635\times 10^7\frac{1}{2}ϵ_{SP},$$
(29)
where $`ϵ_{SP}`$ is the single precision unit of least precision. The fact that the maxima of the metrics agree to 8 decimal places may seem surprising at first, but becomes less mysterious when we recall that we are working at the limits of single precision, where the discrete nature of floating point numbers on a computer becomes apparent. Hence, there is only a discrete set of numbers available that the metrics can take of which the figure quoted above is one.
### B Momentum update
In the momentum update there are two possible sources of reversibility violation. The first is a lack of associativity in the addition $`p(\tau +\delta \tau )=p(\tau )+F(U)\delta \tau `$ required in the update step. The second arises in the computation of the force $`F`$. However, when performing a momentum update forward in time for a step $`\delta \tau `$ followed immediately by a momentum step backwards in time for $`\delta \tau `$, (with no gauge field update in between) the gauge fields, and hence the force, should remain unchanged. Thus, reversibility due to lack of associativity in the addition can be isolated.
Consider a test where one starts with a set of fields $`(U,\pi ,\varphi )`$. First the momentum fields are updated forward in time for a timestep $`\delta \tau `$ to produce fields $`(U,\pi ^{},\varphi )`$ and then a momentum update is performed backwards in timeIn practice this is done by flipping the signs of all the momenta, integrating the equations of motion forward in time and flipping the signs of the momenta again. to produce fields $`(U,\pi ^{\prime \prime },\varphi )`$. We use the same value of the force $`F`$ for both of the updates. One can then define the quantity
$$\mathrm{\Delta }\pi _\mu ^i(x)=\pi _\mu ^i(x)^{\prime \prime }\pi _\mu ^i(x)$$
(30)
as a measure of the reversibility violation incurred by the momentum update step. To improve statistics, one may repeat this several times, in each case using a new set of initial momenta drawn from a Gaussian distribution.
In the numerical tests, we started from some initial gauge field configuration and performed MD in the ordinary sense. Before every momentum update, we performed 100 forward–backward steps with newly drawn momenta in each case. After the test was completed, we restored the original momenta from the end of the last gauge update step and allowed the MD to continue. Thus we obtained an estimate of $`\mathrm{\Delta }\pi _\mu ^i(x)`$, the average reversibility violation due to lack of associativity in the addition. At the end of the complete trajectory, the resulting data was split into 8 sets, one corresponding to each of the Lie algebra indices $`i`$. The data in each set was histogrammed to obtain the distribution of the average reversibility violation for each momentum component.
The results of these momentum update tests are shown in figure 1. We show the histograms of all 8 momentum components. The errors on the data points are small and, to aid clarity, are not displayed. The lattice volume used for these tests was $`V=4^3\times 8`$ sites and physical parameters were $`\beta =5.2`$, $`c=0`$ and $`\kappa =0.1360`$. We performed the tests following each gauge field update along a trajectory consisting of 10 timesteps, each of length $`\delta \tau =0.1`$. We used 500 bins for each momentum component in the histograms. The histogramming process itself was carried out in double precision, allowing us to resolve reversibility violations of $`O(10^1ϵ_{SP})`$.
Figure 1 shows that the distribution of reversibility violations forms a very narrow, apparently symmetric, distribution around $`0`$ with a width that is of $`O(10^1ϵ_{SP})`$. We conclude that the momentum update step in itself is as reversible as it is possible to attain. The apparent symmetry of the distribution may possibly be used to make more general statements about reversibility and area preservation holding stochastically .
### C Reversibility of the force computation
Since gauge fields are unchanged along a momentum update and computation of the force due to gauge fields is an entirely deterministic process, one expects that the force computation will be reversible. However, the pseudofermion contribution to the force requires solution of a linear equations, so further scrutiny is required.
It has been pointed out that the solution process should be reversible, provided that a time symmetric initial guess vector (such as a zero vector or vector with random components) is used to start the solution process. This makes it tempting to carry out such solves with a large target residue $`r`$, and hence save on the computational workload. We discuss this further in sections III G and V.
Another commonly used solver strategy is to use the solution from the force computation of the previous momentum update as an initial guess. This, and variants which use a more elaborate extrapolation of previous solutions, may reduce the computational workload but are inherently non-reversible unless the solutions are effectively exact.
### D Global Reversibility Violations
Having discussed the sources of reversibility violation at a microscopic level, we now turn to the problem of their global accumulation. Consider an MD trajectory with initial fields $`(U,\pi )`$ and a set of pseudofermion fields $`\varphi `$. The latter remain unchanged along an MD trajectory. Suppose we perform an MD trajectory forward to obtain fields $`(U^{},\pi ^{})`$, then having reversed the momenta, perform a second (backward) trajectory and a momentum flip to obtain fields $`(U^{\prime \prime },\pi ^{\prime \prime })`$. One may define the following global reversibility violation metrics:
$`\mathrm{\Delta }\delta U`$ $`=`$ $`\sqrt{{\displaystyle \underset{x,\mu ,a,b}{}}|U_\mu ^{ab}(x)^{\prime \prime }U_\mu ^{ab}(x)|^2},`$ (31)
$`\mathrm{\Delta }\delta \pi `$ $`=`$ $`\sqrt{{\displaystyle \underset{x,\mu ,i}{}}\left(\pi _\mu ^i(x)^{\prime \prime }\pi _\mu ^i(x)\right)^2},`$ (32)
$`|\mathrm{\Delta }\delta H|`$ $`=`$ $`\left|H(U^{\prime \prime },\pi ^{\prime \prime },\varphi )H(U,\pi ,\varphi )\right|.`$ (33)
It is also useful to consider these quantities suitably normalised by their respective degrees of freedom:
$$\mathrm{\Delta }\delta U_{\mathrm{d}.\mathrm{o}.\mathrm{f}}=\frac{\mathrm{\Delta }\delta U}{\sqrt{N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^U}},\mathrm{\Delta }\delta \pi _{\mathrm{d}.\mathrm{o}.\mathrm{f}}=\frac{\mathrm{\Delta }\delta \pi }{\sqrt{N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^\pi }}\text{and}|\mathrm{\Delta }\delta H|_{\mathrm{d}.\mathrm{o}.\mathrm{f}}=\frac{|\mathrm{\Delta }\delta H|}{\sqrt{N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^H}}.$$
(34)
Here $`N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^U=N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^\pi =4\times 8\times V`$ are the respective number of the gauge and momentum degrees of freedom (4 links per site and 8 $`SU(3)`$ generators) and $`N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^H`$ is the number of degrees of freedom involved in computing the Hamiltonian H. In the quenched approximation $`N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^H=N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^U+N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^\pi `$. When dynamical fermions are included, there is an additional factor from the fermions of $`N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^f=24\times V`$ (3 colour and 4 Dirac complex components per site). In the even–odd preconditioned systems, half of the $`N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^f`$ degrees of freedom are represented in the pseudofermion vectors and the remainder absorbed into computing $`\mathrm{Tr}\mathrm{ln}A`$ on sites of the opposite parity.
We also study $`\frac{|\mathrm{\Delta }\delta H|}{|\delta H|}`$, where
$$\delta H=H(U^{},\pi ^{})H(U,\pi ).$$
(35)
This is a measure of the relative error in our energy calculations and is related to the accuracy of the acceptance probability. One would like this relative error to be quite small, certainly no more than a few percent.
### E Volume scaling of global reversibility metrics
According to their definitions, $`\mathrm{\Delta }\delta U`$ and $`\mathrm{\Delta }\delta \pi `$ should scale as $`O(\sqrt{V})`$, since the metrics require the summation of $`O(V)`$ positive definite quantities. We therefore expect that the corresponding normalised (per degree of freedom) metrics should volume-independent. For $`|\mathrm{\Delta }\delta H|`$, the summation involves numbers which are not positive-definite, and one might expect some cancellation. If the numbers are truly random, the cancellations between the terms can be modelled as a random walk and one would expect the sum to scale as $`O(\sqrt{V})`$. Hence one would expect $`|\mathrm{\Delta }\delta H|_{\mathrm{d}.\mathrm{o}.\mathrm{f}}`$ to be independent of the system volume in a manner similar to the $`\mathrm{\Delta }\delta U_{\mathrm{d}.\mathrm{o}.\mathrm{f}}`$ and $`\mathrm{\Delta }\delta \pi `$ metrics.
To satisfy ourselves further that our simulation code is performing as well as can be expected, we carried out reversed trajectories (as described in the definition of the metrics) in the quenched approximation with lattices of different volumes. In each case, we used a single configuration as the starting gauge field for the test and the momentum field was drawn randomly from a heat bath. The trajectory length was $`\tau =1`$ and the length of the timestep was $`\delta \tau =\frac{1}{180}`$. We used $`\beta =5.4`$ and lattices of volume
$$V\{4^4,8^4,10^3\times 16,16^3\times 32\}.$$
(36)
Results of these tests are shown in figure 2 where the volumes have been normalised by the smallest one ($`V_0=4^4`$). We note that the degree of freedom normalised metrics – $`\mathrm{\Delta }\delta U_{\mathrm{d}.\mathrm{o}.\mathrm{f}}`$, $`\mathrm{\Delta }\delta \pi _{\mathrm{d}.\mathrm{o}.\mathrm{f}}`$ and $`\mathrm{\Delta }\delta H_{\mathrm{d}.\mathrm{o}.\mathrm{f}}`$ – are all independent of the volume as expected. We also note that the relative error $`\frac{|\mathrm{\Delta }\delta H|}{|\delta H|}`$ is less than of order $`0.1\%`$, showing that error in computing the acceptance probability is small.
### F Accumulation of rounding errors in MD time
It has been noted by several authors that the MD equations of motion are chaotic and so effects of roundoff error are expected to grow exponentially with MD time along a trajectory. In particular, if one were to carry out reversed trajectory tests, as described in the definition of the metrics $`\mathrm{\Delta }\delta U`$ and $`\mathrm{\Delta }\delta \pi `$, these would be expected to exhibit the leading behaviour
$$\mathrm{\Delta }\delta Ue^{\nu _U\tau }\text{and}\mathrm{\Delta }\delta \pi e^{\nu _\pi \tau }$$
(37)
as a function of the MD trajectory length $`\tau `$. We use this as an operational definition of the effective leading Liapunov exponents $`\nu _U`$ and $`\nu _\pi `$. In our computations we measured only $`\nu _U`$ and, hence, in future discussion we shall drop the subscript $`U`$ and refer to it simply as $`\nu `$. We shall also refer to $`\nu `$ simply as the Liapunov exponent.
The authors of all found positive values for the Liapunov exponents in their studies. In particular it was shown in that as the solver target residue $`r`$ and MD step–size $`\delta \tau `$ were made smaller, the Liapunov exponents appeared to plateau, indicating that chaos was present in the underlying continuum equations of motion for the system and not just a feature of the numerical integration scheme.
For the leading Liapunov exponent $`\nu `$, the authors of found that this plateau came to an end at $`\delta \tau 0.6`$ in the quenched approximation and in the case of dynamical fermion simulations with sufficiently heavy quarks. Beyond this step size, the effective exponent exhibited growth. However, in the case of light quarks, this growth was found to set in significantly earlier, at $`\delta \tau 0.08`$. This sudden growth in Liapunov exponents could signal the onset of instability in the MD. The subject of integrator instabilities will be taken up in section IV.
The authors of also studied the behaviour of the Liapunov exponents as a function of the MD solver target residue $`r`$. They investigated the effects of increasing $`r`$ (using a time symmetric start) as a possible means of improving computational efficiency. Their data indicated a sudden growth in Liapunov exponent as $`r`$ is increased beyond a critical value. The data covered a limited range of $`r`$, and had large statistical errors. However, the sudden apparent growth of the Liapunov exponent coincides with a dramatic drop in acceptance rate, suggesting again that the integrator has become unstable.
### G Tuning the solver target residual
The results of motivated us to measure the Liapunov exponents of our simulations while varying the target residue of a comparatively large volume system, with comparatively light quarks such as those in current production runs.
For the determination of Liapunov exponents, we used 10 configurations taken from one of our large data sets. The lattice volume used was $`V=16^3\times 32`$ and the physical parameters were $`\beta =5.2`$, $`c=2.0171`$ and $`\kappa =0.1355`$. The value of the clover coefficient was calculated using the formula determined by the Alpha collaboration . These parameters correspond to pseudoscalar to vector mass ratio of $`\frac{m_\pi }{m_\rho }0.6`$ and a lattice spacing of $`a=0.097fm`$ where the physical lattice spacing has been determined using the observable $`r_0`$. By current standards, the dynamical fermions are relatively light.
Using the 10 starting configurations, for a given value of $`r`$ we carried out reversed MD trajectories of varying length $`\tau `$ with a constant step–size of $`\delta \tau =\frac{1}{180}`$. This value for $`\delta \tau `$ was the one used in the production of the dataset from which our 10 sample configurations were taken. Our MD solver strategy was to employ a two stage BiCGStab solution to compute the quantity $`X`$ of (19) followed by a restarted CG solution. Hence the target residue used was the accumulated target residue for the CG solver as described in section II E. The target residues used ranged from $`r=10^7`$ to $`r=10^4`$. The smallest of these is near the limit of what may be achieved in a single precision (32bit) computation.
In each test we measured $`\mathrm{\Delta }\delta U`$, $`|\delta H|`$ and $`N_{\mathrm{iters}}`$, where $`N_{\mathrm{iters}}`$ was the total number of solver iterations carried out in both the BiCGStab and CG solves averaged over the forward and reverse trajectories. For each combination of parameters, we also calculated the Metropolis acceptance probability $`P_{\mathrm{acc}}`$.
To evaluate the savings (or losses) in computational cost we defined the cost metric
$$\text{Cost}=\frac{N_{\mathrm{iters}}}{P_{\mathrm{acc}}}.$$
(38)
This heuristic measure reflects the fact that a large number of iterations along an MD trajectory implies high computational cost, as does a low Metropolis acceptance rate. We note that an absolute measure of cost should also take into account the autocorrelation time of the ensemble produced by an HMC computation. Since we are unable to control or measure this quantity on a sample of 10 configurations, we disregard autocorrelation effects in this study where we are interested in the relative cost with different choices of simulation parameters.
Figure 3 shows fits used to extract the (effective) Liapunov exponents. The system is clearly chaotic as $`\mathrm{ln}\mathrm{\Delta }\delta U`$ has a significant positive slope as a function of $`\tau `$. Even with only 10 configurations, the signal for the Liapunov exponents is good except for the cases when $`r=5\times 10^6`$ and when $`r=10^5`$. The data for these latter parameter values seem to show a marked break at $`\tau 0.6`$ and indeed, it was not possible to establish a consistent value of the Liapunov exponent for these two values of $`r`$.
In figure 4 we show $`\delta H`$, the energy change along an MD trajectory averaged over 10 configurations as a function of trajectory length $`\tau `$. One can clearly distinguish three different types of behaviour for $`\delta H`$ depending on the target MD residual $`r`$. For values of $`r<5\times 10^6`$, $`\delta H`$ shows an oscillatory behaviour with $`\tau `$, whereas for $`r>10^5`$ $`\delta H`$ diverges with increasing $`\tau `$, resulting in a corresponding exponential drop in acceptance probability. It is interesting to note that this change in the behaviour of $`\delta H`$ occurs at the value of $`r`$ where the data in figure 3 also show a change.
A summary of results for tuning the solver residue is shown in figure 5. The bottom panel shows the Liapunov exponents $`\nu `$. For each value of $`r`$ we made several determinations of $`\nu `$ by fitting to different ranges of $`\tau `$ in figure 3. We note that the results of these different fits are consistent with each other except for the values of $`r=5\times 10^6`$ and $`r=10^5`$ corresponding to the “break” evident in figure 3.
We note that, overall, the Liapunov exponents appear to show a slow growth with $`r`$. There is no evidence of a plateau as $`r`$ is reduced to $`r=10^7`$. This implies that this manifestation of chaos in the system is not due to the underlying equations of motion, but to the integrator. The behaviour of the exponents near $`r=10^5`$ may perhaps be interpreted as the effect of the integrator changing from being stable to being unstable.
The second panel in figure 5 shows the average acceptance rate $`P_{\mathrm{acc}}`$ for trajectories of length $`\tau 1`$. The acceptance shows a rapid drop for $`r>10^5`$, which is due to the divergent behaviour of $`\delta H`$ for values of $`r`$ in this region. The rapid drop in acceptance rate results in a huge growth in the cost of the algorithm as shown in the third panel of figure 5 where we display the cost metric (38) normalised by its value for the simulation with $`r=10^7`$.
In the top panel of figure 5 we show an enlarged view of the cost function for values of $`r<10^5`$. The cost metrics for values of $`r10^5`$ are too large to fit onto this enlarged plot. We note that the normalised cost has a shallow minimum when $`r=5\times 10^6`$ however at this minimum value the normalised cost has a value of about 0.75 implying a saving of only about $`25\%`$.
## IV Instability in the MD integration
The behaviour of the energy change $`\delta H`$, from oscillatory to divergent, is reminiscent of a known instability in the leapfrog algorithm when applied to the integration of the equations of motion for the simple harmonic oscillator. In this section, we review the simple harmonic oscillator analysis of and compare expectations for interacting theories with our numerical results.
### A Harmonic Oscillator
In what follows we use the notation of . Consider a single oscillator with coordinate $`\varphi `$. The corresponding Hamiltonian function is
$$H=\frac{1}{2}\left(\pi ^2+\omega ^2\varphi ^2\right),$$
(39)
where $`\omega `$ is the angular frequency of the oscillator and $`\pi `$ is the corresponding fictitious momentum.
The leapfrog update for the coordinate and momentum may be written in the form of a matrix $`𝒰_3(\delta \tau )`$ acting on the phase space vector $`(\varphi ,\pi )`$
$$𝒰_3(\delta \tau )=\left(\begin{array}{cc}1\frac{1}{2}\omega ^2\delta \tau ^2& \delta \tau \\ \omega ^2\delta \tau +\frac{1}{4}\omega ^4\delta \tau ^3& 1\frac{1}{2}\omega ^2\delta \tau ^2\end{array}\right).$$
(40)
The update matrix $`𝒰_3`$ can be parameterised as
$$𝒰_3(\delta \tau )=\left(\begin{array}{cc}\mathrm{cos}[\kappa (\delta \tau )\delta \tau ]& \frac{\mathrm{sin}[\kappa (\delta \tau )\delta \tau ]}{\rho (\delta \tau )}\\ \rho (\delta \tau )\mathrm{sin}[\kappa (\delta \tau )\delta \tau ]& \mathrm{cos}[\kappa (\delta \tau )\delta \tau ]\end{array}\right)$$
(41)
where
$`\kappa (\delta \tau )`$ $`=`$ $`{\displaystyle \frac{\mathrm{cos}^1(1\frac{1}{2}\omega ^2\delta \tau ^2)}{\delta \tau }}`$ (42)
$`\rho (\delta \tau )`$ $`=`$ $`\omega \sqrt{1{\displaystyle \frac{1}{4}}\omega ^2\delta \tau ^2}.`$ (43)
Evolution over a whole trajectory of length $`\tau `$ is then given by
$$𝒰_3(\tau )=\left(\begin{array}{cc}\mathrm{cos}[\kappa (\delta \tau )\tau ]& \frac{\mathrm{sin}[\kappa (\delta \tau )\tau ]}{\rho (\delta \tau )}\\ \rho (\delta \tau )\mathrm{sin}[\kappa (\delta \tau )\tau ]& \mathrm{cos}[\kappa (\delta \tau )\tau ]\end{array}\right).$$
(44)
The nature of the instability in the leapfrog scheme may be illustrated by examining the phase space trajectories in this system. The initial phase space vector for an oscillator released from amplitude $`A`$ is $`(\varphi (0),\pi (0))=(A,0)`$. From (44), the phase space vector at time $`\tau `$ is then given by
$$\left(\begin{array}{c}\varphi (\tau )\\ \pi (\tau )\end{array}\right)=\left(\begin{array}{c}A\mathrm{cos}[\kappa (\delta \tau )\tau ]\\ A\rho (\delta \tau )\mathrm{sin}[\kappa (\delta \tau )\tau ]\end{array}\right).$$
(45)
The phase space orbits therefore satisfy
$$\frac{\varphi ^2(\tau )}{A^2}+\frac{\pi ^2(\tau )}{A^2\rho ^2(\delta \tau )}=1.$$
(46)
It can then be seen from (43) and (46) that for $`\omega \delta \tau <2`$ the phase space trajectories are ellipticalIn the exact solution the orbits are circular, the deformation to an ellipse is an effect of the truncation error in the leapfrog scheme even in exact arithmetic., whereas for $`\omega \delta \tau >2`$ they are hyperbolic. The instability at $`\omega \delta \tau =2`$ is the abrupt transition from one class of phase space trajectories to another.
The change in energy
$$\delta H=H(\varphi (\tau ),\pi (\tau ))H(\varphi (0),\pi (0))$$
(47)
may also be computed. Using the same initial conditions,
$$\delta H=\frac{1}{8}\omega ^4A^2\delta \tau ^2\mathrm{sin}^2[\kappa (\delta \tau )\tau ].$$
(48)
When $`\omega \delta \tau <2`$, $`\kappa (\delta \tau )`$ is real and so $`\delta H`$ oscillates with increasing $`\tau `$, in a manner similar to that observed in the bottom panel of figure 4. However, when $`\omega \delta \tau >2`$, $`\kappa (\delta \tau )`$ becomes purely imaginary causing $`\delta H`$ to diverge as $`\mathrm{sinh}^2[\kappa (\delta \tau )\tau ]`$ in a manner similar to that seen in the top panel of figure 4.
### B Generalised treatment of instabilities
We now present a more general method of finding instabilities in the leapfrog algorithm and in higher order schemes of the type discussed in (see section II C) when applied to the case of a harmonic oscillator.
Consider an initial phase space vector $`(\varphi ,\pi )`$ of the harmonic oscillator. This is to be evolved through phase space by the leapfrog matrix $`𝒰_3(\delta \tau )`$ of (40). The area preservation property of the integrator implies that $`det(𝒰_3(\delta \tau ))=1`$. All components of $`𝒰_3(\delta \tau )`$ are real, implying that $`\mathrm{Tr}𝒰_3`$ is also real.
If
$$\lambda _1=u_1+iv_1\text{and}\lambda _2=u_2+iv_2$$
(49)
are the two eigenvalues of $`𝒰_3(\delta \tau )`$, the previous conditions on the trace and the determinant (area preservation) can then be shown to imply that
$$v_1=v_2\text{and}u_1v_2+u_2v_1=0.$$
(50)
We conclude that either:
1. $`u_1=u_2`$ or
2. $`v_1=v_2=0`$
In case 1), the determinant condition ($`\lambda _1\lambda _2=1`$) implies that $`u_1^2+v_1^2=1`$. The eigenvalues have magnitude unity: $`\lambda _{1,2}=e^{\pm i\theta }`$ with $`\theta `$ real, and the update matrices $`𝒰_3(\delta \tau )`$ and $`𝒰_3(\tau )`$ ($`=𝒰_3^{N_{\mathrm{MD}}}(\delta \tau )`$) give stable elliptical trajectories in phase space.
In case 2), by the same condition on the determinant, we have that $`\lambda _1=\eta `$ and $`\lambda _2=\frac{1}{\eta }`$ for some real $`\eta 1`$. On raising $`\lambda _1`$ or $`\lambda _2`$ to the power $`N_{\mathrm{MD}}`$, one of the eigenvalues of $`𝒰_3(\tau )`$ will show an exponential divergence with $`N_{\mathrm{MD}}`$. This implies unstable behaviour in the integrator.
The condition for the onset of instability is that the eigenvalues change from being complex to real. This information can be deduced from the discriminant of the characteristic polynomial of the update matrix $`𝒰_3(\delta \tau )`$. The onset of instability occurs as the discriminant changes sign from negative to positive.
For the leapfrog method, the discriminant is given by
$$D_3=(\omega \delta \tau )^2(\omega \delta \tau 2)(\omega \delta \tau +2).$$
(51)
We note that for $`0<\omega \delta \tau <2`$, the discriminant is negative indicating a stable integrator, whereas for $`\omega \delta \tau >2`$ the discriminant is positive implying an unstable integrator in line with the previous discussion.
### C Instability in Higher Order Schemes
Consider the 5th order scheme of Campostrini and Rossi . This can be constructed from three leapfrog integration steps as
$$𝒰_5(\delta \tau )=𝒰_3(\delta \tau _1)𝒰_3(\delta \tau _2)𝒰_3(\delta \tau _1)$$
(52)
with $`\delta \tau _1=\frac{\delta \tau }{22^{1/3}}`$ and $`\delta \tau _2=\frac{2^{1/3}\delta \tau }{22^{1/3}}`$. This corresponds to $`n=3`$ and $`i=1`$ in (8) and (9).
The discriminant $`D_5`$ is a 12th order polynomial in $`\omega \delta \tau `$ which can easily be found using an algebraic package such as Maple. It is not reproduced here but plotted in figure 6. The nonnegative roots of the $`D_5=0`$ are found to be
$$\omega \delta \tau \{0,\sqrt{126\sqrt[3]{4}}\}.$$
(53)
To three decimal places, the positive root is at $`1.573`$. The discriminant is negative for $`0<\omega \delta \tau <1.573`$ indicating stable behaviour and is positive for $`\omega \delta \tau >1.573`$ for the region where the integrator is unstable.
It is interesting to note that, for the central leapfrog update matrix $`𝒰_3(\delta \tau _2)`$ in the 5th order scheme to become unstable on its own, requires that $`\omega \delta \tau _2=2`$. This implies that this central step should go unstable when
$$\omega \delta \tau =2\frac{\left(22^{1/3}\right)}{2^{1/3}}1.175.$$
(54)
This suggests that, although the central update itself becomes unstable at $`\delta \tau =1.175`$, the other two updates in the scheme stabilize the system until $`\delta \tau 1.57`$.
Following a similar calculation, it can be shown that the discriminant $`D_7`$ of the characteristic polynomial for the update matrix of the 7th order scheme ($`n=5`$, $`i=1`$) has roots at
$$\omega \delta \tau \{0,1.595,1.822,1.869\}$$
(55)
with $`D_7`$ being negative in the intervals $`D_7(0,1.595)`$ and $`D_7(1.822,1.869)`$ indicating two domains of stability. The discriminant is positive for $`D_7(1.592,1.822)`$ and for $`D_7>1.869`$. For the longest constituent 5th order update to go unstable in this scheme requires that $`\omega \delta \tau >1.166`$.
Hence we see that, for the case of the simple harmonic oscillator at least, higher order integration schemes do not help cure the problem of instabilities. Indeed, they become unstable at even smaller values of $`\omega \delta \tau `$ than the simplest leapfrog method.
### D Hypothesis for interacting field theories
Edwards, Horváth and Kennedy advanced the hypothesis that, since the high frequency modes of an asymptotically free field theory can be considered as a collection of weakly coupled oscillator modes, the instability just described in the harmonic oscillator system will also be present for interacting field theories. The onset of the instability will be caused by the mode with highest frequency $`\omega _{\mathrm{max}}`$, when $`\omega _{\mathrm{max}}\delta \tau =2`$. For a single oscillator mode, the onset of instability is abrupt. In the case of an interacting theory, one would expect the effects of the interactions to smooth out this transition.
It is argued in that the instability in lattice QCD with dynamical fermions can be likened to that of a collection of oscillator modes of the sort just described. When applying leapfrog integration to this system, the rôle of $`\omega ^2\varphi `$ in the harmonic oscillator example is played by the MD force $`F_\mu (x)`$. This force can be written as a sum of contributions from the gauge and fermionic pieces of the action as $`F_\mu (x)=F_\mu ^\mathrm{g}(x)+F_\mu ^\mathrm{f}(x)`$, where the subscripts $`g`$ and $`f`$ indicate the gauge and fermionic components of the force respectively.
The fermion force is expected to be proportional to $`m_f^\alpha `$, where $`m_f`$ is the mass of the lightest species of dynamical fermion and $`\alpha `$ is some *negative* parameter. In the case of Wilson (and Clover) fermions the mass in lattice units is defined as
$$am_f=\frac{1}{2}\left(\frac{1}{\kappa }\frac{1}{\kappa _c}\right)$$
(56)
where $`\kappa `$ now stands for the Wilson hopping parameter, and $`\kappa _c`$ is the critical value corresponding to $`m_f=0`$. It is argued that the highest frequency mode (with frequency $`\omega _{\mathrm{max}}`$) is proportional to the fermion force which, in turn, is expected to be proportional to $`m_f^\alpha `$, and thus as $`\kappa \kappa _c`$ ($`m_f0`$), the fermion force will diverge and hence the critical value of $`\delta \tau `$ will decrease. In the following, we evaluate numerical evidence for the validity of this hypothesis.
### E Studies of the force
The forces used in the momentum update belong to the Lie algebra $`su(3)`$. We define the 2–norm $`F`$ in the same manner as for $`\mathrm{\Delta }\delta \pi `$:
$$F=\sqrt{\underset{x,\mu ,i}{}\left(F_\mu ^i(x)\right)^2}.$$
(57)
Again, we can define the 2–norm suitably normalised by the relevant degrees of freedom:
$$F_g_{\mathrm{d}.\mathrm{o}.\mathrm{f}}=\frac{F_g}{\sqrt{N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^U}}\text{and}F_f_{\mathrm{d}.\mathrm{o}.\mathrm{f}}=\frac{F_f}{\sqrt{N_{\mathrm{d}.\mathrm{o}.\mathrm{f}}^f}}$$
(58)
where the subscripts $`g`$ and $`f`$ indicate gauge and fermionic forces respectively.
We can also define an $`\mathrm{}`$–norm for the forces:
$$F_{\mathrm{}}=\underset{x,\mu ,i}{\mathrm{max}}\left|F_\mu ^i(x)\right|.$$
(59)
The $`\mathrm{}`$–norm then is the force component with the maximum magnitude over the lattice and so can be likened to the force mode with the highest frequency, proportional to $`\omega _{\mathrm{max}}^2`$, in the analogous collection of weakly coupled harmonic oscillators. The (degree of freedom) averaged 2–norm on the other hand can be likened to the average frequency–squared of the analogous set of harmonic oscillators.
In our studies we computed the magnitude of the forces at all timesteps of an MD trajectory starting from a single gauge configuration chosen from the same 10 configurations described in section III G (with volume lattice $`V=16^3\times 32`$ sites, and production parameters $`\beta =5.2`$, $`c=2.0171`$, $`\kappa =0.1355`$).
In the first set of tests, we attempted to investigate how the fermion force behaves with the quark mass. We performed MD trajectories consisting of $`N_{MD}=175`$ steps of length $`\delta \tau =\frac{1}{180}`$ for several values of the hopping parameter $`\kappa `$. We measured the norms of the gauge and fermion forces on each timestep. The MD solver target residue was set at $`r=10^6`$. Error bars for the average value of the force were computed by bootstrapping the 175 samples.
It could be argued that a configuration that has been produced in an ensemble equilibrated at some value of $`\kappa `$, will have very small statistical weight at a different value of $`\kappa `$. However, our aim was not to study equilibrium properties of the ensemble, but to test the properties of algorithm components as a function of the external parameter $`\kappa `$.
The average value of $`\kappa _c`$, the critical value of $`\kappa `$ corresponding to massless fermions, is known from separate spectroscopy studies for the ensemble from which the configurations were drawn. It is approximately $`0.1363`$ . Thus, we were able to associate a value of the lattice fermion mass $`am_f`$ with every value of $`\kappa `$ used in our tests through the formula:
$$am_f=\frac{1}{2}\left(\frac{1}{\kappa }\frac{1}{\kappa _c}\right).$$
(60)
Since we expect the fermion mass to vary in some inverse relation to the norm of the force , we attempted to fit the results of our tests with the form
$$F=A(am_f)^\alpha =A\left(\frac{1}{2\kappa }\frac{1}{2\kappa _c}\right)^\alpha ,$$
(61)
where the parameters of the fit were $`A`$, $`\kappa _c`$ and $`\alpha `$.
Results of this test are shown in figure 7. We show both the fits made to the $`\mathrm{}`$–norm and the (degree of freedom) averaged 2–norm of the force. We can see that good fits can be made, which reproduce $`\kappa _c`$ from the spectroscopic studies and that $`\alpha `$ is negative indicating that the magnitudes of the norms do indeed vary in an inverse manner with the fermion mass. The fact that the value of $`\kappa _c`$ is well reproduced and that $`\alpha `$ is negative in sign both lend support to the hypothesis of .
### F Dependence on $`\delta \tau `$ and $`\kappa `$
To investigate further the onset of instability, we computed the averaged forces and $`\delta H`$ along an MD trajectory using the same starting configurations as before. However, this time we varied the MD step size $`\delta \tau `$. The number of steps taken along the trajectory was adjusted to keep the trajectory length constant at $`\tau =175/180`$. The results are plotted in figure 8. From the growth of $`\delta H`$ evident in the plot, one can see that the instability sets in between $`\delta \tau =0.0105`$ and $`\delta \tau =0.0110`$. We can also see that the rapid growth of $`\delta H`$ is accompanied by a growth in the fermionic forces in the system (in both norms) and that the $`\mathrm{}`$–norm of the force appears to grow more rapidly than the degree of freedom averaged 2–norm. This latter behaviour suggests that the onset of instability is driven by a few unstable fermion modes, again in line with the above hypothesis.
In a further investigation of the MD forces, we carried out MD trajectories using the same initial gauge configuration as before, this time varying $`\kappa `$ for two separate values of the step size. The values of the step size were $`\delta \tau =0.010`$ and $`\delta \tau =0.012`$ corresponding to stable and unstable MD at $`\kappa =0.1355`$ respectively, as discussed above.
We show the $`\mathrm{}`$–norms of the gauge and fermion forces in figure 9. This shows that the simulation which was unstable at $`\kappa =0.1355`$ has become stable as $`\kappa `$ is reduced. Once again this seems in line with the hypothesis that the onset of the instability is a function of the combination of the fermionic forces (controlled by $`\kappa `$) and the stepsize $`\delta \tau `$. Recall that the relevant parameter for the SHO was $`\omega \delta \tau `$.
Overall, our studies of the MD forces lend support to the hypothesis that the instability is driven by the $`F\delta \tau `$ term in the momentum update step of the leapfrog algorithm. Since the fermionic force diverges in some inverse relation with the fermion mass, we expect the maximum safe stepsize $`\delta \tau `$ to decrease as the fermion mass is decreased ($`\kappa `$ is increased). Also, having observed a faster rise in the $`\mathrm{}`$–norm of the fermionic force than in the degree of freedom averaged 2–norm, we infer that the instability is driven by a comparatively small number of unstable fermionic modes.
## V Tuning the stepsize and the solver residue
The above conjecture, if correct, can serve to explain the tuning results described in section III G. By increasing the solver residue $`r`$, we are modifying the fermionic force which could then drive the MD integrator unstable. In order to investigate these possibilities, we have carried out a second tuning exercise this, time varying both the step size $`\delta \tau `$ and the solver target residue $`r`$.
We used the 10 configurations used when tuning $`r`$ alone in section III G. Since at this point we were not computing Liapunov exponents, our tests consisted of single MD trajectories in one direction only. For each value of $`\delta \tau `$, we chose the number of steps along the trajectory so as to maintain a constant trajectory length of $`\tau =175/180`$. We also carried out a test with a target residue of $`r=10^9`$ using double precision (64bit) floating point numbers, whereas all other tests used single precision. For each combination of algorithmic parameters, we measured the energy change $`\delta H`$, the corresponding acceptance probability $`P_{\mathrm{acc}}`$ and the cost function of (38).
The results of this tuning exercise are shown in figure 10. First we see in the bottom panel ($`r=10^9`$ symbols) that using double precision does not alleviate the problem of instability. The calculation in double precision appears to become unstable at a similar value of the step size as does that in single precision. Second, we see from the data for $`r=5\times 10^5`$ that, if the solver target residue is too large, one cannot achieve values of $`\delta H`$ of $`O(1)`$, even if $`\delta \tau `$ is made very small.
For our simulations, we are able to achieve non–zero acceptance rates when $`\delta \tau <0.0075`$ and when $`r10^5`$. For parameter values smaller than these, we can attempt to tune our simulation for maximum performance. The top two panels of figure 10 show the variation of the cost function. In this case, the cost function is normalised by its value when $`r=10^6`$ and $`\delta \tau =0.0055`$. These were the parameters used in the production of the dataset from which the configurations were taken. We see that either by tuning the solver residue $`r`$ or the MD step size $`\delta \tau `$, the maximum gain we could make in the cost function is about 25%.
## VI Conclusions and Discussion
### A Stability
We have shown that, for the physical parameters used in our production simulations, the molecular dynamics integrator used becomes unstable at $`\delta \tau 0.01`$ for all studied values of $`r`$, and also for any realistic value of $`\delta \tau `$ when $`r`$ was increased above $`rO(10^5)`$. We identify this instability with the one studied in free field theory for the frequency–step-size combination $`\omega _{\mathrm{max}}\delta \tau =2`$. We have studied numerically the fermion force and found that its behaviour is not inconsistent with the hypothesis of (motivated by free field theory) that the force should grow large as $`\kappa \kappa _c`$. We suppose that a critical value exists for $`F\delta \tau `$ when the leapfrog integrator becomes unstable.
Reducing the value of the MD residual results in an increasingly inaccurate force calculation. If as a result $`F`$ is too large, one may need an extremely small step-size to keep the integrator stable. We found that, for $`r=5\times 10^5`$ at our parameters, one would need a step-size much smaller than $`\delta \tau =0.001`$. (c.f. figure 10).
On the safe side of these limits, one may attempt to tune the algorithm. However, our studies show that on this volume and with these physical parameters, tuning $`\delta \tau `$ and/or $`r`$ is unlikely to produce significant performance gains. We note that it appears entirely safe to carry out computations in single precision in the safer region of parameter space. However, as $`\kappa \kappa _c`$, it may be that the upper limit on $`r`$ decreases beyond the limit of single precision. Alternatively, as the condition number of the fermion matrix increases with increasing $`\kappa `$, the number of iterations in the solver for fixed $`r`$ will increase. This may cause rounding errors to accumulate so that the target residual $`r`$ may not be reached. However, in this latter case, it is only the solve itself that needs to be done in double precision, or restarted in single precision.
### B Higher order integration Schemes
We have demonstrated that, at least for the case of a simple harmonic oscillator, the 5th and 7th order schemes of are not immune to instabilities. We expect that this situation will persist for even higher order schemes of this sort. The source of the problem is that, at the bottom level, these schemes are constructed out of simple leapfrog updates. For any given step–size $`\delta \tau `$ in an integration scheme of order $`n+3`$, there will always be a sub update of order $`n+1`$ which will have a stepsize $`\delta \tau _2>\delta \tau `$. This sub–update, or one of its constituent sub–updates, may eventually drive the whole integration scheme unstable, although the other sub-updates may act as a stabilizing factor at first. We note that, in our harmonic oscillator examples, the smallest positive critical value of $`\omega \delta \tau `$ was always smaller for the higher order integrators than for the leapfrog, indicating that the instability problem is actually worse for the higher order methods.
As the source of the instability appears to come from the fermionic part of the force, we anticipate that a scheme of the type advocated in would not assist avoiding the instability either, as it attempts to improve the truncation error by performing more gauge updates. While this may drive down the truncation error, it does nothing about the problem in the fermionic update.
### C Reversibility
Reversibility itself seems not to be strongly affected by changing $`r`$. The Liapunov exponents of the system seem to show a slow rise before the instability sets in. In the region of transition from stability to instability, the Liapunov exponents are difficult to determine. One might speculate that this behaviour reflects a transition from the Liapunov exponent characterising the underlying continuous equations of motion to that characterising the unstable numerical integrator.
### D Summary
We have investigated the stability and reversibility of the HMC algorithm with two flavours of light dynamical fermions on large lattices as a function of the MD step size $`\delta \tau `$ and the MD target solver residue $`r`$. We have found upper limits on both of these for a fixed set of physical parameters. Beyond these limits, the leapfrog integrator becomes unstable and one cannot carry out a simulation programme, irrespective of the precision of the floating point numbers which one uses. On the safe side of the limits, one can carry out simulations safely in both single and double precision. Parameter tuning seems to give no major performance gains. Reversibility does not seem to be dangerously affected.
## VII Acknowledgements
We gratefully acknowledge financial support from PPARC under grant number GR/L22744. James Sexton would like to thank Hitachi Dublin Laboratory for support. We also wish to thank Z. Sroczynski for helpful discussions and for his assistance in the preparation of this paper.
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# REFERENCES
Spin Tunneling in Mn-12 Cannot be Assisted by Phonons at Millikelvin Temperatures \[Comment on PRL $`\mathrm{𝟖𝟑}`$, 416 (1999) by Bellessa et al\]
Recently, Bellessa et al in a Letter entitled “Phonon-Assisted Tunneling in High-Spin Molecules: Experimental Evidence”, reported observation of a peak in the imaginary part of the ac-susceptibility $`\chi ^{\prime \prime }`$ of Mn-12 molecular nanomagnet, that decreases strongly below 0.1 K down to 0.02 K. They explained this effect by a “phonon-induced tunneling process”. The purpose of this Comment is to show that the explanation of Bellessa et al is logically inconsistent and physically impossible. All formulas of Ref. supporting that explanation are incorrect.
The Letter of Bellessa et al reports observation of the resonant absorption of the energy of the ac magnetic field of frequency $`\omega _0`$ at the value of the transverse magnetic field, $`H_{}`$6 T, that satisfies $`\mathrm{}\omega _0=\mathrm{\Delta }(H_{}`$), where $`\mathrm{\Delta }`$ is the tunneling splitting of the ground state. The corresponding peak in the dependence of $`\chi ^{\prime \prime }`$ on $`H_{}`$ first grows as the temperature is lowered down to 0.1 K but then decreases down to 0.02 K. The authors of Ref. give the following explanation to this effect. Citing p. 122 of Abragam and Bleaney , they write $`\chi ^{\prime \prime }`$ as
$$\mathrm{Eq}.(3):\chi ^{\prime \prime }=CN\mathrm{\Delta }^2f(\omega )T_2\mathrm{tanh}(\mathrm{}\omega /2k_BT).$$
According to Bellessa et al, “$`C`$ is a constant, $`N`$ is the number of spins, and $`\mathrm{\Delta }`$ is the magnetic dipole matrix element between the two states of the fundamental dublet. The shape function $`f(\omega )`$ is a Lorentzian function of the applied frequency $`\omega `$, the resonance frequency $`\omega _0`$ (which depends on the applied magnetic field), and the relaxation time $`T_2`$ describing the linewidth.” Right after that statement the authors of Ref. go on saying “We explain our effect by assuming that the tunneling rate $`\mathrm{\Delta }`$ in Eq. (3) is induced by a two-phonon process: the magnetic moment makes a transition from $`+|\psi >`$ to $`|\psi >`$ (or vice versa) and absorbs (or emits) a quantum $`\mathrm{}\omega `$ from the ac magnetic field only if a phonon of angular frequency $`\omega `$ is absorbed and then reemitted after the transition. This process is quite similar to the Raman process, except that the frequencies of the two phonons are the same”. To account for this effect, Bellessa et al simply insert into Eq. (3) $`n(n+1)`$, where $`n=n(\omega /T)`$ is the phonon occupation number. This gives them the desired decrease of $`\chi ^{\prime \prime }`$ at low temperature. The absence of this effect above 0.1 K is blamed on the unknown temperature dependence of $`T_2`$.
To begin with, how can $`\mathrm{\Delta }`$ be “the magnetic dipole matrix element” and “the tunneling splitting” at the same time? According to the trivial calculation of Abragam and Bleaney , Eq. (3) must contain the dipole matrix element $`|\mu |^2=(g\mu _BS)^2`$ ($`S=10`$ being the spin of the Mn-12 molecule) instead of $`\mathrm{\Delta }^2`$. This equation describes solely the resonant absorption by Mn-12 molecules of photons of frequency $`\omega `$ generated by the ac-field. Why should that process, in the millikelvin range, be accompanied by the absorption and re-emission of phonons? If it was true, Eq. (3), besides the factor $`n(n+1)`$, must have been multiplied by the fourth power of ratio of the matrix element of spin-phonon coupling to the Debye temperature, $`|V_{sph}/(k_B\mathrm{\Theta }_D)|^4`$, which is a very small number. In addition, it would have been multiplied by the phase volume of the two phonons, which is nearly zero if “the frequency is the same”. “Same frequency”, in fact, has nothing to do with the Raman process of emission and absorption of two real phonons satisfying $`\mathrm{}(\omega _1\omega _2)=\mathrm{\Delta }`$, with $`\mathrm{}\omega _1\mathrm{}\omega _2k_BT`$. This latter process should be absolutely negligible in the millikelvin range as compared to the direct photon absorption given by Eq. (3). The suggestion of Ref. is, therefore, total absurd.
It is not the purpose of this Comment to speculate why the resonance value of $`\chi ^{\prime \prime }`$ goes down below 0.1 K. Nevertheless, assuming that the experiment is correct, we will suggest two possibilities. The first is the onset of magnetic ordering in a crystal of Mn-12 clusters due to magnetic dipole interaction between the clusters . The second is the ordering of nuclear spins of Mn atoms inside the cluster . Both types of ordering induce an effective field acting on the cluster. The bias induced by that field drives the cluster off resonance, preventing it from tunneling. Some support to this suggestion comes from the fact that 0.1 K is the right order of magnitude for both dipole ordering temperature and nuclear ordering temperature .
This work has been supported by the NSF Grant No. DMR-9978882.
E. M. Chudnovsky and D. A. Garanin
Physics Department, CUNY Lehman College
Bedford Park Boulevard West, Bronx, NY 10468-1589
chudnov@lehman.cuny.edu
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# Anisotropic conductivity of Nd1.85Ce0.15CuO4-δ films at submillimeter wavelengths
## Abstract
The anisotropic conductivity of thin Nd<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-δ</sub> films was measured in the frequency range 8 cm$`{}_{}{}^{1}<\nu <`$ 40 cm<sup>-1</sup> and for temperatures 4 K $`<T<300`$ K. A tilted sample geometry allowed to extract both, in-plane and c-axis properties. The in-plane quasiparticle scattering rate remains unchanged as the sample becomes superconducting. The temperature dependence of the in-plane conductivity is reasonably well described using the Born limit for a d-wave superconductor. Below $`T_\mathrm{C}`$ the c-axis dielectric constant $`\epsilon _{1c}`$ changes sign at the screened c-axis plasma frequency. The temperature dependence of the c-axis conductivity closely follows the linear in T behavior within the plane.
The electron–doped superconductor Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4</sub> (NCCO) reveals a number of properties, which are rather different from other cuprate superconductors. It has long been believed that the superconductivity in this compound is characterized by a s-wave order parameter (for a review see ). However, recent experiments, including phase-sensitive tricrystal and penetration depth measurements, strongly support d-wave type symmetry.
In-plane microwave properties of NCCO have been investigated using resonator techniques . The infrared conductivity has been obtained via Kramers-Kronig analysis of reflectivity data and by thin film transmission. In contrast to a number of ab-plane experiments, there exists only little information concerning the c-axis properties of NCCO, which is explained by the typically small dimensions of the samples along the c-axis. Most experiments on c-axis dynamics were carried out using polycrystalline NCCO samples.
Recently we have demonstrated the possibility of using a tilted-sample geometry to extract the anisotropic conductivity of layered cuprates in the submillimeter frequency range . This method combines the possibilities of the quasioptical transmission geometry with the high anisotropy of NCCO which may be estimated by the resistivity ratio $`\rho _c/\rho _{ab}10^4`$ . In this paper we present the in-plane and c-axis conductivity of an oxygen reduced Nd<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-δ</sub> film (T<sub>C</sub>=16.9K, sample #A) in the submillimeter frequency range ($`8`$ cm$`{}_{}{}^{1}<\nu <40`$cm$`{}_{}{}^{1})`$ and for temperatures $`4`$ K $`<T<300`$ K. Data on an optimally doped film (sample #B, $`T_C=`$ 20.1 K) are also discussed.
The films were prepared using a two-beam laser deposition on YSZ substrates. X-ray analysis showed the c-axis orientation of the films relative to the crystallographic axes of the substrate. The YSZ substrate of the presented film (#A) was tilted from the (001) orientation by an angle $`\alpha =2.6^o\pm 0.5^o`$. Therefore, the film was also tilted by the same angle from the ideal c-axis orientation. The oxygen concentration in the deposition chamber was reduced compared to the optimum value which resulted in a lower $`T_\mathrm{C}`$ of the film (#A). The ac-susceptibility measurements revealed an onset temperature of 16.9 K and a slightly broader transition width ($`\mathrm{\Delta }T[10\%90\%]=2.0`$ K) as compared to the optimally doped film (sample #B: $`\mathrm{\Delta }T=0.9`$ K ).
The transmission experiments in the frequency range $`6`$ cm$`{}_{}{}^{1}<\nu <40`$ cm<sup>-1</sup> were carried out in a Mach-Zehnder interferometer arrangement which allows both the measurements of transmittance and phase shift. The properties of the blank substrate were determined in a separate experiment. Utilizing the Fresnel optical formulas for the complex transmission coefficient of the substrate-film system, the absolute values of the complex conductivity $`\sigma ^{}=\sigma _1+i\sigma _2`$ were determined directly from the observed spectra. Using the tilted sample geometry at different polarizations of the incident radiation it was possible to separate the conductivity at a given tilt angle into ab-plane and c-axis components. The geometry of the experiments is shown in the insets in Fig. 1.
The conductivity of a tilted sample may be calculated assuming a free-standing film of thickness $`d`$ in a uniform electromagnetic field $`Ee^{i\omega t}`$ parallel to the surface. If the film is thin compared to the penetration depth, $`d\lambda `$, then the current and field distributions may be considered to be uniform. Taking into account the charges formed at the surface, the following equations can be derived for the geometry given in the right inset of Fig. 1:
$$\{\begin{array}{c}j_a=\sigma _a[E\mathrm{cos}\alpha (s/\epsilon _0)\mathrm{sin}\alpha ]\hfill \\ j_c=\sigma _c[E\mathrm{sin}\alpha (s/\epsilon _0)\mathrm{cos}\alpha ]\hfill \end{array}.$$
(1)
Here $`j_a`$$`(j_c)`$ is the current density, $`\sigma _a`$$`(\sigma _c)`$ is the complex conductivity in the ab-plane (along the c-axis), $`s`$ is the surface charge density, $`\epsilon _0`$ is the permittivity of free space, $`\omega =2\pi \nu `$ is the angular frequency, and $`\alpha `$ is the tilt angle. An additional equation, $`i\omega s+j_a\mathrm{sin}\alpha +j_c\mathrm{cos}\alpha =0`$ , follows from the charge conservation. The effective conductivity of the film can be defined through $`\sigma _{eff}E=j_{eff}=j_x\mathrm{cos}\alpha +j_y\mathrm{sin}\alpha `$. Solving these equations one obtains:
$$\sigma _{eff}=\frac{i\epsilon _0\omega (\sigma _a\mathrm{cos}^2\alpha +\sigma _c\mathrm{sin}^2\alpha )+\sigma _a\sigma _c}{i\epsilon _0\omega +\sigma _a\mathrm{sin}^2\alpha +\sigma _c\mathrm{cos}^2\alpha }.\text{ }$$
(2)
Both $`\sigma _{eff}`$ and $`\sigma _a`$ can be determined experimentally using the geometry shown in the right and left insets of Fig. 1, respectively. Therefore Eq. (2) can easily be solved for the c-axis conductivity. Within the approximation $`\alpha \mathrm{sin}\alpha 1`$ and $`|\sigma _a||\sigma _c|`$, Eq. (2) may be simplified to:
$$\sigma _{eff}=\frac{\sigma _a(\sigma _ci\epsilon _0\omega )}{\sigma _a\alpha ^2+(\sigma _ci\epsilon _0\omega )}$$
(3)
As discussed previously, two excitations may be observed within the tilted geometry: i) a peak in the real part of the conductivity if $`Im[\sigma _a\alpha ^2+(\sigma _ci\epsilon _0\omega _1)]=0`$ which corresponds to the mixed ab-plane/c-axis excitation and ii) the longitudinal resonance if $`Im[\sigma _ci\epsilon _0\omega _0]=0`$ which corresponds to the c-axis plasma frequency. Both excitations have been detected for the optimally doped film #B. The c-axis plasma frequency can be identified in the submillimeter frequency range also for the reduced Nd<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-δ</sub> film (see below). Due to the larger tilt angle of the sample #A the frequency of the mixed resonance is shifted to $`\nu _1200`$ cm<sup>-1</sup> (compared to $`\nu _120`$ cm<sup>-1</sup> for sample #B) and occurs as a broad maximum in the effective conductivity spectra at infrared frequencies .
Fig. 1 shows the complex conductivity of the reduced Nd<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-δ</sub> film obtained as described above. The left panels represent the real (lower frame) and imaginary (upper frame) parts of the complex conductivity for currents within the CuO<sub>2</sub> plane. The real part of the in-plane conductivity $`\sigma _{1a}`$ is weakly frequency-dependent in the submillimeter frequency range at all measured temperatures. This indicates that the quasiparticle scattering rate is larger than the frequency of the experiment. Consequently, the imaginary part of the conductivity $`\sigma _{2a}`$ is nearly zero for high temperatures but starts to show distinct frequency dependence on approaching the superconducting transition temperature. From the analysis of $`\sigma _{1a}`$ and $`\sigma _{2a}`$ the quasiparticle scattering rate may be estimated using the Drude expression: $`\sigma ^{}=\sigma _{dc}/(1i\omega \tau )`$. The term $`[i/\omega +\pi \delta (\omega )/2]/[\mu _0\lambda _a^2(T)]`$ has to be added to the Drude expression for $`T<T_\mathrm{C}`$ in order to account for the superconducting condensate. Here $`\lambda _a^2`$ is the in-plane penetration depth, $`\mu _0`$ the permeability of free space. For finite frequencies the additional term influences $`\sigma _{2a}`$ only. The solid lines in the left panels of Fig. 1 were calculated using the Drude expression extended to temperatures below $`T_\mathrm{C}`$ as described above. The fits allow to estimate the quasiparticle scattering rate both, below and above $`T_\mathrm{C}`$. The results are shown in Fig. 2 (full diamonds). The scattering rate is approximately constant for $`T70`$ K and shows only a small anomaly (within experimental errors) at $`T_\mathrm{C}`$. For $`T>T_\mathrm{C}`$, $`1/\tau `$ agrees well with the infrared data of Homes et al. (open circles) and has an approximate linear temperature dependence for $`T>100`$ K. The absence of the anomalous suppression of the quasiparticle scattering is in contrast to the results on other cuprate superconductors. This possibly indicates that $`1/\tau `$ is determined by impurity scattering for $`T<100`$ K. For $`TT_C`$ the $`(1/\omega )`$ frequency dependence dominates the imaginary part of the conductivity $`\sigma _{2a}`$ which allows to estimate the low-frequency in-plane penetration depth, $`\lambda _a(6K)=0.35\mu `$. For the optimally doped sample #B we obtained $`\lambda _a(6K)=0.23\mu `$ .
The right panels of Fig. 1 represent the effective (mixed) conductivity which is described by Eq. (2). Using this equation, the c-axis conductivity may be calculated from the in-plane and mixed conductivity data. The imaginary part of the mixed conductivity crosses zero around $`\nu 20`$ cm<sup>-1</sup>. As will be seen below, this frequency corresponds to the c-axis plasma resonance.
Fig. 3 shows the conductivity $`\sigma _{1c}`$ and the dielectric constant $`\epsilon _{1c}=\sigma _{2c}/(\epsilon _0\omega )`$ of Nd<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-δ</sub> (#A) along the c-axis. The real part of the c-axis conductivity (lower frame) is approximately frequency independent within experimental accuracy and for temperatures well above $`T_\mathrm{C}`$. This behavior agrees well with the low-frequency infrared conductivity of La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>, YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.7</sub>, and Tl<sub>2</sub>BaCuO<sub>6+x</sub>. Only at the lowest temperatures $`\sigma _{1c}`$ does increase with frequency which most probably reflects the vicinity of the c-axis phonon at $`\nu 134`$ cm<sup>-1</sup> . The c-axis dielectric constant is dominated by the high-frequency (phonon) contribution and shows a weak frequency dependence at high temperatures. As the sample becomes superconducting, $`\epsilon _{1c}`$ reveals a $`(1/\omega ^2)`$ behavior, which gives an estimate of the penetration depth, $`\lambda _c(6K)=19.2\mu `$. Consequently, a zero crossing of $`\epsilon _{1c}`$ is observed around $`20`$ cm<sup>-1</sup> which corresponds to the (screened) plasma frequency $`2\pi \nu _p=c/(\lambda _c\epsilon _{\mathrm{}}^{1/2})`$ where $`\epsilon _{\mathrm{}}14`$ is the high-frequency dielectric constant and $`c`$ is the speed of light. For the sample #B we found $`\nu _p=12`$cm <sup>-1</sup> and $`\epsilon _{\mathrm{}}23`$ .
Assuming Josephson coupling between the CuO<sub>2</sub> planes, Basov et al. suggested a correlation between $`\lambda _c(0)`$ and the normal-state conductivity $`\sigma _c(T_\mathrm{C}):\mathrm{}/(\mu _0\lambda _c^2)=\pi \mathrm{\Delta }\sigma _c(T_\mathrm{C})`$. On the basis of this correlation the results on both NCCO samples give an energy gap $`2\mathrm{\Delta }30`$ cm<sup>-1</sup>. This value may be compared to $`2\mathrm{\Delta }60`$ cm<sup>-1</sup> as determined by Raman scattering.
The temperature dependence of the anisotropic conductivity, as measured at $`\nu =10`$ cm<sup>-1</sup>, is represented in Fig. 4. The lower panel of Fig. 4 shows the real and imaginary parts of the in-plane conductivity. For decreasing temperature, $`\sigma _{1a}`$ increases below room temperature, saturates between $`T100`$ K and T<sub>C</sub> and finally decreases after a slight maximum near $`T_\mathrm{C}`$. A peak near $`T_\mathrm{C}`$ observed in $`\sigma _{1a}`$ at microwave frequencies was recently reported for NCCO by Kokales et al. and interpreted as possible evidence for suppression of the quasiparticle scattering. According to Fig. 2, our data suggest a rather temperature independent scattering of quasiparticles below $`T=100`$ K. At high temperatures the imaginary part of the in-plane conductivity (lower panel of Fig. 4, open triangles) has values just above the sensitivity limit of the spectrometer. In the superconducting state $`\sigma _{2a}`$ abruptly increases reflecting the formation of the superconducting condensate.
The lower panel of Fig. 4 shows the comparison of the experimental conductivity with theoretical models. As representative examples we have taken the s-wave BCS expression, as well as Born $`(\stackrel{~}{\sigma }=0)`$ and unitary $`(\stackrel{~}{\sigma }=1)`$ limits of a d-wave superconductor. Here $`\stackrel{~}{\sigma }`$ is the cross section of the impurity scattering. All three models calculate a gap value self-consistently within the weak coupling limit and assume a temperature independent quasiparticle scattering rate of $`1/2\pi \tau =65`$ cm<sup>-1</sup>. It has to be pointed out, that real and imaginary parts of the conductivity have to be fitted simultaneously below and above $`T_\mathrm{C}`$. This condition leaves no free parameters within the models. As documented by the fit results in Fig. 4, the s-wave curve (solid line) shows the poorest agreement with the experiment. In contrast, both limits of the d-wave model describe $`\sigma _{1a}`$ reasonably well. However, the unitary limit (dotted) substantially underestimates the imaginary part $`\sigma _{2a}`$. This is probably because in this limit less spectral weight is shifted to the $`\delta `$-function at zero frequencies. The best description of $`\sigma _{2a}`$ may be obtained using an intermediate scattering cross section $`\stackrel{~}{\sigma }0.2`$. Similar results have been obtained for the optimally doped sample #B for which $`\stackrel{~}{\sigma }=0`$ gave the best description of the data.
The upper panel of Fig. 4 shows the temperature dependence of the c-axis conductivity of Nd<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-δ</sub>. Except for the absolute values, these data closely follow the temperature dependence of the in-plane conductivity. The most prominent difference is caused by the strong phonon contribution on the c-axis conductivity, evidenced by a downward shift of $`\sigma _{2c}`$. Based on the strong in-plane momentum dependence of the scattering rate and of the hopping integral, the anisotropic conductivity for high-T<sub>C</sub> cuprates was recently calculated by van der Marel, and Xiang and Hardy. Parametrizing the in-plane momentum by an angle $`\theta `$ and using $`t_c=t_{}cos^2(2\theta )`$ for the c-axis hopping integral, the c-axis conductivity was found to behave as $`\sigma _{1c}T^3`$ for not too low temperatures. The analysis of Fig. 4 shows, that $`\sigma _{1c}`$ as well as $`\sigma _{1a}`$ for NCCO depend rather linear on temperature below $`T_\mathrm{C}`$. The explanation for this behavior probably is an impurity-induced angular-independent contribution to $`t_c`$. Following the calculations described in Refs. this correction does indeed give a linear temperature dependence of the c-axis conductivity.
In conclusion, the anisotropic conductivity of Nd<sub>1.85</sub>Ce<sub>0.15</sub>CuO<sub>4-δ</sub> films has been obtained using the tilted-sample geometry in the frequency range $`8`$ cm$`{}_{}{}^{1}<\nu <40`$ cm<sup>-1</sup> and for temperatures 4 K $`<T<300`$ K. The in-plane scattering rate is shown to be unchanged as the sample becomes superconducting. The temperature dependence of the in-plane conductivity may be reasonably described within the Born limit of a dirty d-wave superconductor. The c-axis dielectric constant $`\epsilon _{1c}`$ is dominated by a phonon contribution at high temperatures. A zero crossing of $`\epsilon _{1c}`$ is directly observed below $`T_\mathrm{C}`$ which corresponds to the screened c-axis plasma frequency. In contrast to other cuprate superconductors, the temperature dependence of the c-axis conductivity closely follows the in-plane behavior.
This work was supported by BMBF (13N6917/0 - EKM) and in part by the Deutsche Forschungsgemeinschaft through SFB 484.
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# First Steps in Glass Theory
## 1 Introduction
When quenched fast enough so that it avoids the crystallisation transition, almost any liquid becomes a glass. This means that the density profile is not flat as in a liquid, it contains some peaks as in a crystal, but these peaks are not located on the nodes of a periodic or quasi periodic lattice. The understanding of such amorphous ’solid’ states has been recognised for a long time as a major question in condensed matter physics. The sentence by Phil Anderson: ”… there are still fascinating questions of principle about glasses and other amorphous phases…” , written nearly thirty years ago, was once again visionary in that it foresaw the wonderful developments on glassy systems, and particularly on spin glasses. The progress has been particularly difficult in these area, and in particular as far as structural glasses are concerned.
## 2 Mathematics
The first question which comes to mind is whether the glass is a new state of matter. It is not distinguished by any obvious symmetry (a not-obvious symmetry will be discussed later) from the liquid state, and one might think (as many people do) that the density profile would actually become flat on time scales longer than the experimental ones: the glass would just be a liquid with a long relaxation time.
From the statistical physics point of view, one wants to start from a microscopic Hamiltonian. The simplest situation is that of $`N`$ point-like particles in a volume $`V`$, with a pair interaction potential
$$H=\underset{i<j}{}V_{ij}(r_ir_j)$$
(1)
A simple case is that of homogeneous systems where $`V_{ij}`$ is independent of $`i`$ and $`j`$, and can be for instance either a hard sphere potential, a ‘soft sphere’ potential ($`V_{ij}(r)=A/r^{12}`$), or a Lennard-Jones potential ($`V_{ij}(r)=A/r^{12}B/r^6`$). Also much studied numerically, because the crystallisation is more easily avoided, are the binary mixtures where there are two types of particles: each particle $`i`$ has $`ϵ_i\{\pm 1\}`$ and $`V_{ij}(r)=V_{ϵ_iϵ_j}(r)`$, where $`V_{++}`$, $`V_{}`$, and $`V_+=V_+`$ are three potentials of the same type as before, but with different $`A,B`$ parameters corresponding to particles $`+`$ and $``$ having different radii.
Does there exist, in any such case, an independent state of matter which is the glass state? Does it exist as a long-lived metastable state (like the diamond phase of carbon)? Nobody knows the rigorous mathematical answer to these questions. Actually much simpler related questions are unanswered (e.g. proving the existence of a spin glass phase in a finite dimensional short range system ), or have taken many efforts to solve (e.g. proving Kepler’s conjecture that the densest three dimensional packing of hard spheres is the fcc/hcp lattice ).
## 3 Experiments
Experimentally, the liquid falls out of equilibrium on experimental time scales, and becomes a ‘glass’, at a temperature $`T_g`$ called the glass temperature. This glass temperature is conventionnally defined as the one at which the relaxation time $`\tau `$ of the liquid, as obtained e.g. from viscosity or from susceptibility measurements, becomes of the order of $`10^3`$ seconds. Angell’s plot of $`\mathrm{log}\left(\tau /1s\right)`$ versus $`T_g/T`$ allows to distinguish several types of behaviour (fig. 1). So called strong glasses like $`SiO_2`$ have a typical Arrhenius behaviour with one well defined free energy barrier. On the other hand, some glasses, called fragile, show a dramatic increase of the relaxation time when decreasing temperature which is much faster than Arrhenius: the typical free energy barrier thus increases when $`T`$ decreases. This implies a collective behaviour involving more and more particles. An increase of the dynamical correlation, characteristic of the mobile particles (rather than the more natural correlation of frozen particles), has been found in recent simulations . A popular fit of the relaxation time versus temperature is the Vogel Fulcher one,
$$\tau \tau _0\mathrm{exp}\left(\frac{A}{TT_{VF}}\right)$$
(2)
which would predict a phase transition at a temperature $`T_{VF}`$ which is not accessible experimentally (while staying at equilibrium). The more fragile the glass, the closer is $`T_{VF}`$ to $`T_g`$, while strong glass formers have a $`T_{VF}`$ close to zero.
Another interesting experimental signature is that of the specific heat. When one cools the liquid slowly, at a cooling rate $`\mathrm{\Gamma }=dT/dt`$, it freezes into a glass at a temperature which decreases slightly when $`\mathrm{\Gamma }`$ decreases. When this freezing occurs, the specific heat jumps downward, from its value in the equilibrated supercooled liquid state to a glass value which is close to that of the crystal. From the specific heat, one can compute the entropy. The configurational entropy, defined experimentally as the difference $`S_c=S_{liq}S_{crystal}`$, behaves smoothly in the supercooled liquid phase, until the system becomes a glass (see fig.1). It was noted by Kauzmann long ago that, if extrapolated, $`S_c(T)`$ vanishes at a finite temperature $`T_K`$. If cooled more slowly, the system follows the smooth $`S_c(T)`$ curve down to slightly lower temperatures, but then freezes again. One can wonder what could happen at infinitely slow cooling. As a negative $`S_c`$ does not make sense (except for pure hard spheres, where there is no energy), something must happen at temperatures above $`T_K`$. The curve $`S_c(T)`$ could flatten down smoothly, or there might be a phase transition, which in the simplest scenario would lead to $`S_c(T)=0`$ at $`T<T_K`$. This idea of an underlying ”ideal” phase transition, which could be obtained only at infinitely slow cooling, receives some support from the following observation: the two temperatures where the extrapolated experimental behaviour has a singularity, $`T_{VF}`$ and $`T_K`$, turn out to be amazingly close to each other (see the table below). The first phenomenological attempts to explain this fact originate in the work of Kauzmann , and developed among others by Adam, Gibbs and Di-Marzio , which identifies the glass transition as a ‘bona fide’ thermodynamic transition blurred by some dynamical effects.
If there exists a true thermodynamic glass transition at $`T=T_K=T_{VF}`$, it is a transition of a strange type. It is of second order because the entropy and internal energy are continuous. On the other hand the order parameter is discontinuous at the transition, as in first order transitions: the modulation of the microscopic density profile in the glass does not appear continuously from the flat profile of the liquid. As soon as the system freezes, there is a finite jump in this modulation (A more precide definition of the order parameter will be given below).
Comparison of $`T_K`$ and $`T_{VF}`$ in various glass-formers (from )
| Substance | $`T_K(K)`$ | $`T_{VF}(K)`$ | $`T_g(K)`$ |
| --- | --- | --- | --- |
| o-terphenyl | 204.2 | 202.4 | 246 |
| salol | 175.2 | | 220 |
| 2-MTHF | 69.3 | 69.6 | 91 |
| n-propanol | 72.2 | 70.2 | 97 |
| 3-bromopentane | 82.5 | 82.9 | 108 |
## 4 A mean field spin glass analogy
A totally different class of systems where such a 1st-2nd order type transition was found, and studied in great details, is a certain category of mean field spin glasses. A few years after the replica symmetry breaking (RSB) solution of the mean field theory of spin glasses , it was realized that there exists another category of mean-field spin glasses where the static phase transition exists and is due to an entropy crisis . These are now called discontinuous spin glasses because their phase transition, although it is of second order in the Ehrenfest sense, has a discontinuous order parameter . Another name often found in the literature is ‘one step RSB’ spin glasses, because of the special pattern of symmetry breaking involved in their solution. These are spin glasses with infinite range interactions involving a coupling between triplets (or higher order groups) of spins. The simplest among them is the random energy model, which is the $`p\mathrm{}`$ limit version of the p-spin models described by the Hamiltonian
$$H=\underset{i1<\mathrm{}<i_p}{}J_{i_1\mathrm{}i_p}s_{i_1}\mathrm{}s_{i_p}$$
(3)
where the $`J`$’s are (appropriately scaled) quenched random couplings, and the spins can be either of Ising or spherical type .
The analogy between the phase transition of discontinuous spin glasses and the thermodynamic glass transition was first noticed by Kirkpatrick, Thirumalai and Wolynes in a series of inspired papers of the mid-eighties . While some of the basic ideas of the present development were around at that time, there still missed a few crucial ingredients. On one hand one needed to get more confidence that this analogy was not just fortuitous. The big obstacle was the existence (in spin glasses) versus the absence (in structural glasses) of quenched disorder. The discovery of discontinuous spin glasses without any quenched disorder provided an important new piece of information: contrarily to what had been believed for long, quenched disorder is not necessary for the existence of a spin glass phase (but frustration is).
It is important to analyse critically this analogy from the point of view of the dynamical behaviour. In discontinuous mean field spin glasses there exist a dynamical transition temperature at a temperature $`T_c`$ which is larger than the equilibrium transition $`T_K`$. When T decreases and gets near to $`T_c`$, the correlation function relaxes with a characteristic two step forms: a fast $`\beta `$ relaxation leading to a plateau takes place on a characteristic time which does not grow, while the $`\alpha `$ relaxation from the plateau takes place on a time scale which diverges when $`TT_c`$ (see fig. 2). This dynamic transition is exactly described by the schematic mode coupling equations.
However the existence of a dynamic relaxation at a temperature above the true thermodynamic one is possible only in mean field, and the conjecture is that in a realistic system like a glass, the region between $`T_K`$ and $`T_c`$ will have instead a finite, but very rapidly increasing, relaxation time, as explained in fig. 2. A similar behaviour has been found in finite -size mean field models
Another very interesting dynamical regime is the one where the system is out of equilibrium ($`T<T_g`$). Then the system is no longer stationnary: it ages. This is well known for instance from studies in polymeric glasses. If one measures the response of your favorite plastic ruler to some stress, it will behave differently depending on its age. Schematically, new relaxation processes come into play on a time scale comparable to the age of the system: the older the system, the longer the time needed for this ”aging” relaxation to take place. Recent years have seen important developments on the out of equilibrium dynamics of the glassy phases , initiated by the exact solution of the dynamics in a discontinuous spin glass by Cugliandolo and Kurchan . It has become clear that, in realistic systems with short range interactions, the pattern of replica symmetry breaking can be deduced from the measurements of the violation of the fluctuation dissipation theorem . These measurements are difficult. However, numerical simulations performed on different types of glass forming systems have provided an independent and spectacular confirmation of their ‘one step rsb’ structure on the (short) time scales which are accessible. Experimental results have not yet settled the issue, but the first measurements of effective temperatures in the fluctuation dissipation relation have been made recently .
To summarize, the analogy between the phenomenology of fragile glass formers and discontinuous mean field spin glasses accounts for:
* The discontinuity of the order parameter
* The continuity of the energy and the entropy
* The jump in specific heat (and the sign of the jump)
* Kauzmann’s ”entropy crisis”
* The two steps relaxation of the dynamics and the succes of Mode Coupling Theory at relatively high temperatures.
* The aging phenomenon and the pattern of modification of the fluctuation dissipation relation in the low temperature phase
## 5 A lesson from mean field: many valleys
The successes of the above analogy suggest to have a closer look at the mean field models in order to understand, at least at the mean field level, what are the basic ingredients at work in the glass transition. In mean field spin glasses, at temperatures $`T_K<T<T_c`$, the phase space breaks up into ergodic components which are well separated, so-called free energy valleys or TAP states . Each valley $`\alpha `$ has a free energy $`F_\alpha `$ and a free energy density $`f_\alpha =F_\alpha /N`$. The number of free energy minima with free energy density $`f`$ is found to be exponentially large:
$$𝒩(f,T,N)\mathrm{exp}(N\mathrm{\Sigma }(f,T)),$$
(4)
where the function $`\mathrm{\Sigma }`$ is called the complexity. The total free energy of the system, $`\mathrm{\Phi }`$, can be well approximated by:
$$e^{\beta N\mathrm{\Phi }}\underset{\alpha }{}e^{\beta Nf_\alpha (T)}=_{f_{min}}^{f_{max}}𝑑f\mathrm{exp}\left(N[\mathrm{\Sigma }(f,T)\beta f]\right),$$
(5)
where $`\beta =1/T`$. The minima which dominate the sum are those with a free energy density $`f^{}`$ which minimizes the quantity $`\mathrm{\Phi }(f)=fT\mathrm{\Sigma }(f,T)`$. At large enough temperatures the saddle point is at $`f>f_{min}(T)`$. When one decreases $`T`$ the saddle point free energy decreases (see fig.3, with $`m=1`$). The Kauzman temperature $`T_K`$ is that below which the saddle point sticks to the minimum: $`f^{}=f_{min}(T)`$. It is a genuine phase transition . However because $`T_c>T_K`$, the phase space is actually separated into non ergodic components (valleys) at $`T<T_c`$ (actually there exist some non ergodic components also above $`T_c`$, but they are not felt by the system when starting from random initial conditions ). The total equilibrium free energy is analytic at $`T_c`$: in spite of the ergodicity breaking, the system has the same free energy as that of the liquid, as if transitions were allowed between valleys.
What remains of this mean field picture in finite dimensional glasses? When one decreases the temperature, there is a well defined separation of time scales between the $`\alpha `$ and the $`\beta `$ relaxations, which suggest to consider the dynamical evolution of system in phase space as a superposition of two processes: an intravalley (relatively fast) relaxation, and an intervalley (slow, ang getting rapidly much slower when one cools the system) hopping process.
One popular way of making this statement more precise, allowing to study it numerically, is through the use of inherent structures (IS) . Given a configuration of the system, characterized by its phase space position $`x=\{\stackrel{}{x}_1,\mathrm{},\stackrel{}{x}_N\}`$, the corresponding inherent structure $`s(x)`$ is another point of phase space which is the local minimum of the Hamiltonian which is reached from this configuration through a steepest descent dynamics. IS are easily identified numerically. A given trajectory $`x(t)`$ of the glass through phase space maps onto the corresponding trajectory $`s(t)`$ in the space of inherent structures. Looking at the dynamical evolution in the space of IS makes the valley structure slightly more apparent, since one gets rid of the small thermal excitations around each valley minimum. Calling $`𝒟_s`$ the set of those configurations which are mapped to the coherent structure $`s`$, a natural definition of the IS entropy density, $`\mathrm{\Sigma }_{is}`$, is $`N\mathrm{\Sigma }_{is}(T)=_sP(s)\mathrm{ln}(P(s))`$, where the weight of the inherent structure $`s`$ is
$$P(s)=Z(s)/\underset{b}{}Z(b);Z(s)_{x𝒟_s}dx\mathrm{exp}(\beta H(x)).$$
(6)
In a system with short range interactions, it is reasonable to expect that one may have two distinct IS which differ by a local rearrangement of a finite number of atoms. It is then easy to show that the slope of configurational entropy versus free energy is infinite around $`f_{min}`$ , which does not agree with the general scenario, except if the Kauzmann temperature vanishes. This problem is due to te fact that IS are too simple objects, which cannot be identified with the free energy valleys. The difference is very easily seen in spin systems : IS are nothing but configurations which are stable to one spin flip. Zero temperature free energy valleys, defined as TAP states, are stable to the flip of any arbitrarily large number $`k`$ of spins (but the limit $`N\mathrm{}`$ must be taken before the limit $`k\mathrm{}`$). In continuous systems, the generalization is clear: IS are local minima of the energy, so that any infinitesimal move of the positions of all $`N`$ particles raises the energy. Let us generalize the notion of a minimum as follows: define a k-th order local minimum as a configuration of particles such that any infinitesimal move of all $`N`$ particles, together with a move of arbitrary size of $`k`$ particles, raises the energy. The limit $`k\mathrm{}`$ gives the proper definition of a zero temperature free energy valley. The proper definition at finite temperature is slightly more involved . Let us summarize it here briefly. Given two configurations $`x`$ and $`y`$ we define their overlap as before as $`q(x,y)=1/N_{i,k=1,N}w(x_iy_k),`$ where $`w(x)=1\text{for}x`$ small, $`w(x)=0\text{for}x`$ larger than the typical interatomic distance. We add an extra term to the Hamiltonian: we define
$`\mathrm{exp}(N\beta F(y,ϵ))={\displaystyle 𝑑x\mathrm{exp}(H(x)+\beta ϵNq(x,y))},`$
$`F(ϵ)=F(y,ϵ),`$ (7)
where $`f(y)`$ denotes the average value of $`f`$ over equilibrium configurations $`y`$ thermalized at temperature $`\beta ^1`$. Taking the thermodynamic limit before the limit $`ϵ0`$ allows to identify the valley around any generic equilibrium configuration $`y`$ .
In a nutshell, two configurations which differ by the (arbitrarily large) displacement of a finite number of atoms are in the same thermodynamic valley. This definition of the valleys also suffers from some difficulties: Nucleation arguments then forbid the existence of a non-trivial complexity versus free energy curve in a finite dimensional system. The solution consists in noticing that there exist many more metastable valleys, which have a finite but very long lifetime. These can be identified by taking the Legendre transform $`W(q)`$ of the free energy $`F(ϵ)`$:
$$W(q)=F(ϵ)+ϵq;q=\frac{F}{ϵ}.$$
(8)
Analytic computation in mean field models , as well as in glass forming liquids using the replicated HNC approximation , show that $`W(q)`$ is minimal at $`q=0`$, but has a secondary minimum at a certain $`q=q_{EA}`$, in the temperature range $`T_K<T<T_c`$. The behaviour around this second, metastable, minimum corresponds to phenomena that can be observed on time scales shorter than the lifetime of the metastable state. The thermodynamic configurational entropy is the value of the potential $`W(q)`$ at the secondary minimum with $`q0`$ , and it can be defined only if the minimum does exist (i.e. for $`T<T_c`$). Of course the secondary minimum for $`T>T_k`$ is always in the metastable region. However if one would start from a large value of $`ϵ`$ and would decrease $`ϵ`$ to zero not too slowly, the system would not escape from the metastable region and one obtains a proper definition of the thermodynamic configurational entropy in this region $`T>T_K`$. In a similar way one could compute $`q(ϵ)`$ in the region ($`ϵ>ϵ_c`$) where the high $`q`$ phase is thermodynamically stable and extrapolate it to $`ϵ0`$. The ambiguity in the definition of the thermodynamic configurational entropy at temperatures above $`T_k`$ becomes larger and larger when the temperature increases. It cannot be defined for $`T>T_c`$.
## 6 Beyond the analogy: first principles computation
In recent years, it has become possible to go beyond the simple analogy between structural glasses and mean field discontinuous spin glasses. One can actually use the concepts and the techniques which are suggested by this analogy in order to start a systematic first principles study of the glass phase . So far we have focused onto the equilibrium study of the low temperature phase. One main reason is that the direct study of out of equilibrium dynamics is more difficult, and that one might be able to make progress by a careful analysis of the landscape . The strategy is to assume that there exists a phase transition, and that it is of the same type as the one found in discontinuous mean field spin glasses. Within this framework, one tries to compute the properties of the glass phase. This involves several quantitites like the Kauzmann temperature, the radius of the cage which confine the particles in the glass phase, the configurational entropy etc… The validity of the scenario is checked from the comparison of various predictions with numerical simulations of well equilibrated systems.
The first task is to define an order parameter. This is not trivial in an equilibrium theory where we have no notion of time persistent correlations. The best way is to introduce two copies of the system, with a weak interaction. The two sets of particles have positions $`x_i`$ and $`y_i`$ respectively, the total Hamiltonian is
$$E=\underset{1ijN}{}(v(x_ix_j)+v(y_iy_j))+ϵ\underset{i,j}{}w(x_iy_j)$$
(9)
where we have introduced a small attractive potential $`w(r)`$ between the two systems. The precise shape of $`w`$ is irrelevant, insofar as we shall be interested in the limit $`ϵ0`$, but its range should be of order or smaller than the typical interparticle distance. The order parameter is then the correlation function between the two systems:
$$g_{xy}(r)=\underset{ϵ0}{lim}\underset{N\mathrm{}}{lim}\frac{1}{\rho N}\underset{ij}{}<\delta (x_iy_jr)>$$
(10)
In the liquid phase this correlation function is identically equal to one, while it has a nontrivial structure in the glass phase, reminiscent of the pair correlation of a dense liquid, but with an extra peak around $`r0`$. Let us notice that we expect a discontinuous jump of this order parameter at the transition, in spite of its being second order in the thermodynamic sense. The existence of a non trivial order parameter is associated with the spontaneous breaking of a symmetry: For $`ϵ=0`$, with periodic boundary conditions, the system is symmetric under a global translation of the $`x`$ particles with respect to the $`y`$ particles. This symmetry is spontaneously broken in the low temperature phase, where the particles of each subsystem tend to sit in front of each other.
Generalizing this approach to a system of $`m`$ coupled replicas, sometimes named ‘clones’ in this context (the order parameter used only $`m=2`$), provides a wonderful method for studying analytically the thermodynamics of the glass phase . In the glass phase, the attraction will force all $`m`$ systems to fall into the same glass state, so that the partition function is:
$$Z_m=\underset{\alpha }{}e^{\beta Nmf_\alpha (T)}=_{f_{min}}^{f_{max}}𝑑f\mathrm{exp}\left(N[\mathrm{\Sigma }(f,T)m\beta f]\right)$$
(11)
In the limit where $`m1`$ the corresponding partition function $`Z_m`$ is dominated by the correct saddle point $`f^{}`$ for $`T>T_K`$. The interesting regime is when the temperature is $`T<T_K`$, and the number $`m`$ is allowed to become smaller than one. The saddle point $`f^{}(m,T)`$ in the expression (11) is the solution of $`\mathrm{\Sigma }(f,T)/f=m/T`$. Because of the convexity of $`\mathrm{\Sigma }`$ as function of $`f`$, the saddle point is at $`f>f_{min}(T)`$ when $`m`$ is small enough, and it sticks at $`f^{}=f_{min}(T)`$ when $`m`$ becomes larger than a certain value $`m=m^{}(T)`$, a value which is smaller than one when $`T<T_K`$ (see fig. 3). The free energy in the glass phase, $`F(m=1,T)`$, is equal to $`F(m^{}(T),T)`$. As the free energy is continuous along the transition line $`m=m^{}(T)`$, one can compute $`F(m^{}(T),T)`$ from the region $`mm^{}(T)`$, which is a region where the replicated system is in the liquid phase. This is the clue to the explicit computation of the free energy in the glass phase. It may sound a bit strange because one is tempted to think of $`m`$ as an integer number. However the computation is much clearer if one sees $`m`$ as a real parameter in (11). As one considers low temperatures $`T<T_K`$ the $`m`$ coupled replicas fall into the same glass state and thus they build some molecules of $`m`$ atoms, each molecule being built from one atom of each ’colour’. Now the interaction strength of one such molecule with another one is basically rescaled by a factor $`m`$ (this statement becomes exact in the limit of zero temperature where the molecules become point like). If $`m`$ is small enough this interaction is small and the system of molecules is liquid. When $`m`$ increases, the molecular fluid freezes into a glass state at the value $`m=m^{}(T)`$. So our method requires to estimate the replicated free energy, $`F(m,T)=\mathrm{log}(Z_m)/(\beta mN)`$, in a molecular liquid phase, where the molecules consist of $`m`$ atoms and $`m`$ is smaller than one. For $`T<T_K`$, $`F(m,T)`$ is maximum at the value of $`m=m^{}`$ smaller than one, while for $`T>T_K`$ the maximum is reached at a value $`m^{}`$ is larger than one. The knowledge of $`F_m`$ as a function of $`m`$ allows to reconstruct the configurational entropy function $`Sc(f)`$ at a given temperature $`T`$ through a Legendre transform, using the parametric representation (easily deduced from a saddle point evaluation of (11)) :
$$f=\frac{\left[mF(m,T)\right]}{m};\mathrm{\Sigma }(f)=\frac{m^2}{T}\frac{F(m,T)}{m}.$$
(12)
The Kauzmann temperature (’ideal glass temperature’) is the one such that $`m^{}(T_K)=1`$. For $`T<T_K`$ the equilibrium configurational entropy vanishes. Above $`T_K`$ one obtains the equilibrium configurational entropy $`\mathrm{\Sigma }(T)`$ by solving (12) at $`m=1`$.
This gives the main idea which allows to compute the free energy in the glass phase, at a temperature $`T<T_K`$, from first principles: it is equal to the free energy of a molecular liquid at the same temperature, where each molecule is built of $`m`$ atoms, and an appropriate analytic continuation to $`m=m^{}(T)<1`$ has been taken. The whole problem is reduced to a computation in a liquid. This is not trivial, and requires to develop some specific approximations. I shall not elaborate on that here, but refer the reader to the original papers . The basic idea of the approximation is that the size of the molecules is directly related to the thermal wandering of an atom in its cage. Therefore at low temperatures one can use some small cage approximation. it is natural to write the partition function in terms of the center of mass and relative coordinates $`\{r_i,u_i^a\}`$, with $`x_i^a=r_i+u_i^a`$ and $`_au_i^a=0`$, and to expand the interaction in powers of the relative displacements $`u`$. Keeping only the term quadratic in $`u`$ (harmonic vibrations of the molecules), and integrating over these vibration modes, one gets the ”harmonic resummation” approximation where the partition function is given by:
$$Z_m=Z_m^0𝑑r\mathrm{exp}\left(\beta mH(r)\frac{m1}{2}Tr\mathrm{log}M\right)$$
(13)
where $`Z_m^0=m^{Nd/2}\sqrt{2\pi T}^{Nd(m1)}/N!`$, and the matrix $`M`$, of dimension $`dN\times dN`$, is given by:
$$M_{(i\mu )(j\nu )}=\frac{^2H(r)}{r_i^\mu r_j^\nu }=\delta _{ij}\underset{k}{}v_{\mu \nu }(r_ir_k)v_{\mu \nu }(r_ir_j)$$
(14)
and $`v_{\mu \nu }(r)=^2v/r_\mu r_\nu `$ (the indices $`\mu `$ and $`\nu `$ denote space directions). Now we are back to a real problem of liquid theory, since we have only $`d`$ degrees of freedom per molecule (the center of mass coordinates), and the number of clones, $`m`$, appears as a parameter in (13).
Once one has derived an expression for the replicated free energy, one can deduce from it the whole thermodynamics, as described above (Notice that the ‘technical’ approximation of neglecting the exchange of atoms between different molecules, as well as using a harmonic model, means that one really studies the IS in this computation, rather than the real free energy valleys). In all three cases, one finds an estimate of the Kauzmann temperature which is in reasonable agreement with simulations, with a jump in specific heat, from a liquid value at $`T>T_K`$ to the Dulong-Petit value $`C=3/2`$ (we have included only positional degrees of freedom) below $`T_K`$. This is similar to the experimental result, where the glass specific heat jumps down to the crystal value when one decreases the temperature (Our approximations so far are similar to the Einstein approximation of independent vibrations of atoms, in which case the contribution of positional degrees of freedom to the crystal specific heat is $`C=3/2`$). The parameter $`m^{}(T)`$ and the cages sizes are nearly linear with temperature in the whole glass phase. This means, in particular, that the effective temperature $`T/m`$ is always close to $`T_K`$, so in our theoretical computation we need only to evaluate the expectation values of observables in the liquid phase, at temperatures where the HNC approximation for the liquid still works quite well.
A more detailed numerical check of these analytical predictions involves the measurement of the complexity,
$$\mathrm{\Sigma }_t=S(T)S_{valley}(T)$$
(15)
The liquid entropy is estimated by a thermodynamic integration of the specific heat from the very dilute (ideal gas) limit. It turns out that in the deeply supercooled region the temperature dependence of the liquid entropy is well fitted by the law predicted in : $`S_{liq}(T)=aT^{2/5}+b`$, which presumably allows for a good extrapolation at temperatures $`T`$ which cannot be simulated. As for the ’valley’ entropy, it can be estimated as that of an harmonic solid. One needs however the vibration frequencies of the solid. These have been approximated by several methods, most of which are based on some evaluation of the Instantaneous Normal Modes (INM) in the liquid phase, and the assumption that the spectrum of frequencies does not depend much on temperature below $`T_K`$. Starting from a typical configuration of the liquid, one can look at the INM around it. In general there exist some negative eigenvalues (the liquid is not a local minimum of the energy) which one must take care of. Several methods have been tried: either one keeps only the positive eigenvalues, or one considers the absolute values of the eigenvalues . Alternatively one can also consider the INM around the nearest inherent structure which has by definition a positive spectrum . This procedure really measures the configurational entropy rather than the thermodynamic complexity. The computation of the thermodynamic complexity, using its definition as a system coupled to a reference thermalized configuration, has also been computed in and turns out to be not very different from the configurational entropy, on the time and temperature scales which have been studied so far (they must differ on infinitely long time scales, as we discussed in the previous section).
The results for the configurational entropy as a function of temperature are shown in fig.4, for binary mixtures of soft spheres and of Lennard-Jones particles. The agreement with the analytical result obtained from the replicated fluid system is rather satisfactory, considering the various approximations involved both in the analytical estimate and in the numerical ones.
## 7 Conclusion
Our knowledge on first principle computations of glasses is still rather primitive. Basically we have obtained, for the glass, the equivalent of the Einstein approximation for the crystal. Even within this simple scheme, doing the actual computation for the glass turns out to be rather involved. What is most needed next is: on the analytical side, some better approximations of the molecular liquid state, allowing to go beyond the small cage expansions, and some reliable estimation of time scales in the regime $`T_K<T<T_c`$; on the numerical side, some precise results in the glass phase at equilibrium ; on the experimental side, some more measurements of the fluctuation dissipation ratio in the out of equilibrium dynamics. No doubt: ”… there are still fascinating questions of principle about glasses and other amorphous phases…” .
## 8 Acknowledgments
It is a great pleasure to thank Giorgio Parisi for the collaboration which led to the works described here, as well as G.Biroli and R. Monasson for useful discussions.
## 9 Bibliography
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# 1 Introduction
## 1 Introduction
High luminosities of $`B`$-factories and hadron colliders open a real experimental possibility to observe doubly heavy baryons $`\mathrm{\Xi }_{QQ^{}}`$ and $`\mathrm{\Omega }_{QQ^{}}`$. This prospective stimulates the theoretical interest to study the physics of such the baryons: the spectroscopy in the framework of both potential models and QCD sum rules , the lifetimes and inclusive decay modes , the mechanism of production in various collisions and the rate of yield at accelerators .
The pair production can be essential at the energies close to the threshold, that we consider in the present paper. The physical approximation for the calculations are caused by the apparent diquark structure of doubly heavy baryons, since the diquark size is essentially less than the radius of confinement determining the motion of light quark inside the $`QQ^{}q`$ system. We suppose that, at first, the calculations of heavy diquark production should be done, and further, the models of diquark fragmentation into the baryon can be applied.
In paper the differential and total cross sections for exclusive production of meson pairs in $`e^+e^{}`$ annihilation was calculated in the franework of constituent quark model. The results were obtained for the case of final particles in pseudoscalar and vector states close to the threshold. Following the same procedure, we examinate the processes of $`e^+e^{}𝔡\overline{𝔡}`$ and $`q\overline{q}𝔡\overline{𝔡}`$, where $`𝔡`$, $`\overline{𝔡}`$ are diquark and antidiquark, respectively (in our calculations we neglect masses of annihilating particles). Differential and total cross sections for the exclusive production of heavy diquark pairs in axial-axial, axial-scalar and scalar-scalar states are determined for the diquark composed of different heavy quarks. We also consider the case of axial diquark composed of equivalent quarks. We do not concern for the annihilation into the pseudoscalar and vector diquarks, since their production does not contribute in the leading order of $`1/m`$ expansion. If one supposes the fragmentation of diquarks into the doubly heavy baryons, the formulae derived may be useful in the calculation of cross sections for the pair production of baryons.
In section 2, we describe basic points of the constituent quark model. Section 3 and 4 contains matrix elements, differential and total cross sections for the processes of $`e^+e^{}`$ and $`q\overline{q}`$ annihilation into the pairs of diquarks. In section 5 numerical results are given. In Conclusion we summarize the consideration.
## 2 Basic points of the model
In this paper we use the constituent quark model , which considers the quark masses and leptonic constants as the only input parameters.
According to the model the diquark $`𝔡=(Q_1Q_2)`$ mass is equal to
$$M=m_1+m_2.$$
Four-momenta of the quarks entering the diquark are given by
$$k_{Q_1}=\frac{m_1}{M}P+q,$$
$$k_{Q_2}=\frac{m_2}{M}Pq,$$
where $`P`$ is the meson momentum.
To represent The Fock state of diquark with two different quarks in the model, we use the principle of superposition for the wave packages with distributions $`\mathrm{\Psi }`$. Therefore, we find the following:
in the case of scalar diquark
$$|S_𝔡^i=\frac{ϵ_{ijk}}{\sqrt{2}}\frac{d^3q}{(2\pi )^3}\mathrm{\Psi }_s(q)\underset{\lambda _1\lambda _2}{}\frac{(\mathrm{\Psi }_{\lambda _1}^{}\widehat{C}\gamma _5\mathrm{\Psi }_{\lambda _2})^{}}{\sqrt{2}}\widehat{a}_{\lambda _1}^{j+}\widehat{b}_{\lambda _2}^{k+}|0,$$
(1)
and for the axial diquark
$$|A_𝔡^i=\frac{ϵ_{ijk}}{\sqrt{2}}\frac{d^3q}{(2\pi )^3}\mathrm{\Psi }_a(q)\underset{\lambda _1\lambda _2}{}\frac{(\mathrm{\Psi }_{\lambda _1}^{}\widehat{C}\gamma _\mu \mathrm{\Psi }_{\lambda _2})^{}}{\sqrt{2}}e_\mu \widehat{a}_{\lambda _1}^{j+}\widehat{b}_{\lambda _2}^{k+}|0,$$
(2)
where $`\widehat{C}`$ is a matrix of charge conjugation, $`i`$ and $`j`$ are the colour indices, $`\widehat{e}=e_\mu \gamma _\mu `$, $`e_\mu `$ is a polarisation vector of axial diquark, $`\widehat{a}`$ and $`\widehat{b}`$ denote the operators of quark creation, so that the axial and scalar diquarks in the Fock space are normilized in the following way:
$`S_𝔡^i(P)|S_𝔡^j(P^{})`$ $`=`$ $`(2\pi )^3\delta _{ij}\delta (\stackrel{}{P}\stackrel{}{P}^{}),`$ (3)
$`A_𝔡^i(P,\lambda )|A_𝔡^j(P^{},\lambda ^{})`$ $`=`$ $`(2\pi )^3\delta _{ij}\delta _{\lambda \lambda ^{}}\delta (\stackrel{}{P}\stackrel{}{P}^{}),`$ (4)
$`\lambda `$ and $`\lambda ^{}`$ are diquark polarisation indices.
If the heavy quarks are identical, then the Pauli principle must be taken into account. So, the above formulae have to be divided by $`\sqrt{2}`$ and antisymmetrized over the permutaions of $`\widehat{a}`$ and $`\widehat{b}`$ operators.
The heavy quark propagator has the form
$$S(k)=(k_\mu \gamma _\mu +m)D(k),$$
where
$$D^1(k)=k^2m^2.$$
## 3 Amplitudes and cross sections for the $`e^+e^{}`$ annihilation
In this section we consider the $`e^+e^{}`$ annihilation into the diquark and antidiquark pairs. The diagrams, which contribute into the processes in the leading order, are shown in Figs. 1a and 1b. The colour factor for the processes involved is equal to
$$\mathrm{Colour}_{ij}=\frac{2}{3}\delta _{ij},$$
where $`i`$, $`j`$ are the colours indices of diquark and antidiquark, correspondingly.
### 3.1 Annihilation into pairs of scalar diquarks
The matrix element for the pair production of scalar diquarks may be writen as follows:
$$_{ss}=i\frac{64\pi ^2}{3}\frac{f_{ss}}{s^2}\delta _{ij}|\mathrm{\Psi }_s(0)|^2(P_\mu ^{}P_\mu )l_\mu $$
(5)
where $`f_{ss}`$ is equal to
$`f_{ss}`$ $`=`$ $`M\left(\alpha _s\left({\displaystyle \frac{m_1^2}{M^2}}s\right){\displaystyle \frac{q_2}{m_1^2}}+\alpha _s\left({\displaystyle \frac{m_2^2}{M^2}}s\right){\displaystyle \frac{q_1}{m_2^2}}\right)\alpha _{em}\left(s\right)`$ (6)
$`{\displaystyle \frac{2M^3}{s}}\left(\alpha _s\left({\displaystyle \frac{m_1^2}{M^2}}s\right){\displaystyle \frac{q_2m_2}{m_1^3}}+\alpha _s\left({\displaystyle \frac{m_2^2}{M^2}}s\right){\displaystyle \frac{q_1m_1}{m_2^3}}\right)\alpha _{em}\left(s\right),`$
and $`l_\mu `$ denotes the leptonic vector current, $`\mathrm{\Psi }_s(0)`$ is a wave function of scalar diquark at the origin. The charges of the quarks $`Q_1`$ and $`Q_2`$ are equal to $`q_1`$ and $`q_2`$. $`P^{}`$, $`P`$ are the four-momenta of scalar diquark and antidiquark, correspondingly.
After simple algebraic calculations for the differential cross section $`d\sigma _{ss}`$/ $`d\mathrm{cos}\theta `$ we get the following expression:
$$\frac{d\sigma _{ss}}{d\mathrm{cos}\theta }=64\pi ^3\frac{f_{ss}^2}{3s^3}|\mathrm{\Psi }_s(0)|^4\left(1\frac{4M^2}{s}\right)^{3/2}(1\mathrm{cos}^2\theta ),$$
(7)
where $`\theta `$ is the angle between the momenta of lepton and diquark.
The total cross section for the exclusive production of heavy scalar diquark pairs in $`e^+e^{}`$ annihilation is equal to
$$\sigma _{ss}=256\pi ^3\frac{f_{ss}^2}{9s^3}|\mathrm{\Psi }_s(0)|^4\left(1\frac{4M^2}{s}\right)^{3/2}.$$
(8)
### 3.2 Annihilation into axial and scalar diquarks
The matrix element for the pair production of axial antidiquark and scalar diquark may be writen as follows:
$$_{as}=\frac{128\pi ^2}{3s^3}\delta _{ij}f_{as}\mathrm{\Psi }_s^{}(0)\mathrm{\Psi }_a(0)ϵ_{\mu \alpha \beta \gamma }e_\alpha P_\beta q_\gamma l_\mu ,$$
(9)
where $`f_{as}`$ is equal to
$$f_{as}=M^3\left(\alpha _s\left(\frac{m_1^2}{M^2}s\right)\frac{q_2}{m_1^3}\alpha _s\left(\frac{m_2^2}{M^2}s\right)\frac{q_1}{m_2^3}\right)\alpha _{em}\left(s\right).$$
$`\mathrm{\Psi }_a(0)`$ is a wave function of axial diquark at the origin, $`q=P+P^{}`$.
Carring out algebraic calculations for the differential cross section $`d\sigma _{as}/d\mathrm{cos}\theta `$ we find the following expression:
$$\frac{d\sigma _{as}}{d\mathrm{cos}\theta }=64\pi ^3\frac{f_{as}^2}{3s^4}|\mathrm{\Psi }_s(0)\mathrm{\Psi }_a(0)|^2\left(1\frac{4M^2}{s}\right)^{3/2}(2\mathrm{sin}^2\theta ).$$
(10)
The total cross section for the exclusive production of heavy scalar diquark and axial antidiquark pairs in $`e^+e^{}`$ annihilation is given by the following formula:
$$\sigma _{as}=512\pi ^3\frac{f_{as}^2}{9s^4}|\mathrm{\Psi }_s(0)\mathrm{\Psi }_a(0)|^2\left(1\frac{4M^2}{s}\right)^{3/2}.$$
(11)
In the above formulae the difference between the masses of axial and scalar diquarks is neglected.
### 3.3 Annihilation into pairs of axial diquarks
The matrix element for the exclusive pair production of two axial diquarks may be represented as follows:
$$_{aa}=i\frac{128\pi ^2}{3s^3}\delta _{ij}|\mathrm{\Psi }_a(0)|^2\left(f_{aa}^{[1]}(P_\mu ^{}P_\mu )(e^{}e)+f_{aa}^{[2]}((e^{}q)e_\mu (eq)e_\mu ^{})\right)l_\mu ,$$
(12)
where $`f_{aa}^{[1]}`$ and $`f_{aa}^{[2]}`$ are equal to
$`f_{aa}^{[1]}`$ $`=`$ $`M^3\left(\alpha _s\left({\displaystyle \frac{m_1^2}{M^2}}s\right){\displaystyle \frac{q_2m_2}{m_1^3}}+\alpha _s\left({\displaystyle \frac{m_2^2}{M^2}}s\right){\displaystyle \frac{q_1m_1}{m_2^3}}\right)\alpha _{em}\left(s\right),`$ (13)
$`f_{aa}^{[2]}`$ $`=`$ $`M^4\left(\alpha _s\left({\displaystyle \frac{m_1^2}{M^2}}s\right){\displaystyle \frac{q_2}{m_1^3}}+\alpha _s\left({\displaystyle \frac{m_1^2}{M^2}}s\right){\displaystyle \frac{q_1}{m_2^3}}\right)\alpha _{em}\left(s\right).`$ (14)
For the differential cross section $`d\sigma _{aa}/d\mathrm{cos}\theta `$ we calculate the following expression:
$$\frac{d\sigma _{aa}}{d\mathrm{cos}\theta }=\frac{512\pi ^3}{3s^5}|\mathrm{\Psi }_a(0)|^4\left(1\frac{4M^2}{s}\right)^{3/2}(𝒜\mathrm{cos}^2\theta ),$$
where $`𝒜`$ and $``$ have the form
$$𝒜=(f_{aa}^{[1]})^2(8+(\eta 2)^2)2f_{aa}^{[1]}f_{aa}^{[2]}\eta (\eta 2)+(f_{aa}^{[2]})^2(\eta ^2+2\eta ),$$
$$=(f_{aa}^{[1]})^2(8+(\eta 2)^2)2f_{aa}^{[1]}f_{aa}^{[2]}\eta (\eta 2)+(f_{aa}^{[2]})^2(\eta ^22\eta ),$$
and $`\eta =s/M^2`$ .
The total cross section for the exclusive production of heavy axial diquarks pairs in $`e^+e^{}`$ annihilation is given by
$$\sigma _{aa}=\frac{1024\pi ^3}{9s^5}|\mathrm{\Psi }_a(0)|^4\left(1\frac{4M^2}{s}\right)^{3/2}(3𝒜).$$
(15)
Actually, the behaviour of this cross section at large $`s`$ is $`\sigma _{aa}1/s^3`$.
### 3.4 The diquark containing two equivalent particles
If the diquarks are composed of identical particles, the above formulae must be changed. Evidently, we have to take into account the production of axial diquark, only. Since the smallest angular momentum for the scalar diquark is equal to unit, it does not contribute in the leading order of $`1/m`$ expansion. In contrast, the angular momentum of axial diquark can be equal to zero.
All formulae writen down for the annihilation into two axial diquarks remains valid except two form factors $`f_{aa}^{[1]},f_{aa}^{[2]}`$. They must be transformed in the following way:
$`f_{aa}^{[1]}`$ $`=`$ $`2q\alpha _s(M^2)M\alpha _{em}\left(s\right),`$ (16)
$`f_{aa}^{[2]}`$ $`=`$ $`4q\alpha _s(M^2)M\alpha _{em}\left(s\right),`$ (17)
where q is the quark charge.
## 4 Amplitudes and cross sections for the $`q\overline{q}`$ annihilation
In this section we consider the $`q\overline{q}`$ annihilation into the diquark and antidiquark pairs. The diagrams, which contribute in the leading order, are shown in Figs. 1a, 1b, 2, 3a and 3b. The diagram shown in Fig. 2 results in zero in our approach, when the quarks are on the mass shells.
The colour factor for the diagrams shown in Figs. 1a, 1b is equal to
$$\mathrm{Colour}_{(ij)(lm)}^{[1]}=\frac{1}{3}t_{ij}^at_{lm}^a,$$
where $`m,l`$ are the colour indices of annihilating quark and antiquark, correspondingly. It is easy to notice that this is a colour octet state. The colour factor for the diagrams shown in Figs. 3a, 3b is equal to
$$\mathrm{Colour}_{(ij)(lm)}^{[2]}=\frac{5}{12}t_{ij}^at_{lm}^a\frac{1}{9}\delta _{ij}\delta _{lm},$$
that appears to be a mixture of octet and singlet colour states.
### 4.1 Annihilation into pairs of scalar diquarks
The matrix element for the exclusive pair production of two scalar diquarks may be written as follows:
$$_{ss}=\frac{32\pi ^2i}{s^2}|\mathrm{\Psi }_s(0)|^2\left(\stackrel{~}{f}_{ss}^{[1]}\frac{2t_{if}^at_{lm}^a}{3}P_\mu ^{}\stackrel{~}{f}_{ss}^{[2]}\left(\frac{5t_{if}^at_{lm}^a}{6}\frac{2\delta _{if}\delta _{lm}}{9}\right)\frac{(p,P^{}P)}{s}P_\mu \right)l_\mu ,$$
(18)
where $`\stackrel{~}{f}_{ss}^{[1]}`$ and $`\stackrel{~}{f}_{ss}^{[2]}`$ are equal to
$`\stackrel{~}{f}_{ss}^{[1]}`$ $`=`$ $`M\left({\displaystyle \frac{\alpha _s\left(\frac{m_1^2}{M^2}s\right)}{m_1^2}}+{\displaystyle \frac{\alpha _s\left(\frac{m_2^2}{M^2}s\right)}{m_2^2}}\right)\alpha _s\left(s\right){\displaystyle \frac{2M^3}{s}}\left(\alpha _s\left({\displaystyle \frac{m_1^2}{M^2}}s\right){\displaystyle \frac{m_2}{m_1^3}}+\alpha _s\left({\displaystyle \frac{m_2^2}{M^2}}s\right){\displaystyle \frac{m_1}{m_2^3}}\right)\alpha _s\left(s\right),`$ (19)
$`\stackrel{~}{f}_{ss}^{[2]}`$ $`=`$ $`{\displaystyle \frac{M^5}{m_1^3m_2^3}}\alpha _s\left({\displaystyle \frac{m_1^2}{M^2}}s\right)\alpha _s\left({\displaystyle \frac{m_2^2}{M^2}}s\right),`$ (20)
and $`p`$ is the four-momentum of annihilating quark.
For the differential cross section $`d\sigma _{ss}/d\mathrm{cos}\theta `$ we get the following expression:
$`{\displaystyle \frac{d\sigma _{ss}}{d\mathrm{cos}\theta }}`$ $`=`$ $`{\displaystyle \frac{8\pi ^3}{81s^3}}|\mathrm{\Psi }_s(0)|^4((2\stackrel{~}{f}_{ss}^{[1]}+{\displaystyle \frac{5}{4}}\stackrel{~}{f}_{ss}^{[2]}\sqrt{1{\displaystyle \frac{4M^2}{s}}}\mathrm{cos}\theta )^2+{\displaystyle \frac{(\stackrel{~}{f}_{ss}^{[2]})^2}{2}}(1{\displaystyle \frac{4M^2}{s}})\mathrm{cos}^2\theta )`$ (21)
$`\left(1{\displaystyle \frac{4M^2}{s}}\right)^{3/2}(1\mathrm{cos}^2\theta ).`$
The total cross section for the exclusive production of two heavy scalar diquarks in the $`q\overline{q}`$ annihilation is given by the formula
$$\sigma _{ss}=\frac{8\pi ^3}{81s^3}|\mathrm{\Psi }_s(0)|^4\left(1\frac{4M^2}{s}\right)^{3/2}\left(\frac{16}{3}(\stackrel{~}{f}_{ss}^{[1]})^2+\frac{11}{20}\left(1\frac{4M^2}{s}\right)(\stackrel{~}{f}_{ss}^{[2]})^2\right).$$
(22)
### 4.2 Annihilation into axial and scalar diquarks
The matrix element for the pair production of axial antidiquark and scalar diquark may be represented in the following way:
$$_{as}=\frac{32\pi ^2}{s^3}\mathrm{\Psi }_s^{}(0)\mathrm{\Psi }_a(0)\left(\stackrel{~}{f}_{as}^{[1]}\frac{2t_{if}^at_{lm}^a}{3}ϵ_{\mu \alpha \beta \gamma }P_\beta e_\alpha q_\gamma \stackrel{~}{f}_{as}^{[2]}\left(\frac{5t_{if}^at_{lm}^a}{6}\frac{2\delta _{if}\delta _{lm}}{9}\right)ϵ_{\mu \nu \alpha \beta }q_\alpha e_\beta p_\nu \right)l_\mu ,$$
(23)
where $`\stackrel{~}{f}_{as}^{[1]}`$and $`\stackrel{~}{f}_{as}^{[2]}`$ are equal to
$`\stackrel{~}{f}_{as}^{[1]}`$ $`=`$ $`M^3\left({\displaystyle \frac{\alpha _s\left(\frac{m_1^2}{M^2}s\right)}{m_1^3}}{\displaystyle \frac{\alpha _s\left(\frac{m_2^2}{M^2}s\right)}{m_2^3}}\right)\alpha _s(s),`$ (24)
$`\stackrel{~}{f}_{as}^{[2]}`$ $`=`$ $`\alpha _s\left({\displaystyle \frac{m_1^2}{M^2}}s\right)\alpha _s\left({\displaystyle \frac{m_2^2}{M^2}}s\right){\displaystyle \frac{M^5(m_2m_1)}{m_1^3m_2^3}}.`$ (25)
The differential cross section $`d\sigma _{as}/d\mathrm{cos}\theta `$ are given by
$`{\displaystyle \frac{d\sigma _{as}}{d\mathrm{cos}\theta }}`$ $`=`$ $`{\displaystyle \frac{64\pi ^3}{81s^4}}|\mathrm{\Psi }_s(0)\mathrm{\Psi }_a(0)|^2\sqrt{1{\displaystyle \frac{4M^2}{s}}}({\displaystyle \frac{(\stackrel{~}{f}_{as}^{[1]})^2}{2}}(1{\displaystyle \frac{4M^2}{s}})(1+\mathrm{cos}^2\theta )`$
$`+{\displaystyle \frac{33(\stackrel{~}{f}_{as}^{[2]})^2}{16}}(1+{\displaystyle \frac{s4M^2}{8M^2}}\mathrm{sin}^2\theta )+{\displaystyle \frac{5}{2}}\stackrel{~}{f}_{as}^{[1]}\stackrel{~}{f}_{as}^{[2]}\sqrt{1{\displaystyle \frac{4M^2}{s}}}\mathrm{cos}\theta ).`$
The total cross section for the exclusive production of heavy axial antidiquark and scalar diquark pairs in the $`q\overline{q}`$ annihilation equals
$$\sigma _{as}=\frac{64\pi ^3}{243s^4}|\mathrm{\Psi }_s(0)\mathrm{\Psi }_a(0)|^2\sqrt{1\frac{4M^2}{s}}\left(4(\stackrel{~}{f}_{as}^{[1]})^2\left(1\frac{4M^2}{s}\right)+\frac{33(\stackrel{~}{f}_{as}^{[2]})^2}{8}\left(2+\frac{s}{4M^2}\right)\right).$$
(27)
Here we have also neglected the difference between masses of axial and scalar diquarks.
### 4.3 Annihilation into pairs of axial diquarks
The amplitude for the exclusive pair production of two axial diquarks in the process may be represented as follows:
$`_{aa}`$ $`=`$ $`{\displaystyle \frac{32\pi ^2i}{3}}{\displaystyle \frac{\alpha _s(4m_1^2)\alpha _s(4m_2^2)}{m_1^3m_2^3s^3}}M^5|\mathrm{\Psi }_a(0)|^2\{((Pe^{})((pP^{})(le)+(lP^{})(pe))`$
$`+(P^{}e)\left((Pp)(le^{})+(lP)(pe^{})\right)+(ee^{})\left((lP)(pP^{})+(lP^{})(pP)\right)`$
$`+{\displaystyle \frac{s}{2}}((pe)(le^{})+(le)(pe^{})))(t_{if}^at_{lm}^a{\displaystyle \frac{4}{15}}\delta _{if}\delta _{lm})`$
$`+(\stackrel{~}{f}_{aa}^{[2]}((el)(Pe^{})(P^{}e)(e^{}l))+\stackrel{~}{f}_{aa}^{[1]}(Pl)(ee^{}))\left(t_{if}^at_{lm}^a\right)\}.`$
The differential cross section $`d\sigma _{aa}/d\mathrm{cos}\theta `$ is given by the following expression:
$`{\displaystyle \frac{d\sigma _{aa}}{d\mathrm{cos}\theta }}`$ $`=`$ $`{\displaystyle \frac{200\pi ^3}{81}}\alpha _s^2(4m_1^2)\alpha _s^2(4m_2^2){\displaystyle \frac{M^{10}}{m_1^6m_2^6s^7}}|\mathrm{\Psi }_a(0)|^4\sqrt{1{\displaystyle \frac{4M^2}{s}}}`$ (29)
$`(a_4\mathrm{cos}^4\theta +a_3\mathrm{cos}^3\theta +a_2\mathrm{cos}^2\theta +a_1\mathrm{cos}\theta +a_0).`$
The total cross section for the exclusive production of heavy axial diquarks pairs in the $`q\overline{q}`$ annihilation has the form
$$\sigma _{aa}=\frac{200\pi ^3}{81}\alpha _s^2(4m_1^2)\alpha _s^2(4m_2^2)\frac{M^{10}}{m_1^6m_2^6s^7}|\mathrm{\Psi }_a(0)|^4\sqrt{1\frac{4M^2}{s}}\left(\frac{2}{5}a_4+\frac{2}{3}a_2+2a_0\right),$$
(30)
where we have used the following definitions:
$`\stackrel{~}{f}_{aa}^{[1]}`$ $`=`$ $`{\displaystyle \frac{\alpha _s\left(s\right)}{\alpha _s(4m_1^2)\alpha _s(4m_2^2)}}{\displaystyle \frac{8m_1^3m_2^3}{5M^2}}\left({\displaystyle \frac{\alpha _s(4m_1^2)m_2}{m_1^3}}+{\displaystyle \frac{\alpha _s(4m_2^2)m_1}{m_2^3}}\right),`$ (31)
$`\stackrel{~}{f}_{aa}^{[2]}`$ $`=`$ $`{\displaystyle \frac{\alpha _s\left(s\right)}{\alpha _s(4m_1^2)\alpha _s(4m_2^2)}}{\displaystyle \frac{4m_1^3m_2^3}{5M}}\left({\displaystyle \frac{\alpha _s(4m_1^2)}{m_1^3}}+{\displaystyle \frac{\alpha _s(4m_2^2)}{m_2^3}}\right),`$ (32)
$`a_4`$ $`=`$ $`{\displaystyle \frac{99}{400}}s^2(s4M^2)^2,`$ (33)
$`a_3`$ $`=`$ $`{\displaystyle \frac{1}{8}}{\displaystyle \frac{s^3}{M^2}}\left(2\stackrel{~}{f}_{aa}^{[2]}s+\stackrel{~}{f}_{aa}^{[1]}(6M^2+s)\right)\left(14{\displaystyle \frac{M^2}{s}}\right)^{3/2},`$ (34)
$`a_2`$ $`=`$ $`{\displaystyle \frac{1}{16M^4}}s(s4M^2)\{4\stackrel{~}{f}_{aa}^{[1]}\stackrel{~}{f}_{aa}^{[2]}s(s2M^2)+(\stackrel{~}{f}_{aa}^{[1]})^2(12M^44M^2s+s^2)`$ (35)
$`+s({\displaystyle \frac{33}{25}}(12M^6M^4s)+4(\stackrel{~}{f}_{aa}^{[2]})^2(s2M^2))\},`$
$`a_1`$ $`=`$ $`{\displaystyle \frac{s^2}{8M^2}}\sqrt{14{\displaystyle \frac{M^2}{s}}}\left(2\stackrel{~}{f}_{aa}^{[2]}s(s+4M^2)+\stackrel{~}{f}_{aa}^{[1]}(24M^4+2M^2s+s^2)\right),`$ (36)
$`a_0`$ $`=`$ $`{\displaystyle \frac{s}{16M^4}}(4\stackrel{~}{f}_{aa}^{[1]}\stackrel{~}{f}_{aa}^{[2]}s(8M^46M^2s+s^2)+(\stackrel{~}{f}_{aa}^{[1]})^2(48M^6+28M^4s8M^2s^2+s^3)`$ (37)
$`+2s({\displaystyle \frac{33}{25}}M^4s(4M^2+s)+\stackrel{~}{f}_{aa}^{[2]}(16M^44M^2s+2s^2))).`$
### 4.4 The diquark with two identical particles
In this section we consider the $`q\overline{q}`$ annigilation into pair of diquarks composed of two identical heavy quarks. So, as in $`e^+e^{}`$ annihilation, we have to take into account the axial diquark production, only.
The amplitude has the form
$`_{aa}`$ $`=`$ $`{\displaystyle \frac{512\pi ^2i}{3}}{\displaystyle \frac{\alpha _s^2(M^2)}{Ms^3}}|\mathrm{\Psi }_a(0)|^2\{((Pe^{})((pP^{})(le)+(lP^{})(pe))`$ (38)
$`+(P^{}e)\left((Pp)(le^{})+(lP)(pe^{})\right)+(ee^{})\left((lP)(pP^{})+(lP^{})(pP)\right)`$
$`+{\displaystyle \frac{s}{2}}((pe)(le^{})+(le)(pe^{})))(t_{if}^at_{lm}^a{\displaystyle \frac{4}{15}}\delta _{if}\delta _{lm})`$
$`+(\stackrel{~}{f}_{aa}^{[2]}((el)(Pe^{})(P^{}e)(e^{}l))+\stackrel{~}{f}_{aa}^{[1]}(Pl)(ee^{}))\left(t_{if}^at_{lm}^a\right)\}.`$
The differential cross section $`d\sigma _{aa}/d\mathrm{cos}\theta `$ equals
$$\frac{d\sigma _{aa}}{d\mathrm{cos}\theta }=2^{11}\frac{25\pi ^3}{81}\frac{\alpha _s^4(M^2)}{M^2s^7}|\mathrm{\Psi }_a(0)|^4\sqrt{1\frac{4M^2}{s}}(a_4\mathrm{cos}^4\theta +a_3\mathrm{cos}^3\theta +a_2\mathrm{cos}^2\theta +a_1\mathrm{cos}\theta +a_0).$$
(39)
The total cross section for the exclusive production is given by
$$\sigma _{aa}=2^{11}\frac{25\pi ^3}{81}\frac{\alpha _s^4(M^2)}{M^2s^7}|\mathrm{\Psi }_a(0)|^4\sqrt{1\frac{4M^2}{s}}\left(\frac{2}{5}a_4+\frac{2}{3}a_2+2a_0\right),$$
(40)
where we have used the following definitions:
$`\stackrel{~}{f}_{aa}^{[1]}`$ $`=`$ $`{\displaystyle \frac{\alpha _s\left(s\right)}{\alpha _s(M^2)}}{\displaystyle \frac{M^2}{5}},`$ (41)
$`\stackrel{~}{f}_{aa}^{[2]}`$ $`=`$ $`{\displaystyle \frac{\alpha _s\left(s\right)}{\alpha _s(M^2)}}{\displaystyle \frac{M^2}{5}}.`$ (42)
## 5 Numerical estimates
The cross section ratios of the exclusive heavy diquark pair production to the respective heavy quark cross section in $`e^+e^{}`$ annihilation,
$$\sigma (e^+e^{}Q\overline{Q})=\frac{4\pi \alpha _{em}^2q_Q^2}{s}\sqrt{1\frac{4m_Q^2}{s}}\left(1+\frac{2m_Q^2}{s}\right),$$
are shown in Figs. 4a, 4b.
We see that the most optimistic expectations for the pair production is given for the $`cc`$-diquarks at $`B`$-factories, when the hogh luminosities can yield several thousands pairs of doubly charmed baryons.
The cross section ratios of the exclusive heavy diquark pair production to the respective heavy quark cross section in $`q\overline{q}`$ annihilation,
$$\sigma (q\overline{q}Q\overline{Q})=\frac{8\pi \alpha _s^2}{27s}\sqrt{1\frac{4m_Q^2}{s}}\left(1+\frac{2m_Q^2}{s}\right)$$
are shown in Figs. 4c, 4d.
Figs. 4a, 4c represent the case of diquark composed of different quarks and Figs. 4b, 4d give the case of diquark composed of identical quarks. We put
$`m_c`$ $`=`$ $`1.55GeV,`$
$`m_b`$ $`=`$ $`4.9GeV,`$ (43)
$`\mathrm{\Lambda }_{QCD}`$ $`=`$ $`0.2GeV.`$
The values of diquark wave functions at the origin have been taken from .
In the quark-antiquark annihilation the pair production of doubly charmed baryons gives about $`10^5`$ fraction of charm yield at $`\sqrt{s}<100`$ GeV in the hadron collisions, while at higher energies the gluon-gluon fusion will dominate. In experoments at fixed target the threshold behaviour results in an additional suppression.
The inclusive production of doubly heavy baryons was considered in , so that we see that the pair production results in the suppression factor about 0.1.
## 6 Conclusion
In this paper we have considered the exclusive production of doubly heavy diquark pairs for the axial-axial, scalar-scalar and axial-scalar states in the framework of the constituent quark model. The matrix elements, differential and total cross sections for the processes of $`e^+e^{}`$ and $`q\overline{q}`$ annihilation have been given. We have calculated the pair production of diquarks composed of identical heavy quarks. The obtained formulae can be used in the calculation of cross sections for the diquark fragmentation into the doubly heavy baryons in the processes of $`e^+e^{}`$ annihilation and proton-antiproton inelastic collisions.
We have found that the yield of doubly heavy baryon pairs could really reach $`10^3÷10^6`$ depending on the energies and luminosities of accelerators, $`e^+e^{}`$ colliders and fixed target experiments with hadron beams.
The authors thank Dr. V.V.Kiselev, who asked us for the calculations of pair production of doubly heavy diquarks and explained some points in the problem.
This work is in part supported by the Russian Foundation for Basic Research, grants 99-02-16558 and 00-15-96645.
## Figure Captions
* The diagrams of $`e^+e^{}`$ ($`q\overline{q}`$) annihilation into the pair of doubly heavy diquarks due to the single photon (gluon) exchange.
* The $`q\overline{q}`$ annihilation into the diquarks due to the three-gluon interaction.
* The diagrams of double gluon annihilation of light quarks into the pair of doubly heavy diquarks.
* The ratios of total cross sections for the exclusive production of $`bc`$ diquark pair and the $`c\overline{c}`$ production in $`e^+e^{}`$ annihilation for the axial-axial, axial-scalar and scalar scalar states.
* The ratios of cross sections for the exclusive production of $`bb`$ and $`cc`$ diquark pair and the respective $`b\overline{b}`$ and $`c\overline{c}`$ production in $`e^+e^{}`$ annihilation.
* The ratios of cross sections for the exclusive production of $`bc`$ diquark pair and the $`c\overline{c}`$ production in $`q\overline{q}`$ annihilation. The notations are the same as in Fig. 4a.
* The ratios of cross sections for the the exclusive production of $`bb`$ and $`cc`$ diquark pair and the respective $`b\overline{b}`$ and $`c\overline{c}`$ production in $`q\overline{q}`$ annihilation.
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# On space-times admitting shear-free, irrotational, geodesic null congruences
## 1 Introduction
In this article we wish to extend earlier work on shear-free, irrotational and geodesic (SIG) timelike and spacelike congruences to SIG null congruences. The fact that we are dealing with null congruences means that we have to approach the problem in a completely different way; we must make extensive use of the Newman-Penrose formalism.
Thus, we wish to study a congruence of curves whose tangent vector $`𝐤`$ is null and geodesic. Hence, we have a family of null geodesics $`x^a=x^a(y^\alpha ,v)`$, where $`y^\alpha `$ distinguishes the different geodesics, and $`v`$ is the affine parameter along a fixed geodesic. The null tangent vector is $`k^a=\frac{x^a}{v}`$, and satisfies $`k_{}^{a}{}_{;b}{}^{}k^b=0`$. The spin coefficients are defined in , where $`\rho =(\theta +i\omega )`$ is called the complex divergence and $`\sigma `$ is the complex shear. The geodesic condition implies that the spin coefficient $`\kappa `$ vanishes and $`ϵ+\overline{ϵ}=0`$ follows from the choice of an affine parameter along the congruence. The congruence is said to be shear-free if $`\sigma =0`$. Also, from the relation $`k_{[a;b}k_{c]}=(\overline{\rho }\rho )\overline{m}_{[a}m_bk_{c]}`$ , it follows that $`w=0`$ (i.e., zero twist) is a necessary and sufficient condition for $`𝐤`$ to be hypersurface orthogonal.
First we shall briefly review some of the results of relevance to this work. Goldberg and Sachs proved that if a gravitational field contains a shear-free, geodesic, null congruence $`𝐤`$, then $`\kappa =\sigma =0`$, and if
$$R_{ab}k^ak^b=R_{ab}k^am^b=R_{ab}m^am^b=0,$$
(1)
then the field is algebraically special (i.e., $`\mathrm{\Psi }_0=\mathrm{\Psi }_1=0`$), and $`𝐤`$ is a degenerate eigendirection. In addition, a vacuum metric is algebraically special if and only if it contains a shear-free geodesic null congruence.
A space-time admits a geodesic, shear-free, twist-free ($`\kappa =\sigma =\omega =0`$) and diverging ($`\rho =\overline{\rho }=\theta =1/r`$) null congruence $`𝐤`$, and satisfies (1), if and only if the metric can be written in the form
$$ds^2=2r^2P^2(z,\overline{z},u)dzd\overline{z}2dudr2H(z,\overline{z},r,u)du^2.$$
(2)
Robinson-Trautman models with this metric have been found for vacuum, Einstein-Maxwell and pure radiation fields with or without a cosmological constant .
For geodesic null vector fields we have that $`(\theta +i\omega )_{,a}k^a+(\theta +i\omega )^2+\sigma \overline{\sigma }=R_{ab}k^ak^b/2`$. Therefore, in the non-diverging case (i.e., $`\rho =(\theta +i\omega )=0`$), if the energy condition $`T_{ab}k^ak^b0`$ is satisfied, it follows that $`\sigma =0=R_{ab}k^ak^b`$. Thus, non-twisting (and therefore geodesic) and non-expanding null congruences must be shear-free. Hence, the space-time is algebraically special, and it corresponds to vacuum, Einstein-Maxwell, and pure radiation field. Perfect fluid solutions violate $`R_{ab}k^ak^b=0`$ unless $`\mu +p=0`$. This class of solutions has been studied by Kundt .
Another important case corresponds to the Kerr-Schild metric, which is given by $`g_{ab}=\eta _{ab}2\varphi k_ak_b`$. The null vector $`𝐤`$ of a Kerr-Schild metric is geodesic if and only if the energy-momentum tensor obeys the condition $`T_{ab}k^ak^b=0`$, and then $`𝐤`$ is a multiple principal null direction of the Weyl tensor and the space-time is algebraically special. The general properties of the Kerr-Schild metrics and their applications to vacuum, Einstein-Maxwell, and pure radiation space-times can be found in .
Finally, we note the algebraically special perfect fluid space-times corresponding to the generalized Robinson-Trautman solutions investigated by Wainwright . They are characterized by a multiple null eigenvector $`𝐤`$ of the Weyl tensor which is geodesic, shear-free, and twist-free but expanding (i.e., $`\mathrm{\Psi }_o=\mathrm{\Psi }_1=0`$, $`\kappa =\sigma =\omega =0`$, $`\rho =\overline{\rho }0`$), and the four-velocity obeys $`u_{[a;b}u_{c]}=0`$, $`k_{[c}k_{a];b}u^b=0`$. The line-element of the space-time can be written in the form
$$ds^2=\frac{1}{2}\chi ^2(r,u)P^2(z,\overline{z},u)dzd\overline{z}+2du(drUdu),$$
(3)
with
$$U=r(\mathrm{ln}P)_{,u}+U^0(z,\overline{z},u)+S(r,u),\chi _{,r}0,\frac{\chi _{,rr}}{\chi }0.$$
(4)
In this case no dust solutions nor solutions of Petrov types $`III`$ and $`N`$ are possible.
## 2 Analysis
Let us consider space-times admitting a shear-free, irrotational, geodesic null congruence in which the source of the gravitational field is a combination of a perfect fluid and null radiation, so that the energy-momentum tensor has the form
$$T_{ab}=(\mu +p)u_au_bpg_{ab}+\varphi ^2k_ak_b,$$
(5)
where $`u^a`$ is the four-velocity of the fluid, $`\mu `$ and $`p`$ are the density and the pressure of the fluid, respectively, and $`𝐤`$ is a null vector. The null radiation is geodesic, twist-free, and shear-free, and defines the null congruence. Wainwright proved that for a space-time in which there exists a SIG null congruence, coordinates can be chosen so that the metric takes on the simplified form (3) with $`u=x^1`$, $`r=x^2`$, $`z=x^3+ix^4`$, the tangent field of the null congruence is given by $`k^a=\delta _2^a`$, $`k_a=\delta _a^1`$, and we can introduce the null tetrad
$`k^a=\delta _r^a,`$ $`l^a=\delta _u^a+U\delta _r^a,`$ $`m^a=P\chi ^1(\delta _3^a+i\delta _4^a),`$ (6)
$`k_a=\delta _a^u,`$ $`l_a=U\delta _a^u+\delta _a^r,`$ $`m_a=P^1\chi (\delta _a^3+i\delta _a^4)/2.`$ (7)
With the sign convention used here we have that $`u^au_a=k^al_a=1=m^a\overline{m}_a`$. Note that the null radiation is everywhere tangent to the repeated null congruence of the space-time.
First, since $`\mathrm{\Phi }_{01}\frac{1}{2}R_{ab}k^am^b=0`$, we conclude that the four-velocity satisfies $`u^am_a=0`$, and hence it can be expressed in terms of the null tetrad by
$$u^a=\frac{1}{\sqrt{2}B}(B^2k^a+l^a)\mathrm{and}u_a=\frac{1}{\sqrt{2}B}[(B^2U)\delta _a^u+\delta _a^r],$$
(8)
for some function $`B`$. The conditions $`\mathrm{\Phi }_{02}\frac{1}{2}R_{ab}m^am^b=0`$ and $`\mathrm{\Phi }_{12}\frac{1}{2}R_{ab}m^al^b=0`$ are satisfied identically. The non-zero components of the Ricci tensor are
$`\mathrm{\Phi }_{00}{\displaystyle \frac{1}{2}}(R_{ab}{\displaystyle \frac{1}{4}}Rg_{ab})k^ak^b={\displaystyle \frac{1}{2}}(\mu +p)(𝐤𝐮)^2,`$ (9)
$`\mathrm{\Phi }_{11}{\displaystyle \frac{1}{4}}(R_{ab}{\displaystyle \frac{1}{4}}Rg_{ab})(k^al^b+m^a\overline{m}^b)={\displaystyle \frac{1}{4}}(\mu +p)(𝐤𝐮)(𝐥𝐮),`$ (10)
$`\mathrm{\Phi }_{22}{\displaystyle \frac{1}{2}}(R_{ab}{\displaystyle \frac{1}{4}}Rg_{ab})l^al^b={\displaystyle \frac{1}{2}}(\mu +p)(𝐥𝐮)^2+{\displaystyle \frac{1}{2}}\varphi ^2.`$ (11)
In addition, since $`𝐤𝐮=\frac{1}{\sqrt{2}B}`$ and $`𝐥𝐮=\frac{1}{\sqrt{2}}B`$ implies $`𝐥𝐮=B^2(𝐤𝐮)`$, we obtain
$`B^2\mathrm{\Phi }_{00}`$ $`=`$ $`2\mathrm{\Phi }_{11},`$ (12)
$`B^4\mathrm{\Phi }_{00}`$ $`=`$ $`\mathrm{\Phi }_{22}{\displaystyle \frac{1}{2}}\varphi ^2.`$ (13)
If we now assume that the fluid is non-rotating, then $`B^2=U+F(r,u)`$, and the compatibility condition (12) can be written as
$$(U+F)\mathrm{\Phi }_{00}=2\mathrm{\Phi }_{11}.$$
(14)
On differentiating this equation successively with respect to $`z`$ and $`r`$, we obtain the restriction
$$(\chi ^2)_{,rrr}[U_{}^{0}{}_{,z}{}^{}+r(\mathrm{ln}P)_{,uz}]=0.$$
(15)
There are consequently two different cases to consider.
In the first case $`U_{}^{0}{}_{,z}{}^{}+r(\mathrm{ln}P)_{,uz}=0`$, which is equivalent to $`U_{}^{0}{}_{,z}{}^{}=(\mathrm{ln}P)_{,uz}=0`$, so that $`P=P(z,\overline{z})`$ and $`U^0=U^0(u)`$. Obviously, the solutions admit a multiply transitive group of motions, $`G_3`$, acting on the 2-spaces $`r=`$const, $`u=`$const, of constant curvature, and belong to class $`II`$ of Stewart and Ellis . The metric (3) can then be rewritten as
$$ds^2=\chi ^2(r,u)\frac{2dzd\overline{z}}{(1+\frac{k}{2}z\overline{z})^2}+2du(drU(r,u)du).$$
(16)
The non-zero Ricci components are given by
$`\mathrm{\Phi }_{00}={\displaystyle \frac{\chi _{,rr}}{\chi }},`$ (17)
$`\mathrm{\Phi }_{11}={\displaystyle \frac{\chi _{,r}\chi _{,u}}{2\chi ^2}}+{\displaystyle \frac{(\chi _{,r})^2U}{2\chi ^2}}{\displaystyle \frac{U_{,rr}}{4}}+{\displaystyle \frac{k}{4\chi ^2}},`$ (18)
$`\mathrm{\Phi }_{22}={\displaystyle \frac{\chi _{,u}U_{,r}}{\chi }}{\displaystyle \frac{\chi _{,uu}}{\chi }}2{\displaystyle \frac{\chi _{,ur}U}{\chi }}{\displaystyle \frac{\chi _{,r}U_{,u}}{\chi }}{\displaystyle \frac{\chi _{,rr}U^2}{\chi }},`$ (19)
and the Ricci scalar is given by
$$\frac{R}{2}=12\mathrm{\Lambda }=4\frac{\chi _{,r}U_{,r}}{\chi }+2\frac{\chi _{,r}\chi _{,u}}{\chi ^2}+2\frac{(\chi _{,r})^2U}{\chi ^2}+4\frac{\chi _{,ur}}{\chi }+U_{,rr}+4\frac{\chi _{,rr}U}{\chi }+\frac{k}{\chi ^2}.$$
(20)
Hence, the metric (16) can be interpreted as pure radiation plus a perfect fluid where $`\mu `$ and $`p`$ are given by
$$\mu =\frac{R}{4}+6\mathrm{\Phi }_{11},p=\frac{R}{4}+2\mathrm{\Phi }_{11},$$
(21)
$`u_a`$ is determined by (8) with $`B^2=2\mathrm{\Phi }_{11}/\mathrm{\Phi }_{00}`$, and $`\varphi ^2`$ is given by
$$\varphi ^2=2\left(\mathrm{\Phi }_{22}4\frac{\mathrm{\Phi }_{11}^2}{\mathrm{\Phi }_{00}}\right).$$
(22)
In the second case (i.e., $`\chi _{}^{2}{}_{,rrr}{}^{}=0`$) two possibilities arise:
$`(i)`$ $`\chi ^2=ϵr,`$ $`ϵ=\pm 1`$ (23)
$`(ii)`$ $`\chi ^2=ϵ(r^2k^2),`$ $`k=const.`$ (24)
In both subcases $`\chi =\chi (r)`$, and they can be written together as $`\chi ^2=ar^2+2br+c`$, with $`a`$, $`b`$, $`c`$ taken to be appropriate constants. From equation (14) we obtain
$$aU^0b(\mathrm{ln}P)_{,u}+K=G(u),$$
(25)
and
$$\frac{1}{2}[\chi ^2S_{,r}S(\chi ^2)_{,r}]_{,r}+\frac{F\mathrm{\Sigma }}{\chi ^2}=G(u),$$
(26)
where $`K4P^2(\mathrm{ln}P)_{z\overline{z}}`$, $`\mathrm{\Sigma }b^2ac`$, and $`G(u)`$ is an arbitrary function of $`u`$.
Subcase $`(i)`$: $`a=c=0`$, $`b=ϵ/2`$. Integrating equation (26) we see that $`S`$ can be written in the form
$$S=rh(u)+2ϵG(u)r\mathrm{ln}rf(u)\frac{1}{2}r\frac{dr}{r^2}^r\frac{d\widehat{r}}{\widehat{r}}F(\widehat{r},u),$$
(27)
where $`h(u)`$ and $`f(u)`$ are arbitrary functions of $`u`$.
Subcase $`(ii)`$: $`a=ϵ`$, $`b=0`$, $`c=ϵk^2`$, $`\mathrm{\Sigma }=k^2`$. We obtain
$$S=ϵG(u)+f(u)\chi ^2\frac{dr}{\chi ^4}+h(u)\chi ^22k^2\chi ^2\frac{dr}{\chi ^4}^r\frac{d\widehat{r}}{\chi ^2(\widehat{r})}F(\widehat{r},u).$$
(28)
Therefore, the metric (3) with $`\chi (r)`$ given by (23) or (24), $`S(r,u)`$ given by (27) or (28), and $`P(z,\overline{z},u)`$ satisfying (25) can be interpreted as pure radiation plus a perfect fluid, in which the four-velocity is determined by (8) and $`\varphi ^2`$, $`\mu `$ and $`p`$ are determined by (21) and (22), respectively.
## Acknowledgments
This work was supported by the European Union, TMR Contract No. ERBFMBICT961479 (AMS), the Natural Sciences and Engineering Research Council of Canada (AAC) and the Canadian Institute for Theoretical Astrophysics (DJM).
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# A Continuous Time Asynchronous Model of the Stock Market;
## 1 The Market Mechanism
A trader in our model is not continuously active, but rather ”sleeps” most of the time. Of course the time intervals when nothing happens are not actually spent by the computer waiting idly: each operation/event in the virtual market world is performed by the computer in the appropriate time order and with the appropriate advancement of the ”market world clock” without actually waiting for the a corresponding real time to pass.
The trader wakes up at times which are the result of his own (previously taken) decisions: e.g. the trader may decide to wake-up authomatically at a pre-defined time interval after the last order or in reaction to specific market events (news, price changes). Once awake, the trader decides whether (s)he wishes to issue new Buy/Sell orders (see figure 1).
When a trader first enters the market, (s)he is given an initial amount of cash and shares. In the current version of the software implementing the model, a trader cannot have more than one order in the market at any given time (and cannot borrow or short-sale).
### 1.1 The Stacks Algorithm
In order to make the traders actions and decisions asynchronous and taking place at arbitrary continous time, we implemented the market using a set of stacks, that contain traders in different decision order stages (see figure 2):
* The ”Dormant Stack” \- A trader’s default state is dormant. Each trader defines the conditions on which (s)he should be waken up: (s)he can ask for instance to be waken up every time that a time K has passed since the last wake-up, or (s)he can ask to be waken up when something happens in the market (news arrived to the market, price has changed by more than $`5\%`$ etc.). External news events are given on an ordinal scale (e.g. -5 to 5), such that the trader can ask to be waken up only from a specific news level.
* The ”Decision Stack” \- Once the trader is awake, he is transferred to the ”Decision Stack”. Each trader has an individual delay period, during which he stays in this stack (and ”examines” the market). Once the delay period is over, the trader can decide to issue a new order, to change / delete an existing order (if it has not already been executed), or to go back to sleep. He can also change his waking parameters. An ”order” in our model is actually a ”LIMIT” order, i.e. a desired price and quantity of shares. Future versions may also allow ”MARKET” orders (i.e. buy/sell at the best market price).
* The ”Transferred Orders Stack” \- When a trader decides to issue a new order, the order is transferred to this stack before it enters the market (see 2.2 below). A transferred order is delayed in this stack for a constant period of time (In future version each trader will have an individual delay period, since some traders may be ”closer” to the market).
Note that when running a simulation of the model there is no relation between computer time intervals and simulated time interval, since computer operations are made as soon as possible without waiting.
### 1.2 Matching Orders
Once an order leaves the ”transferred orders” stack, it enters either the BUY stack or the SELL stack of the market matching orders mechanism. The orders in these stacks are sorted according to price.
If the minimal price in the SELL stack is lower than the maximal price in the BUY stack:
MIN(SELL Order Price) $`<`$ MAX (BUY Order Price)
then a transaction of shares-for-cash will be executed between the issuers of these orders. Currently, short selling of shares is not allowed in the model.
Of the two orders that were executed, the first that was issued to the market also determines the new market price. If the two orders differ in size (number of shares), the number of shares in the transaction equals the minimum size of the two orders. The order that was not fully executed is returned to the ”transferred orders” stack, with a decreased size.
In future versions of the program, we may limit the Market Orders Stacks in either size (worst orders are cancelled) or time (orders with expiry times).
### 1.3 The World Manager
All this process of waking up the traders, moving their orders between the different stacks and operating the market matching mechanism is performed by a central ”World manager” (or Market Manager and a wider sense). Other than being a technical aid to the process, the World manager has no influence on it whatsoever. In future versions of the program, the World manager will also have the ability to introduce news to the market.
## 2 Traders Decision Functions
### 2.1 Traders Types
When entering the ”Decision Stack”, a trader uses his individual decision function to process inputs (current price, historical price quotes, news, fundamental price) and produce a decision (Go back to sleep, Change waking parameters and issue an order).
The decision function can be programmed to include any decision making algorithm. The current version of the market simulation offers three types of pre-defined configurable decision functions:
* ”Random Trader” \- Issues buy/sell orders randomly in the proximity of the current market price.
* ”Fundamentalist Trader” \- decides according to a ”fundamental price”, which follows a random walk.
* ”Chartist Trader” \- Uses simple technical analysis (In our case - curve fitting for the last N price quotes).
Future versions of the model will include more complex ”psychological” features of traders , including the possibility to react to external news.
### 2.2 Notation
Following are descriptions of the three configurable decision functions. But first some notation:
When issuing an order, a trader has to decide on two parameters:
OrderP = The price the trader sets for his order
OrderShares = Shares to buy/sell (i.e. the size of the order)
The decision process also uses the following variables:
NewP = This is the trader’s best estimate for the market price in the next time tick.
MP(t), FP(t) = Market price and fundamental price (respectively) at the current time.
CurrShares = Current number of shares the trader owns.
CurrCash = Current amount of cash the trader has.
PriceTick = The market price is progressing in discrete ticks of size PriceTick (default: 1/32).
RAND\[a,b\] = A random number in the interval \[a,b\].
RANDN\[m,s\] = A number drawn from a gaussian distribution with mean m and standard deviation s.
### 2.3 Random Trader
In order to calculate the price for his order, a Random Trader chooses a random price near the current market price:
OrderP = MP(t) + MP(t) * Sigma * RAND\[-1,1\]
Where:
Sigma = Price variations around Market Price (configurable, default: $`10\%`$).
Now the Random Trader chooses randomly the desired value of shares he wishes to hold after the order is executed (in Currency value) :
InvestInShares = (OrderP * CurrShares + CurrCash) * RAND
He therefore needs to own this amount of shares:
InvestInShares / OrderP
So he will buy /sell (determined by the sign) the following number of shares:
OrderShares = $`|`$ CurrShares - InvestInShares / OrderP $`|`$
3.4 Fundamentalist Trader
The basis for the fundamentalist trader’s behavior is the fundamental price, FP, which is updated every pre-defined update period:
FP(t+1) = FP(t) * ( RANDN\[0,Sigma\] +1 )
Where
Sigma = Maximum price variations around FP (configurable, default: 0.4)
What the Fundamentalist Trader actually sees is a noisy Fundamental Price, bf FP’(t) : FP’(t) = FP(t) * ( 1+ RAND\[-1,1\] * Sigma )
Note: If FP’(t) $``$ 0, we set FP’(t) = FP(t).
The following ratio is used to determined how much the Fundamentalist is willing to trade (He is more willing to trade as the difference between Market Price and his individual Fundamental Price is higher):
RATIO = $`|`$ (MP(t) - FP’(t) ) / MP(t) $`|`$
Now the decision is easy:
If FP’(t) $`<`$ MP(t) ,
the trader wishes to sell, in a price as close as possible to the market price:
OrderP = MP(t) - PriceTick
If MP (t )= PriceTick then OrderP=MP(t )
OrderShares = RATIO * CurrShares.
If FP’(t) $`>`$ MP(t), the trader wishes to buy shares, in the best price possible:
OrderP = MP(t) + PriceTick
OrderShares = CurrCash * RATIO / OrderP
### 2.4 Chartist Trader
A chartist trader has a memory span of N (default is 3). He extrapolates the next market price (NewP) using a simple polynomial fit for the last N market price quotes he remembers (which means that these price quotes need not be three consecutive market price quotes).
Once NewP is determined:
If MP(t) $`>`$ NewP the trader wants to sell stocks:
OrderP = MP(t) - PriceTick
OrderShares = CurrShares * (OrderP - NewP) / MP(t)
If MP(t) $`<`$ NewP the trader wants to buy stocks:
OrderP = MP(t) + PriceTick
OrderShares = CurrCash * (NewP - OrderP) / (MP(t )\* OrderP )
## 3 Implementing the Model
Using the stacks system, our model can simulate a market with millions of traders acting continuously and asynchronously according to individual decision functions. There are no limitations on the number or types of decision functions that are used in the market.
We implemented the market model using a PC-based software. Here is a description of a typical market simulation: First, we need to introduce traders to the market. This is done by iteratively choosing a decision function and configuring it for a specific trader (see fig. 3). The decision functions to choose from are pre-programmed, and can vary from simple chartists to complicated ”psychological” decision functions. For the sake of simplicity, we allow several traders to use the same configured decision function (forming a ”group”), e.g. 10,000 traders that use a fundamentalist decision function with the same parameters. However, if the decision function contains random variables (as most do), each trader in a group will act independently in the market, with no correlation to the other members of the group (besides his decision function). The market in the current example (figure 3) consists of 6 different groups of fundamentalists (some are ”richer”, i.e. have more cash/shares when entering the market, some fluctuate more around the fundamental price), 3 groups of chartists (”richer”, ”longer memory”, etc.) and two groups of random traders. Thus, a total of 161,014 traders were introduced to the market in 11 different groups.
After choosing the traders, the market simulation is initiated. The market form (see fig. 4) shows market parameters in real-time, including market price and data on the traders that are at the top of the BUY and SELL stacks. One can get simple graphs and statistics from the user interface, and a more thorough analysis using a direct link to matlab. Future versions of the software will also allow tracking a single trader’s behavior in the market (current version allows it only manually). Figure 5 shows the price behavior of the simulation under the current settings.
## 4 Concluding Remarks
The continuous asynchronous model’s main advantages are its flexibility, robustness and efficiency. The model poses no limitation on the number of traders that enter the market, while allowing each trader to use a unique decision function. Moreover, the program implementing the model is fast enough to simulate a year of stock market activity with millions of traders in a matter of seconds. Results of the model will be published in subsequent papers.
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# Untitled Document
KCL-MTH-00-33
CERN-TH/2000-147
hep-th/0005270
Hidden Superconformal Symmetry in M Theory
P. West
Department of Mathematics
King’s College, London, UK
and
Theory Division
CERN
CH1211, Geneva 23
Abstract
The bosonic sectors of the eleven dimensional and IIA supergravity theories are derived as non-linear realisations. The underlying group includes the conformal group, the general linear group and as well as automorphisms of the supersymmetry algebra. We discuss the supersymmetric extension and argue that Osp(1/64) is a symmetry of M theory. We also derive the effective action of the closed bosonic string as a non-linear realisation.
1. Introduction
Much of the progress in recent years in our understanding of the non-perturbative effects of string theory has relied on the structure of the supergravity theories in eleven and ten dimensions. While there is only one eleven dimensional supergravity theory , in ten dimensions there exist the IIA and IIB supergravity theories as well as the type I supergravity theory coupled to the Yang-Mills theory . These theories are essentially uniquely determined by the type of supersymmetry that they possess. Hence given a string theory in ten dimensions its complete low energy effective action must be the supergravity with the same space-time supersymmetry. One intriguing feature of supergravity theories is the occurance of coset space symmetries that control the way the scalars in these theories behave. The four dimensional $`N=4`$ supergravity theory possess a SL(2,R)/ U(1) symmetry , the IIA theory a SO(1,1) symmetry , the IIB theory a SL(2,R)/ U(1) and the further reductions of the eleven dimensional supergravity theory possess cosets based on the exceptional groups . These symmetries have also played an important role in string dualities in recent years and any further elucidation of the symmetries of supergravity theories could prove useful.
Although these symmetries can be viewed as a consequence of supersymmetry, it is desirable to have a deeper understanding. One step in this direction has been the extension of the coset space description of the scalars to include the gauge fields . This has been achieved by introducing a group with Grassmann odd as well as Grassmann even generators. All of these generators are scalars under the Lorentz group and the indices of the gauge fields are incorporated by writing them as forms. The group elements $`g`$ of the coset are then exponentials of these forms each of which is multiplied by a generator with the corresponding Grassmann character. This has the advantage that one automatically finds the gauge field strengths when taking the Cartan forms, $`g^1dg`$ using the method of non-linear realisations . The result is an elegant formulation of these sectors of the supergravity theories, but it is not apparent how this method can be extended to include the graviton or indeed the fermionic sectors of the theory.
Recently, it was shown that part of the $`GL(32)`$ automorphism group of the supersymmetry algebra was found to be a symmetry of the fivebrane equations of motion. This symmetry was also found to play an important role in formulating the branes of M theory in terms of a non-linear realisation, indeed the world-volume gauge field strengths are the Goldstone bosons for part of this automorphism symmetry . It was conjectured that this symmetry could play a role in M theory and should occur in eleven dimensional supergravity.
Long ago , Ogievetsky realised that the group of general coordinate realisations was the closure of the conformal group and the group of affine transformations, IGL(4) in four dimensions. As a consequence, in reference , gravity was reformulated as the non-linear realisation of these two groups. Subsequently, it was shown that the Sokatchev-Ogievetsky superspace formulation of the $`N=1`$ supergravity in four dimensions could be expressed as a non-linear realisation
In this paper we wish to revive this old idea of realising gravity as a non-linear realisation and, by combining it with the presence of the automorphism symmetries, show that the bosonic sector of eleven dimensional supergravity can be expressed as a non-linear realisation. In particular, in section two, we will show how the bosonic sector of sector of eleven-dimensional supergravity, that is the graviton and the rank three gauge field, is a non-linear realisation of the conformal group, SO(2,11) and a group which is generated by the generators of the group of affine transformations IGL(11) and two further generators which are of rank three and rank six. While the graviton is the Goldstone boson for the group GL(11), the gauge field and its dual are the Goldstone bosons associated with these two additional generators. We will argue that these new generators arise as part of the GL(32) automorphism group of the supersymmetry algebra in the fully supersymmetric theory. One puzzle with realising gauge fields as Goldstone fields is that one does not obviously find their field strengths when following the standard method of non-linear realisations. In fact, the field strengths of these gauge fields arise only as a result of demanding that the theory be invariant under both groups.
We will also show, in section three, that the bosonic sector of IIA supergravity can also be derived as a non-linear realisation. In section four, we show that if one starts with a generic theory of Goldstone fields, some of which carry anti-symmetrised space-time indices, and also demands conformal symmetry then one finds that these Goldstone fields must possess gauge symmetries. In section five, we show that the low energy effective action for the closed bosonic string can be written as a non-linear realisation. We explain, in section six, how one may derive the dynamics of branes in a background using the theory of non-linear realisations. We sketch, in section seven, how the non-linear realisation of the bosonic sectors of supergravity theories can be extended to the fully theory including the fermions. We conclude in section eight.
2. Eleven Dimensional Supergravity The Lagrangian of eleven dimensional supergravity written in the signature $`diag(\eta _{ab})=(1,1,1,\mathrm{}1)`$ is given by
$$L=+\frac{e}{4\kappa ^2}R\left(\mathrm{\Omega }(e,\psi )\right)\frac{e}{48}F_{\mu _1\mathrm{}\mu _4}F^{\mu _1\mathrm{}\mu _4}\frac{e}{2}\overline{\psi }_\mu \mathrm{\Gamma }^{\mu \nu \varrho }D_\nu \left(\frac{1}{2}(\widehat{\widehat{\mathrm{\Omega }}}+\widehat{\mathrm{\Omega }})\right)\psi _\varrho $$
$$\frac{1}{192}e\kappa (\overline{\psi }_{\mu _1}\mathrm{\Gamma }^{\mu _1\mathrm{}\mu _6}\psi _{\mu _2}+12\overline{\psi }^{\mu _3}\mathrm{\Gamma }^{\mu _4\mu _5}\psi ^{\mu _6})(F_{\mu _3\mathrm{}\mu _6}+\widehat{F}_{\mu _3\mathrm{}\mu _6})$$
$$+\frac{2\kappa }{(12)^4}ϵ^{\mu _1\mathrm{}\mu _{11}}F_{\mu _1\mathrm{}\mu _4}F_{\mu _5\mathrm{}\mu _8}A_{\mu _9\mu _{10}\mu _{11}}$$
$`(2.1)`$
where
$$F_{\mu _1\mathrm{}\mu _4}=4_{[\mu _1}A_{\mu _2\mu _3\mu _4]}$$
$$\widehat{F}_{\mu _1\mathrm{}\mu _4}=F_{\mu _1\mathrm{}\mu _4}+3\overline{\psi }_{[\mu _1}\mathrm{\Gamma }_{\mu _2\mu _3}\psi _{\mu _4]},$$
$`(2.2)`$
and
$$\widehat{\widehat{\mathrm{\Omega }}}_{\mu bc}=\widehat{\mathrm{\Omega }}_{\mu bc}+\frac{1}{4}\overline{\psi }_\nu \mathrm{\Gamma }_{\mu bc}^{\nu \lambda }\psi _\lambda ,$$
$$\widehat{\mathrm{\Omega }}_{\mu bc}=w_{\mu bc}(e)\frac{1}{2}(\overline{\psi }_\nu \mathrm{\Gamma }_c\psi _b\overline{\psi }_\nu \mathrm{\Gamma }_b\psi _c+\overline{\psi }_c\mathrm{\Gamma }_\mu \psi _b)$$
and
$$w_{\mu bc}(e)=\frac{1}{2}(e_b{}_{}{}^{\rho }_{\mu }^{}e_{\rho c}e_c{}_{}{}^{\rho }_{\mu }^{}e_{\rho b})\frac{1}{2}(e_b{}_{}{}^{\rho }_{\rho }^{}e_{\mu c}e_c{}_{}{}^{\rho }_{\rho }^{}e_{\mu b})\frac{1}{2}(e_b{}_{}{}^{\lambda }e_{c}^{}{}_{}{}^{\rho }_{\lambda }^{}e_{\rho a}e_c{}_{}{}^{\lambda }e_{b}^{}{}_{}{}^{\rho }_{\lambda }^{}e_{\rho a})e_\mu ^a$$
$`(2.3)`$
The symbol $`w_{\mu mn}(e)`$ is the usual expression for the spin-connection in terms of the vielbein $`e_\mu ^n`$.
The equation of motion of the gauge field is given by
$$D_{\mu _1}F^{\mu _1\mathrm{}\mu _4}+\frac{(dete)^1}{4.(12)^2}ϵ^{\mu _2\mathrm{}\mu _4\nu _1\mathrm{}\nu _8}F_{\nu _1\nu _2\nu _3\nu _4}F_{\nu _5\nu _6\nu _7\nu _8}=0$$
$`(2.4)`$
This equation may be rewritten in first order form as
$$F_{\mu _1\mathrm{}\mu _4}=\frac{dete}{7!}ϵ_{\mu _1\mathrm{}\mu _4\nu _1\mathrm{}\nu _7}F^{\nu _1\mathrm{}\nu _7}$$
$`(2.5)`$
where
$$\stackrel{~}{F}_{\nu _1\mathrm{}\nu _7}7(_{[\nu _1}A_{\nu _2\mathrm{}\nu _7]}+5A_{[\nu _1\mathrm{}\nu _3}F_{\nu _4\mathrm{}\mu _7]})$$
$`(2.6)`$
We consider the Lie algebra whose non-vanishing commutators are
$$[K^a{}_{b}{}^{},K^c{}_{d}{}^{}]=\delta _b^cK^a{}_{d}{}^{}\delta _d^aK^c{}_{b}{}^{},[K^a{}_{b}{}^{},P_c]=\delta _c^aP_b$$
$$[K^a{}_{b}{}^{},R^{c_1\mathrm{}c_6}]=\delta _b^{c_1}R^{ac_2\mathrm{}c_6}+\mathrm{},[K^a{}_{b}{}^{},R^{c_1\mathrm{}c_3}]=\delta _b^{c_1}R^{ac_2c_3}+\mathrm{},$$
$`(2.7)`$
$$[R^{c_1\mathrm{}c_3},R^{c_4\mathrm{}c_6}]=2R^{c_1\mathrm{}c_6}$$
$`(2.8)`$
where $`+\mathrm{}`$ denote the appropriate anti-symmetrisations. The generators $`K^a_b`$ and $`P_c`$ generate the affine group IGL(11) while the generators $`R^{c_1\mathrm{}c_3}`$ and $`R^{c_1\mathrm{}c_6}`$ form a subalgebra that is the same as that which was found to be a symmetry of the fivebrane . This subalgebra was also required in the description of the fivebrane as a non-linear realisation . In these references it was identified as part of the GL(32) automorphism algebra of the eleven dimensional supersymmetry algebra. We denote by $`G_{11}`$ the group whose Lie algebra is that of equations (2.7) and (2.8).
We now construct the non-linear realisation corresponding to the group $`G_{11}`$ taking the Lorentz group to be a local symmetry. The generators of the Lorentz group are given by $`J_{ab}=K_{ab}K_{ba}`$ where the indices are lowered with the Minkowski metric. We therefore consider group elements of the form
$$g=e^{x^\mu P_\mu }e^{h_a{}_{}{}^{b}K_{}^{a}_b}exp(\frac{A_{c_1\mathrm{}c_3}R^{c_1\mathrm{}c_3}}{3!}+\frac{A_{c_1\mathrm{}c_6}R^{c_1\mathrm{}c_6}}{6!})$$
$`(2.9)`$
The fields $`h^a_b`$, $`A_{c_1\mathrm{}c_3}`$ and $`A_{c_1\mathrm{}c_6}`$ depend on $`x^\mu `$. Although we use the exponential parameterization the reader who prefers a globally valid expression can readily rewrite the above group element in the appropriate form. In fact, one could take $`x^\mu `$, $`h^a_b`$, $`A_{c_1\mathrm{}c_3}`$ and $`A_{c_1\mathrm{}c_6}`$ to depend on D parameters, thus leading to a kind of democracy between fields and coordinates. To recover the above form from this formulation, one uses the reparameterisation invariance inherent in the construction to choose the $`x^\mu `$ equal to the parameters.
The theory is to be invariant under
$$gg_0gh^1$$
$`(2.10)`$
where $`g_0`$ is a rigid element of the full group generated by the above Lie algebra and $`h`$ is a local element of the Lorentz group. The corresponding $`g_0`$ invariant forms are given by
$$𝒱=g^1dgw$$
$`(2.11)`$
where $`w\frac{1}{2}dx^\mu w_{\mu b}{}_{}{}^{c}J_{}^{b}_c`$ is the Lorentz connection and so transforms as
$$whwh^1+hdh^1$$
$`(2.12)`$
As a result
$$𝒱h𝒱h^1$$
$`(2.13)`$
This approach differs from that of reference where the Lorentz symmetry was a rigid symmetry and the field $`h_a^b`$ was symmetric. The advantage of the approach adopted here is that one finds directly the vielbein formulation of general relativity and so the identification of part of the theory with general relativity is readily apparent.
Evaluating $`𝒱`$ we find that
$$𝒱=dx^\mu (e_\mu {}_{}{}^{a}P_{a}^{}+\mathrm{\Omega }_a{}_{}{}^{b}K_{}^{a}{}_{b}{}^{}+\frac{1}{3!}\stackrel{~}{D}_\mu A_{c_1\mathrm{}c_3}R^{c_1\mathrm{}c_3}+\frac{1}{6!}\stackrel{~}{D}_\mu A_{c_1\mathrm{}c_6}R^{c_1\mathrm{}c_6})$$
$`(2.14)`$
where
$$e_\mu {}_{}{}^{a}(e^h)_\mu {}_{}{}^{a},\stackrel{~}{D}_\mu A_{c_1\mathrm{}c_3}_\mu A_{c_1c_2c_3}+((e^1_\mu e)_{c_1}{}_{}{}^{b}A_{bc_2c_3}^{}+\mathrm{}),$$
$$\stackrel{~}{D}_\mu A_{c_1\mathrm{}c_6}_\mu A_{c_1\mathrm{}c_6}+((e^1_\mu e)_{c_1}{}_{}{}^{b}A_{bc_2\mathrm{}c_6}^{}+\mathrm{})(A_{[c_1\mathrm{}c_3}\stackrel{~}{D}_\mu A_{c_4\mathrm{}c_6]})$$
$$\mathrm{\Omega }_{\mu b}{}_{}{}^{c}(e^1_\mu e)_b{}_{}{}^{c}w_{\mu b}{}_{}{}^{c},$$
$`(2.15)`$
where $`+\mathrm{}`$ denotes the action of $`(e^1_\mu e)`$ on the other indices of $`A_{c_1\mathrm{}c_3}`$ and $`A_{c_1\mathrm{}c_6}`$.
The covariant derivatives of the Goldstone fields associated with this non-linear realisation, that is of the field $`h_a^b`$ and the fields $`A_{c_1\mathrm{}c_3}`$ and $`A_{c_1\mathrm{}c_6}`$, are given by
$$\mathrm{\Omega }_{ab}{}_{}{}^{c}(e^1)_a{}_{}{}^{\mu }(e^1_\mu e)_{b}^{}{}_{}{}^{c}w_{ab}{}_{}{}^{c},$$
$$\stackrel{~}{D}_aA_{c_1\mathrm{}c_3}(e^1)_a{}_{}{}^{\mu }\stackrel{~}{D}_{\mu }^{}A_{c_1\mathrm{}c_3},\stackrel{~}{D}_aA_{c_1\mathrm{}c_6}(e^1)_a{}_{}{}^{\mu }\stackrel{~}{D}_{\mu }^{}A_{c_1\mathrm{}c_6}$$
$`(2.16)`$
We note that these quantities are not field strengths as the indices are not anti-symmetrised, nor are the derivatives those that occur in general relativity. As we shall see, we will only recover the field strengths of the equations of motion once we consider the simultaneous non-linear realisation with the conformal group.
Under $`h=e^{(\frac{1}{2}w_a{}_{}{}^{b}(x)J^a{}_{b}{}^{})}`$, we find that $`e_\mu ^a`$ transforms as a elfbein should under a local Lorentz transformation and indeed all the indices of the fields in $`𝒱`$ which are contracted with the generators are rotated in the correct way as to be interpreted as tangent space indices. A matter field $`B`$ transforms as $`BB^{}=D(h)B`$ where $`D`$ is the representation of the Lorentz group to which $`B`$ belongs. The covariant derivative of the matter field $`B`$ is given by
$$\stackrel{~}{D}_aB(e^1)_a{}_{}{}^{\mu }_{\mu }^{}B+\frac{1}{2}w_{ab}^c\mathrm{\Sigma }^b{}_{c}{}^{}B,$$
$`(2.17)`$
where $`w_{ab}^c(e^1)_a{}_{}{}^{\mu }w_{\mu b}^{c}`$ and $`\mathrm{\Sigma }_a^b`$ is the representation of the generators of the Lorentz group associated with $`B`$.
As we have mentioned above, the bosonic sector of eleven dimensional supergravity is not uniquely determined by just taking a non-linear realisation of the group $`G_{11}`$. We must find a non-linear realisation of the closure of this group with the conformal group. We could calculate the closure and find the non-linear realisation of this infinite dimensional group. However, it is easier to find the simultaneous realisation of the two groups, that is calculate the Cartan forms for both groups and then use only those combinations that are invariant under both groups. We must also take into account any duplication of Goldstone fields between the two groups. We will now follow this second stratergy.
We now construct the non-linear realisation for the conformal group SO(2,11), taking the now rigid Lorentz group as the istropy group. This procedure is well known, , but for completeness we summarise the derivation. The generators of the conformal group obey the relations
$$[J_{ab},K_c]=\eta _{ac}K_b+\eta _{bc}K_a,[P_a,D]=P_a,[K_a,D]=K_a,[P_a,K_b]=+2\eta _{ab}D2J_{ab}$$
$`(2.18)`$
in addition to those of the Poincare group and relations where the commutators vanish. We take as our coset representative
$$g=e^{x^\mu P_\mu }e^{\varphi ^\mu K_\mu }e^{\sigma D}$$
$`(2.19)`$
and the Cartan forms are given by
$$g^1dg=dx^a(e^\sigma P_a+e^\sigma (_a\varphi ^b\varphi ^c\varphi _c\delta _a^b+2\varphi _a\varphi ^b)K_b+(_a\sigma +2\varphi _a)D+(\delta _a^c\varphi ^d+\delta _a^d\varphi ^c)J_{cd})$$
$`(2.20)`$
The covariant derivatives of the Goldstone fields are obtained by taking all the above expressions, with the exception of the first term, and multiplying by $`e^\sigma `$. These transform only under the Lorentz group and we can set the covariant derivative for $`\sigma `$ to zero and still preserve the group. As a result, we can eliminate $`\varphi _\mu `$ in terms of $`_\mu \sigma `$, namely $`2\varphi _\mu =_\mu \sigma `$. In effect, this leaves $`\sigma `$ as the only Goldstone field.
It is straightforward to find the transformations under dilations and special conformal transformations of $`\sigma `$ and a field $`B`$ which transforms under the representation $`\mathrm{\Sigma }_{ab}`$ of the Lorentz group. The result is
$$\delta \sigma =(2x\beta xx^2\beta )\sigma +2\beta _\mu x^\mu +\lambda ,\delta B=(2x\beta xx^2\beta )B+(\beta ^ax^b\beta ^bx^a)\mathrm{\Sigma }_{ab}B$$
$`(2.21)`$
No dilation term occurs in the above variation since dilations are not part of the isotropy group.
The covariant derivative with respect to conformal transformations, denoted $`\mathrm{\Delta }_a`$, of the field $`B`$ is given by
$$\mathrm{\Delta }_aB=e^\sigma (_a+^b\sigma \mathrm{\Sigma }_{ab})B$$
$`(2.22)`$
In particular, for a vector field $`A_a`$ we have
$$\mathrm{\Delta }_aA_b=e^\sigma (_aA_b+\eta _{ab}^c\sigma A_c_b\sigma A_a)$$
$`(2.23)`$
Following the procedure of Borisov and Ogievetski , we must now construct quantities from the derivatives of the Goldstone fields of the first group $`G_{11}`$ which are also covariant with respect to the conformal group. In view of the identical transformation of $`x^\mu `$ under dilations and the determinant part of GL(11) we must identify $`h_\mu {}_{}{}^{a}=\overline{h}_\mu ^a+\sigma \delta _\mu ^a`$ where $`\overline{h}_\mu {}_{}{}^{\mu }=0`$. While the field $`\sigma `$ identified in this way must transform in the relavent way determined by the conformal group, the fields $`\overline{h}_\mu ^a`$, $`A_{c_1\mathrm{}c_3}`$ and $`A_{c_1\mathrm{}c_6}`$ transform under conformal transformations as their indices suggest. It is simplest to first carry out this procedure for the fields $`A_{c_1\mathrm{}c_3}`$ and $`A_{c_1\mathrm{}c_6}`$ . The $`G_{11}`$ covariant derivative of $`A_{c_1\mathrm{}c_3}`$ can be rewritten as
$$\stackrel{~}{D}_aA_{c_1\mathrm{}c_3}=(\overline{e})_a{}_{}{}^{\mu }(\mathrm{\Delta }_\mu A_{c_1\mathrm{}c_3}+e^\sigma (\eta _{\mu c_1}^d\sigma A_{d\mathrm{}c_3}$$
$$+D_{c_1}\sigma A_{\mu \mathrm{}c_3}+_\mu \sigma A_{c_1\mathrm{}c_3}+(\overline{e}^1_\mu \overline{e})_{c_1}{}_{}{}^{d}A_{d\mathrm{}c_3}^{}+\mathrm{}))$$
$`(2.24)`$
In this equation $`\overline{e}=e^{\overline{h}}`$ and $`+\mathrm{}`$ denotes the terms that arise when the connections of the derivatives contract with the other indices on $`A_{c_1\mathrm{}c_3}`$. Even at order $`(\overline{h})^0`$ it is apparent that only by completely anti-symmetrising in $`a,c_1\mathrm{}c_3`$ can one obtain an expression such that all $`\sigma `$ dependence is through the conformal derivative $`\mathrm{\Delta }_a`$ alone and as a result is an expression that is simultaneously covariant under both groups. Thus the unique simultaneously covariant quantity is
$$\stackrel{~}{F}_{c_1\mathrm{}c_4}4(e_{[c_1}{}_{}{}^{\mu }_{\mu }^{}A_{c_2\mathrm{}c_4]}+e_{[c_1}{}_{}{}^{\mu }(e^1_\mu e)_{[c_2}^{}^bA_{bc_3c_4]}+\mathrm{})$$
$`(2.25)`$
A similar calculation for the gauge field $`A_{c_1\mathrm{}c_6}`$ leads to the unique simultaneously covariant expression
$$\stackrel{~}{F}_{c_1\mathrm{}c_7}7(e_{[c_1}{}_{}{}^{\mu }(_\mu A_{c_2\mathrm{}c_7]})+e_{[c_1}{}_{}{}^{\mu }(e^1_\mu e)_{[c_2}^{}^bA_{bc_3\mathrm{}c_7]}+\mathrm{}+5\stackrel{~}{F}_{[c_1\mathrm{}c_4}\stackrel{~}{F}_{c_5\mathrm{}c_7]})$$
$`(2.26)`$
What is not apparent from the above expressions is that the covariant derivatives are those that one should find in general relativity. To verify this, and indeed to recover general relativity itself, we must recover the usual expression for the spin connection in terms of the elfbein. We can use the inverse Higgs effect to place constraints on the covariant derivative $`\mathrm{\Omega }_{ab}^c`$ of the field $`h_\mu ^a`$. Within the context of the group $`G_{11}`$ there is no unique way to do this. However, as explained above, we must do this in just such a way that the $`G_{11}`$ covariant derivative of equation (2.17) can be rewritten in terms of the covariant derivatives of the conformal group of equation (2.21). Consider a matter field $`B`$ as discussed above, its $`G_{11}`$ covariant derivative can be expressed as
$$\stackrel{~}{D}_aB=(\overline{e})_a{}_{}{}^{\mu }(\mathrm{\Delta }_\mu Be^\sigma _\nu \sigma \mathrm{\Sigma }_\mu {}_{}{}^{\nu }B+\frac{1}{2}e^\sigma w_{\mu b}^c\mathrm{\Sigma }^b{}_{c}{}^{}B)$$
$`(2.27)`$
This equation tells us that $`w_{ab}^c`$ must be expressible as in terms of the conformal covariant derivatives of $`\overline{e}_\mu ^a`$ as well as a derivatives of $`\sigma `$ which must cancel the second term in the right-hand side of the above equation. The unique solution is to take the constraint
$$\mathrm{\Omega }_{a[bc]}\mathrm{\Omega }_{b(ac)}+\mathrm{\Omega }_{c(ab)}=0$$
$`(2.28)`$
This results in the well known expression for the spin connection in terms of the elfbein given in equation (2.3). Although the connection term that appears in the field strengths of equations (2.25) and (2.26) looks incorrect, when one takes into account the anti-symmetry on all the indices one finds that the covariant derivative can be re-expressed in terms of the usual spin connection appropriate to tangent indices. Thus the field strengths of these equations when written in terms of the components appropriate for the coordinates of the space-time are just the curl of the gauge potential written in the same components.
The equations of motion for the simultaneous non-linear realisation must be written in terms of $`\stackrel{~}{F}_{c_1\mathrm{}c_4}`$ and $`\stackrel{~}{F}_{c_1\mathrm{}c_7}`$ given in equations (2.25) and (2.26) and the spin-connection in such a way that the equations are covariant under the local Lorentz transformations. Clearly, the spin connection can only enter either in $`\stackrel{~}{F}_{c_1\mathrm{}c_4}`$ and $`\stackrel{~}{F}_{c_1\mathrm{}c_7}`$, in the way which is already specified, or through the Riemann tensor
$$R_{\mu \nu b}^c_\mu w_{\nu b}^c+w_{\mu b}^dw_{\nu d}^c(\mu \nu )$$
$`(2.29)`$
The unique first order equation for the gauge field which is not trivial is
$$\stackrel{~}{F}_{c_1\mathrm{}c_4}=\frac{1}{7!}ϵ_{c_1\mathrm{}c_{11}}\stackrel{~}{F}^{c_5\mathrm{}c_{11}}$$
$`(2.30)`$
in agreement with equation (2.5) when written in the local coordinate frame. The only other non-trivial equation is
$$R_{\mu \nu b}^ce_c{}_{}{}^{\nu }e_{a}^{}{}_{}{}^{\mu }\frac{1}{2}\eta _{ab}R_{\mu \nu b}^ce_c{}_{}{}^{\nu }e_{}^{b\mu }\frac{c}{4}(\stackrel{~}{F}_{ac_1\mathrm{}c_3}\stackrel{~}{F}_b{}_{}{}^{c_1\mathrm{}c_3}\frac{1}{6}\eta _{ab}\stackrel{~}{F}_{c_1\mathrm{}c_4}\stackrel{~}{F}^{c_1\mathrm{}c_4})=0$$
$`(2.31)`$
as it should be. The constant $`c`$, of proportionality can only be determined once the full supersymmetric treatment is given. It has value 1. 3. Gauge Symmetry It is instructive to trace more precisely how the gauge invariance of the fields $`A_{a_1\mathrm{}a_3}`$ and $`A_{a_1\mathrm{}a_6}`$ arises as a consequence of the simultaneous realisation of $`G_{11}`$ and the conformal group. Taking
$$g_0=exp(\frac{c_{\mu _1\mathrm{}\mu _3}\delta _{a_1\mathrm{}a_3}^{\mu _1\mathrm{}\mu _3}R^{a_1\mathrm{}a_3}}{3!}+\frac{c_{\mu _1\mathrm{}\mu _6}\delta _{a_1\mathrm{}a_6}^{\mu _1\mathrm{}\mu _6}R^{a_1\mathrm{}a_6}}{6!})$$
$`(3.1)`$
where $`\delta _{a_1\mathrm{}a_n}^{\mu _1\mathrm{}\mu _n}=\delta _{\mu _1}^{a_1}\mathrm{}\delta _{\mu _n}^{a_n}`$ and $`c_{\mu _1\mathrm{}\mu _3}`$ and $`c_{\mu _1\mathrm{}\mu _6}`$ are constants in equation (2.10), we find that the vielbein is inert and the other fields transform as
$$\delta A_{a_1\mathrm{}a_3}=c_{a_1\mathrm{}a_3},\delta A_{a_1\mathrm{}a_6}=c_{a_1\mathrm{}a_6}+20c_{[a_1\mathrm{}a_3}A_{a_4\mathrm{}a_6]}$$
$`(3.2)`$
where $`c_{\mu _1\mathrm{}\mu _3}=e_{\mu _1}{}_{}{}^{a_1}\mathrm{}e_{\mu _3}{}_{}{}^{a_3}c_{a_1\mathrm{}a_3}^{}`$ and similarly for $`c_{a_1\mathrm{}a_6}`$. The vielbeins occur because the factor in $`g`$ which contains the fields $`A_{a_1\mathrm{}a_3}`$ and $`A_{a_1\mathrm{}a_6}`$ is to the right of that containing the $`h_a^b`$ fields. Thus, it is the fields with curved indices that transform most simply. To find the conformal transformation of the fields with curved indices we write them as $`A_{\mu _1\mathrm{}\mu _p}=(e^{\overline{h}})_{\mu _1}{}_{}{}^{a_1}\mathrm{}(e^{\overline{h}})_{\mu _p}{}_{}{}^{a_p}e_{}^{p\sigma }A_{a_1\mathrm{}a_p}`$ for $`p=3,6`$. Taking into account the conformal transformation of $`\sigma `$ and $`B`$ of equation (2.21) we find that
$$\delta A_{\mu _1\mathrm{}\mu _p}=(2(x\beta )(x)x^2(\beta ))A_{\mu _1\mathrm{}\mu _p}+(2\beta _{\mu _1}x^\kappa A_{\kappa \mu _2\mathrm{}\mu _p}2x_{\mu _1}\beta ^\kappa A_{\kappa \mu _2\mathrm{}\mu _p}+\mathrm{})$$
$$+2p(\beta x)A_{\mu _1\mathrm{}\mu _p}$$
$`(3.3)`$
where $`+\mathrm{}`$ stands for the other terms where the induced Lorentz rotation acts on the other indices of $`A_{\mu _1\mathrm{}\mu _p}`$. We note that this is the variation of a matter field that we would have obtain had we included the dilations in the isotropy group and assigned the field dilation weight p.
For simplicity, we will illustrate the mechanism of how gauge symmetry arises for a single form field $`A_{\mu _1\mathrm{}\mu _p}`$ that has a constant shift, i.e. $`\delta A_{\mu _1\mathrm{}\mu _p}=c_{\mu _1\mathrm{}\mu _p}`$, under a non-linear realisation. Carrying out the commutation of this shift with a special conformal transformation we find that
$$[\delta _c,\delta _\beta ]A_{\mu _1\mathrm{}\mu _p}=p_{[\mu _1}\stackrel{~}{\mathrm{\Lambda }}_{\mu _2\mathrm{}\mu _p]}^{(2)}$$
$`(3.4)`$
We recognise this as a gauge transformation with parameter
$$\stackrel{~}{\mathrm{\Lambda }}_{\mu _2\mathrm{}\mu _p}^{(2)}=2px\beta x^\kappa c_{\kappa \mu _2\mathrm{}\mu _p}x^2\beta ^\kappa c_{\kappa \mu _2\mathrm{}\mu _p}+(2x_{\mu _2}\beta ^\kappa x^\rho c_{\rho \kappa \mu _3\mathrm{}\mu _p}+\mathrm{})$$
$`(3.5)`$
We may write the original shift of the field as a gauge transformation with parameter $`\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _p}^{(1)}=x^\kappa c_{\kappa \mu _2\mathrm{}\mu _p}`$ and taking its commutator with special conformal transformations we find another gauge transformation which is quadratic in $`x^\mu `$. Thus starting from a gauge transformation that is linear in $`x^\mu `$ we obtain one which is bilinear.
By induction, we will now show that taking repeated commutators with special conformal transformations leads to a gauge transformation with an arbitrary local parameter. Let us suppose that after taking $`r1`$ commutators we find a transformation which can be written as a gauge transformation of $`A_{\mu _1\mathrm{}\mu _p}`$, denoted $`\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _p}^{(r)}`$, which is a polynomial in $`x^\mu `$ of degree r. Taking the commutator of this with another special conformal transformation we find that
$$[\delta _\mathrm{\Lambda },\delta _\beta ]A_{\mu _1\mathrm{}\mu _p}=p_{[\mu _1}\stackrel{~}{\mathrm{\Lambda }}_{\mu _2\mathrm{}\mu _p]}^{(r+1)}$$
$`(3.6)`$
where
$$\stackrel{~}{\mathrm{\Lambda }}_{\mu _2\mathrm{}\mu _p]}^{(r+1)}=(2x\beta xx^2\beta )\mathrm{\Lambda }_{\mu _2\mathrm{}\mu _p}^{(r)}$$
$$+(2\beta _{\mu _2}x^\kappa \mathrm{\Lambda }_{\kappa \mu _3\mathrm{}\mu _p}^{(r)}2x_{\mu _2}\beta ^\kappa \mathrm{\Lambda }_{\kappa \mu _3\mathrm{}\mu _p}^{(r)}+\mathrm{})+2(p1)x\beta \mathrm{\Lambda }_{\mu _2\mathrm{}\mu _p}^{(r)}$$
$`(3.7)`$
Hence we recover another gauge transformation, which is a polynomial in $`x^\mu `$ of one degree higher. It is clear that proceeding in this way we can find an arbitrary local gauge transformation. Thus, we have shown that taking the closure of a Goldstone shift and the conformal group leads to a local gauge transformation. In fact, if we start from a non-linear realisation of the fields we can regard gauge invariance as a consequence of conformal invariance.
The standard $`U(1)`$ field has been been previously considered as a Goldstone boson by considering an infinite dimensional algebra . 4. IIA Supergravity The bosonic part of the ten-dimensional IIA supergravity theory is given by ,,
$$\begin{array}{cc}\hfill L^B& =eR\left(w(e)\right)\frac{1}{12}ee^{\frac{\sigma }{2}}F_{\mu _1\mathrm{}\mu _4}^{}F^{\mu _1\mathrm{}\mu _4}\frac{1}{3}ee^\sigma F_{\mu _1\mathrm{}\mu _3}F^{\mu _1\mathrm{}\mu _3}\hfill \\ & ee^{\frac{3}{2}\sigma }F_{\mu _1\mu _2}F^{\mu _1\mu _2}\frac{1}{2}_\mu \sigma ^\mu \sigma \hfill \\ & +\frac{1}{2(12)^2}ϵ^{\mu _1\mathrm{}\mu _{10}}F_{\mu _1\mathrm{}\mu _4}F_{\mu _5\mathrm{}\mu _8}A_{\mu _9\mu _{10}}\hfill \end{array}$$
$`(4.1)`$
where
$$F_{\mu _1\mu _2}=2_{[\mu _1}A_{\mu _2]}$$
$`(4.2)`$
$$F_{\mu _1\mu _2\mu _3}=3_{[\mu _1}A_{\mu _2\mu _3]}$$
$`(4.3)`$
$$F_{\mu _1\mathrm{}\mu _4}^{}=4(_{[\mu _1}A_{\mathrm{}\mu _4]}+2A_{[\mu _1}F_{\mu _2\mu _3\mu _4]})$$
$`(4.4)`$
The equations of motion of the bosonic, non-gravitational sector of the theory were expressed as a non-linear realisation in reference . Although we use a different group and strategy, some of the steps below have analogues in those of reference . We proceed much as for eleven-dimensional supergravity, we take the group $`G_{IIA}`$ to have the generators of IGL(10) together with the generators $`R_{a_1\mathrm{}a_p}`$ for $`p=0,1,2,3,5,6,7,8`$ that obey equations analoguous to equation (2.7 ) as well as the relations
$$[R,R^{a_1\mathrm{}a_p}]=c_pR^{a_1\mathrm{}a_p},[R^{a_1\mathrm{}a_p},R^{a_1\mathrm{}a_q}]=c_{p,q}R^{a_1\mathrm{}a_{(p+q)}}$$
$`(4.5)`$
where
$$c_1=c_7=\frac{3}{4},c_2=c_6=\frac{1}{2},c_3=c_5=\frac{1}{4}:$$
$$c_{1,2}=c_{2,3}=c_{3,3}=c_{2,5}=c_{1,5}=2,c_{1,7}=3,c_{2,6}=2,c_{3,5}=1$$
$`(4.6)`$
and all other $`c`$’s vanish. We take the group element of $`G_{IIA}`$ to be of the form
$$g=e^{x^\mu P_\mu }g_hg_A$$
$`(4.6)`$
where
$$g_h=e^{h_a{}_{}{}^{b}K_{}^{a}_b},$$
and
$$g_A=exp(\frac{A_{a_1\mathrm{}a_8}}{8!}R^{a_1\mathrm{}a_8})$$
$$\mathrm{}exp(\frac{A_{a_1\mathrm{}a_3}}{3!}R^{a_1a_2a_3})exp(\frac{A_{a_1a_2}}{2!}R^{a_1a_2})exp(A_aR^a)exp(AR)$$
$`(4.7)`$
We also take the Lorentz group to be a local symmetry and so consider the quantity
$$𝒱=g^1dgw$$
$`(4.8)`$
We may rewrite this as
$$𝒱=(g_h^1dg_h)+(g_A^1dg_A+g_A^1(g_h^1dg_h)g_Ag_h^1dg_h)$$
$$dx^\mu (e_\mu {}_{}{}^{a}P_{a}^{}+\mathrm{\Omega }_{\mu a}^bK^a{}_{b}{}^{})+dx^\mu (\underset{p=1}{\overset{8}{}}\frac{1}{p!}\stackrel{~}{D}_\mu A_{a_1\mathrm{}a_p}R^{a_1\mathrm{}a_p})$$
$`(4.9)`$
where the definition applies to each of the terms in the brackets separately.
The next step is to demand conformal invariance. In particular, we should take combinations of derivatives of the Goldstone fields that are also conformally covariant. The procedure follows closely those of sections two and three. Indeed, it is an inevitable consequence of section three that the fields $`A_{a_1\mathrm{}a_p}`$ will appear in quantities that are gauge invariant. Thus the quantities which involve the fields $`A_{a_1\mathrm{}a_p}`$ that are $`G_{IIA}`$ and conformally covariant are
$$\stackrel{~}{F}_{a_1\mathrm{}a_p}pe^{c_{(p1)}A}\stackrel{~}{D}_{[a_1}A_{a_2\mathrm{}a_3]}$$
$`(4.10)`$
One finds that
$$\stackrel{~}{F}_a=D_aA,\stackrel{~}{F}_{a_1a_2}=2e^{\frac{3}{4}A}D_{[a_1}A_{a_2]},\stackrel{~}{F}_{a_1a_2a_3}=3e^{\frac{1}{2}A}D_{[a_1}A_{a_2a_3]},$$
$$\stackrel{~}{F}_{a_1a_2a_3a_4}=4e^{\frac{1}{4}A}(D_{[a_1}A_{a_2a_3a_4]}+2e^{\frac{1}{2}A}A_{[a_1}\stackrel{~}{F}_{a_2a_3a_4]}),\stackrel{~}{F}_{a_1a_2a_3a_4a_5}=0,$$
$$\stackrel{~}{F}_{a_1\mathrm{}a_6}=6e^{\frac{1}{4}A}(D_{[a_1}A_{a_2\mathrm{}a_6]}+5e^{\frac{1}{4}A}A_{[a_1a_2}\stackrel{~}{F}_{a_3\mathrm{}a_6]}),$$
$$\stackrel{~}{F}_{a_1\mathrm{}a_7}=7e^{\frac{1}{2}A}(D_{[a_1}A_{a_2\mathrm{}a_7]}20A_{[a_1a_2a_3}D_{[a_4}A_{a_5\mathrm{}a_7]}+2e^{\frac{1}{4}A}A_{[a_1}\stackrel{~}{F}_{a_2\mathrm{}a_7]}),$$
$$\stackrel{~}{F}_{a_1\mathrm{}a_8}=8e^{\frac{3}{4}A}(D_{[a_1}A_{a_2\mathrm{}a_8]}7.6A_{[a_1a_2}(D_{a_3}A_{a_4\mathrm{}a_8]}+10A_{[a_3a_4}(D_{[a_5}A_{a_6\mathrm{}a_8]})),$$
$$\stackrel{~}{F}_{a_1\mathrm{}a_9}=9(D_{[a_1}A_{a_2\mathrm{}a_9]}8.7A_{[a_1a_2}D_{[a_3}A_{a_4\mathrm{}a_9]}+\frac{7.4}{3}e^{\frac{1}{4}A}A_{[a_1a_2a_3}\stackrel{~}{F}_{a_4\mathrm{}a_9]}+3e^{\frac{3}{4}A}A_{[a_1}\stackrel{~}{F}_{a_2\mathrm{}a_9]})$$
$`(4.11)`$
The unique equations of motion of the forms that are first order in derivatives and constructed from the simultaneously covariant derivatives of the Goldstone fields are
$$\stackrel{~}{F}^{a_1\mathrm{}a_p}=\frac{1}{(10p)!}ϵ^{a_1\mathrm{}a_{10}}\stackrel{~}{F}_{a_{(p+1)}\mathrm{}a_{10}},p=1,2,3,4$$
$`(4.12)`$
as well as the equation for the vielbein. These are the equations of motion of IIA supergravity.
One might have thought that the easiest way to obtain the non-linear realisation of IIA supergravity would be to directly perform the dimensional reduction on the non-linear realisation of eleven dimensional supergravity. However, the correct number of fields and generators does not arise in a natural way as becomes apparent if one tries. This suggests that the formulation of eleven dimensional supergravity given in section two may not be the most natural one and that there should exist a first order formulation of the vielbein equation of motion by introducing higher rank fields. 5. The Closed Bosonic String Effective Action One can also apply the theory of non-linear realisations to the low energy effective action of the closed string. This has been found to be
$$d^Dx𝑑ete(R\frac{4}{(D2)}_\mu \varphi ^\mu \varphi \frac{1}{3}e^{\frac{8\varphi }{D2}}F_{\mu _1\mu _2\mu _3}F^{\mu _1\mu _2\mu _3})$$
$`(5.1)`$
where $`F_{\mu _1\mu _2\mu _3}=3_{[\mu _1}A_{\mu _2\mu _3]}`$. In this action, $`D`$ is the dimension of space-time, but we must take $`D=26`$ to obtain the consistent closed bosonic string. In principle we should include the cosmological term for $`D26`$ and there will also be corrections to the dilaton potential from higher order genus surfaces . Presumeably, these terms could be accounted for by taking into account an anomaly in the symmetry below.
We then consider the group $`G_D`$ whose generators are $`K^a_b`$, $`R`$, $`R^{a_1a_2}`$, $`R^{a_1\mathrm{}a_{(D4)}}`$ and $`R^{a_1\mathrm{}a_{(D4)}}`$. They obey equations analogous to equation (2.7)as well as the relations
$$[R,R^{a_1\mathrm{}a_p}]=c_pR^{a_1\mathrm{}a_p},[R^{a_1\mathrm{}a_p},R^{a_1\mathrm{}a_q}]=c_{p,q}R^{a_1\mathrm{}a_{(p+q)}}$$
$`(5.2)`$
where
$$c_2=c_{D4}=\frac{4}{(D2)},c_{2,D4}=2.$$
$`(5.3)`$
We take and all the other $`c`$’s vanish.
The non-linear realisation of $`G_D`$ is built out of the group element $`g=g_hg_A`$ where
$$g_A=exp(\frac{A_{a_1\mathrm{}a_{(D2)}}R^{a_1\mathrm{}a_{(D2)}}}{(D2)!})exp(\frac{A_{a_1\mathrm{}a_{(D4)}}R^{a_1\mathrm{}a_{(D4)}}}{(D4)!})exp(\frac{A_{a_1a_2}R^{a_1a_2}}{(2)!})exp(AR)$$
$`(5.4)`$
Calculating the Cartan forms $`g^1dgw`$ and demanding simultaneous invariance under the conformal group, we find that the equations of motion must be built out of $`w_{ab}^c`$ and
$$\stackrel{~}{F}_a=D_aA,\stackrel{~}{F}_{a_1a_2a_3}=3e^{\frac{4}{(D2)}A}D_{[a_1}A_{a_2a_3]},$$
$$\stackrel{~}{F}_{a_1\mathrm{}a_{(D3)}}=(D3)e^{\frac{4}{(D2)}A}D_{[a_1}A_{a_2\mathrm{}a_{(D3)}]},$$
$$\stackrel{~}{F}_{a_1\mathrm{}a_{(D1)}}=(D1)(D_{[a_1}A_{a_2\mathrm{}a_{(D1)}]}+(D2)e^{\frac{4}{(D2)}A}A_{[a_1a_2}F_{a_3\mathrm{}a_{(D1)]}})$$
$`(5.5)`$
The equations of motion are given by
$$\stackrel{~}{F}^{a_1\mathrm{}a_p}=\frac{1}{(Dp)!}ϵ^{a_1\mathrm{}a_D}\stackrel{~}{F}_{a_{(p+1)}\mathrm{}a_D},p=1,2$$
$`(5.6)`$
provided we identify $`\varphi `$ with $`A`$, as well as the vielbein equation. 6. Branes in a Background In a recent paper , the branes of M theory were derived as a non-linear realisation. Since in this paper we have shown that the background supergravity to which they couple can also be formulated as a non-linear realisation, it is straightforward, at least in principle, to describe the dynamics of branes in a background as a non-linear realisation. We now illustrate the procedure for the case of a bosonic brane coupled to gravity.
We begin with a group element of IGL(D) of the form
$$g=e^{X^{\underset{¯}{a}}(\xi )P_{\underset{¯}{a}}}e^{h_{\underset{¯}{a}}{}_{}{}^{\underset{¯}{b}}(X)K^{\underset{¯}{a}}_{\underset{¯}{b}}}$$
$`(6.1)`$
We use the same index notation as in reference , where $`D`$ is the dimension of space-time and $`\xi ^n`$ are the coordinates of the brane worldvolume. We consider the forms
$$𝒱=g^1dgwd\xi ^n(e_n{}_{}{}^{a}P_{a}^{}+f_n{}_{}{}^{a^{}}P_{a^{}}^{}+\omega _{n\underset{¯}{a}}^{\underset{¯}{b}}K^{\underset{¯}{a}}{}_{\underset{¯}{b}}{}^{})$$
$`(6.2)`$
The spin connection $`w`$ takes values in the Lie algebra of $`SO(1,p)\times SO(Dp1)`$.
In fact, in reference , we took the isotropy group to be SO(1,p), although we could have taken the above group. Making this latter choice simplifies the analysis of reference a bit, but the results are the same.
Returning to the bosonic brane in a background, we find that
$$e_n{}_{}{}^{a}=_nX^{\underset{¯}{m}}e_{\underset{¯}{m}}{}_{}{}^{a},f_n{}_{}{}^{a^{}}=_nX^{\underset{¯}{m}}e_{\underset{¯}{m}}^a^{}$$
$`(6.3)`$
Since $`f_n^a^{}`$ transforms in a covariant manner, we can set it to zero. This solves for $`_nX^a^{}`$ in terms of the $`e_n^a^{}`$ of the background gravity. Hence, in the case of a local background we find that the parts of the vielbein that belong to the coset $`\frac{SO(1,D1)}{SO(1,p)\times SO(Dp1)}`$ play the role of the Goldstone bosons of the Lorentz group that were solved for in reference . Proceeding as in that paper, and using the constraint $`f_n{}_{}{}^{a^{}}=0`$, we find that
$$e_n{}_{}{}^{a}\eta _{ab}^{}e_m{}_{}{}^{b}=_nX^{\underset{¯}{p}}g_{\underset{¯}{p}\underset{¯}{q}}_mX^{\underset{¯}{q}}$$
$`(6.4)`$
where $`g_{\underset{¯}{n}\underset{¯}{m}}=e_{\underset{¯}{m}}{}_{}{}^{\underset{¯}{a}}\eta _{\underset{¯}{a}\underset{¯}{b}}^{}e_{\underset{¯}{n}}^{\underset{¯}{b}}`$. Thus, we find that the invariant action is
$$d^p\xi 𝑑ete_n^a=d^p\xi \sqrt{det(_nX^{\underset{¯}{p}}g_{\underset{¯}{p}\underset{¯}{q}}_mX^{\underset{¯}{q}})}$$
$`(6.5)`$
as it should be. 7. Supersymmetric Extension It would be interesting to extend the analysis of this paper to the full supergravity theories, that is incorporate supersymmetry. Let us first sketch how this would go for eleven dimensional supergravity. To extend the group IGL(11), it is natural to consider the group IGL(11/32). The generators of GL(11/32) group can be labelled by $`K^A_B`$ where $`A=(a,\alpha )`$ and similarly for $`B`$ etc. We can then denote the generators by $`K^A{}_{B}{}^{}=(K^a{}_{b}{}^{},K^a{}_{\alpha }{}^{},K^\alpha {}_{a}{}^{},K^\alpha {}_{\beta }{}^{})`$ and the generators of inhomogeneous transformations by $`P_a,K_\alpha `$. The non-linear realisation is then built from the group elements of the form
$$g=e^{(X^aP_a+K^\alpha \theta _\alpha )}e^{(h_a{}_{}{}^{b}K_{}^{a}{}_{b}{}^{}+A_\alpha {}_{}{}^{\beta }K_{}^{\alpha }{}_{\beta }{}^{})}e^{(\psi _a{}_{}{}^{\alpha }K_{}^{a}{}_{\alpha }{}^{}+\zeta _\alpha {}_{}{}^{a}K_{}^{\alpha }{}_{a}{}^{})}$$
$`(7.1)`$
In the group element of equation (7.1), the fields $`h_a^b`$, $`A_\alpha ^\beta `$, $`\psi _a^\alpha `$ and $`\zeta _\alpha ^a`$ are functions of $`x^a`$ and $`\theta _\alpha `$ and have geometric dimensions 0, 0, 1/2 and -1/2 respectively. Thus the lowest components of $`h_a^b`$, $`A_\alpha ^\beta `$ and $`\psi _a^\alpha `$ have the correct dimensions to be identified with the graviton, gauge fields and gravitino respectively. Indeed, their shifts under the appropriate symmetries of the non-linear realisations make this identification inevitable.
We must also consider a non-linear realisation of a supersymmetric generalisation of the conformal group. It is known , that there is a unique generalisation of SO(2,11) that also contains the supersymmetry algebra: it is $`Osp(1/64)`$. The Lie algebra of this group can be written as
$$[R_{\widehat{\alpha }\widehat{\beta }},R_{\widehat{\gamma }\widehat{\delta }}]=C_{\widehat{\beta }\widehat{\gamma }}R_{\widehat{\alpha }\widehat{\delta }}C_{\widehat{\alpha }\widehat{\delta }}R_{\widehat{\gamma }\widehat{\beta }}C_{\widehat{\beta }\widehat{\delta }}R_{\widehat{\alpha }\widehat{\gamma }}C_{\widehat{\alpha }\widehat{\gamma }}R_{\widehat{\beta }\widehat{\delta }}$$
$$\{\rho _{\widehat{\alpha }},\rho _{\widehat{\beta }}\}=R_{\widehat{\alpha }\widehat{\beta }},[\rho _{\widehat{\gamma }},R_{\widehat{\alpha }\widehat{\beta }}]=C_{\widehat{\beta }\widehat{\gamma }}\rho _{\widehat{\alpha }}+C_{\widehat{\alpha }\widehat{\gamma }}\rho _{\widehat{\beta }}$$
$`(7.2)`$
where $`C_{\widehat{\beta }\widehat{\gamma }}=C_{\widehat{\gamma }\widehat{\beta }}`$ is the metric that occurs in the invariant line element of this group and $`\widehat{\alpha },\widehat{\beta }=1,2\mathrm{},64`$.
We now decompose the 64 component spinor, $`\rho _{\widehat{\alpha }}`$ in this group into two 32 component spinors using the index decomposition $`\widehat{\alpha }=(\alpha ,\alpha ^{})`$ where $`\alpha =1,\mathrm{},32,\alpha ^{}=1,\mathrm{},32`$ etc. In particular, we set $`Q_\alpha =\rho _\alpha `$ and $`S^\alpha =\rho _\beta ^{}C^{\beta ^{}\alpha }`$, $`R_{\alpha \beta }=Z_{\alpha \beta }`$, $`R_\alpha {}_{}{}^{\beta }=R_{\alpha \beta ^{}}C^{\beta ^{}\beta }`$ and $`Z^{\alpha \beta }=R_{\alpha ^{}\beta ^{}}C^{\alpha ^{}\alpha }C^{\beta ^{}\beta }`$. Taking $`C_{\alpha \beta }=0=C_{\alpha ^{}\beta ^{}}`$ we may write the algebra of Osp(1/64) in the form
$$\{Q_\alpha ,Q_\beta \}=Z_{\alpha \beta },[Q_\alpha ,Z_{\gamma \delta }]=0,[Z_{\alpha \delta },Z_{\gamma \beta }]=0,$$
$$[Q_\alpha ,R_\gamma {}_{}{}^{\delta }]=\delta _\alpha ^\delta Q_\gamma ,[Z_{\alpha \beta },R_\gamma {}_{}{}^{\delta }]=\delta _\alpha ^\delta Z_{\gamma \beta }\delta _\beta ^\delta Z_{\gamma \alpha }$$
$`(7.3)`$
and
$$\{S^\alpha ,S^\beta \}=Z^{\alpha \beta },[S^\alpha ,Z^{\gamma \delta }]=0,[Z^{\alpha \delta },Z^{\gamma \beta }]=0,$$
$$[S^\gamma ,R_\alpha {}_{}{}^{\beta }]=\delta _\alpha ^\gamma S^\beta ,[Z^{\alpha \beta },R^\gamma {}_{\delta }{}^{}]=\delta _\gamma ^\beta Z^{\alpha \delta }+\delta _\gamma ^\alpha Z^{\delta \beta }$$
$`(7.4)`$
as well as
$$\{Q_\alpha ,S^\beta \}=R_\alpha {}_{}{}^{\beta },[Z_{\alpha \beta },Z^{\gamma \delta }]=\delta _\beta ^\gamma R_\alpha {}_{}{}^{\delta }\delta _\alpha ^\delta R_\beta {}_{}{}^{\gamma }\delta _\beta ^\delta R_\alpha {}_{}{}^{\gamma }\delta _\alpha ^\gamma R_\beta ^\delta $$
$`(7.5)`$
We recognise that Osp(1/64) contains a sub-algebra, given in equation (7.3), which is precisely the usual supersymmetry algebra in eleven dimensions with all its central charges, plus the GL(32) automorphism group that was found to play a role in the fivebrane equations of motion and in the branes of M theory realised as a non-linear realisation . Expanding in $`\gamma `$-matrices, we can express
$$R^\alpha {}_{\beta }{}^{}=\underset{n}{}R^{a_1\mathrm{}a_n}(\gamma _{a_1\mathrm{}a_n})^\alpha {}_{\beta }{}^{}.$$
$`(7.6)`$
The generators $`R`$ and $`R_{a_1a_2}`$ are to be identified with dilations and Lorentz rotations.
The considerations of this paper show that Osp(1/64) must be a symmetry of eleven dimensional supergravity. This group has previously been considered in the context of M theory with two times and mentioned as a possible unifying group in reference .
When taking the simultaneous realisation of the two groups IGL(11/32) and Osp(1/32) we must identify the Goldstone fields whose corresponding generators have the same action on the coordinates $`x^a`$ and $`\theta ^\alpha `$. In principle, one should also consider the action of Osp(1/64) on the central charges, but for the present discussion we shall ignore this subtilty. As for the bosonic sector consider in section two, the dilations are in common and so we must identify $`h^\mu _\mu `$ with $`\sigma `$, However, the generators $`K^\alpha _\beta `$ and $`R_\alpha ^\beta `$ act the same way on the coordinates $`x^a`$ and $`\theta ^\alpha `$ with the exception of the scalar and rank two generators that behave differently. In Osp(1/64) these are the dilations and Lorentz rotations and so their actions on the coordinates $`x^a`$ and $`\theta ^\alpha `$ are related. In contrast, the GL(11/32) action on the coordinates $`x^a`$ and $`\theta ^\alpha `$ is unrelated. Thus, we should identify all of the generators in $`K^\alpha _\beta `$ with those in $`R_\alpha ^\beta `$ with the exception of the two generators of rank zero and two. As a result, in the simultaneous non-linear realisation we find that we have generators of every rank in $`R_\alpha ^\beta `$, including the dilations and Lorentz rotations, as well as two additional generators of rank zero and two, which we may denote by $`K`$ and $`K_{ab}`$. Hence, in eleven dimensions, we find the automorphisms $`R^{a_1\mathrm{}a_n}`$ for $`n=0,1,2,3,4,6`$ and the two additional generators, $`K`$ and $`K_{a_1a_2}`$. It was observed on reference that the algebra of equation (2.8) was a contraction of the GL(32) automorphism algebra. Hence, it would seem natural to identify the generators $`R^{a_1\mathrm{}a_3}`$ and $`R^{a_1\mathrm{}a_6}`$, which are the generators, whose Goldstone fields are the gauge fields of eleven dimensional supergravity in the non-linear realisation of section two, with the automorphisms that arise in the groups GL(11/32) and Osp(1/64).
The supergravity action of equation (2.1) essentially contains three contributions, the kinetic terms in the first line, the Noether term in the second line and the Chern-Simmons term in the last line. We have already accounted for the first and last terms and in the supersymmetric extension we must account for the Noether term. In the Cartan forms of the group element of equation (7.1) we find a term of the form
$$e^{\psi _a{}_{}{}^{\alpha }K_{}^{a}_\alpha }(\stackrel{~}{D}_aA_{a_1\mathrm{}a_3}R^{a_1\mathrm{}a_3}+\stackrel{~}{D}_aA_{a_1\mathrm{}a_6}R^{a_1\mathrm{}a_6})e^{\psi _a{}_{}{}^{\alpha }K_{}^{a}_\alpha }$$
$`(7.7)`$
Taking the commutator of the $`R^{a_1\mathrm{}a_3}`$ and $`R^{a_1\mathrm{}a_6}`$ generators with $`K^a_\alpha `$ to be the obvious $`\gamma `$-matrix times $`K^a_\alpha `$ we do indeed find a term that has, at least in form, that of the Noether term.
We now briefly comment on the supersymmetric extension of the IIA theory considered in section four. The extension of the group IGL(10) is presumably the group IGL(10/32). Although reference was concerned with $`N=1`$ supersymmetry, it would seem inevitable that the unique extension of the conformal group to include a type II superalgebra is the group Osp(1/64). We should consider the simultaneous realisation of both of these groups. In the later group, we will find the GL(32) automorphism group and so the generators $`R^{a_1\mathrm{}a_p}`$ for $`p=0,1\mathrm{}10`$. Identifying the generators in the same way as above we find that the simultaneous non-linear realisation of these two groups includes the generators $`R^{a_1\mathrm{}a_n}`$ for $`n=0,1,2,\mathrm{},10`$ as well as two additional generators, which we can denote by $`K`$ and $`K_{a_1a_2}`$. We note that in contrast to the eleven dimensional theory most of the generators are need to ensure the necessary Goldstone bosons. The correspondence with the automorphisms of the supersymmetry algebra is less obvious in this case and it is possible that one may have to introduce generators in addition to those of the above two groups. In particular, the generator $`R`$, which leads to the SO(1,1) transformations of the IIA supergravity theory, does not seem to have an obvious identification with these generators.
The non-linear realisation of the IIb supergravity theory follows a similar pattern, but the automorphisms that are active are different from those in the eleven dimensional and IIA supergravity theories.
The above is a sketch of the extension to the supergravity theory, however, until one actually carries out the full calculation one cannot be sure that all the considerations in this section are correct. 8. Conclusion It is clear that all supergravity theories can be formulated as non-linear realisations. The bosonic part of the group underlying these constructions will include, the conformal group, the general linear group and certain automorphisms. In the complete theory, the conformal group will be embedded in the relevant Osp group which automatically contain the automorphisms of the Poincare supersymmetry algebras with all their central charges.
For many years it has been a puzzle to understand why the scalars that occur in supergravity theories belong to a non-linear realisation. However, from the perspective of this paper this it could be viewed as just a consequence of the whole theory being a non-linear realisation. It is known that if one reduces eleven dimensional supergravity on a torus one finds the group GL(11-d) in d dimensions. From the view point of the non-linear realisation of eleven dimensional supergravity given in this paper, this is hardly surprising since it is just part of the original GL(11) group of the original theory. However, it would be good to understand the emergence of the exceptional groups from the dimensional reduction of the non-linear realisations given in this paper.
One intriguing feature of the constructions of this paper is that the group is apparently different for each supergravity theory. This would be compatible with the suggestion, in references , that the full automorphism group is a symmetry of M theory and that as one goes to the limits of M theory such as eleven dimensional supergravity, IIA and IIB theory one finds different contractions of this automorphism group. In fact, it is inevitable that Osp(1/64) is a symmetry of M theory as this group is the unique extension of the conformal group to include supersymmetry and is in required in both the IIA and eleven dimensional supergravities. As the IIB supersymmetry algebra can be otained form the IIA supersymmetry algebra by an invertible transformation , it is likely that Osp(1/64) is also required in the IIB case. It is interesting to note that by taking this group one automatically encodes all the central charges and the GL(32) automorphism. As such, this group implicitly includes all the branes. Since it includes brane rotating symmetries one would have to restrict the field of the group to ensure it was compatible with the charge quantization conditions.
The maximal supergravities in ten dimensions are the low energy limits of the coresponding string theories. As such, it is perhaps not surprising that they should possess a non-linear realisation in that this has been the traditional role for such formulations. However, it does suggest that there is an alternative formulation of these string theories, perhaps M theory, in which all the symmetries discussed in this paper, including Osp(1/64) are linearly realised.
¿From a practical view point, it would be interesting to see if one could use the conformal symmetry, and its superextension, to derive constraints on the Greens functions of the supergravity theories. Such a calculation would utilise our knowledge of solving conformal Ward identities with theorems about the behaviour of Greens functions of Goldstone particles. In a sense gravity and supergravity can be thought of as the analogues of the conformally invariant two dimensional models. One can think of the symmetry of these latter models as being found by starting with the finite dimensional globally defined conformal group and generating an infinite dimensional group by taking its closure with the group whose generators are $`L_2,L_0,L_2`$. In the theories considered in this paper, one also starts with a finite dimensional, globally defined, extension of conformal group and generates an infinite dimensional group by taking its closure with an extension of the affine group. Acknowledgment The author would like to thank Bernard Julia for explaining the content of reference , Toine van Proeyen for discussions on superconformal symmetry in eleven dimensions, Arkardy Tseytlin for discussions on the closed bosonic string effective action and George Papadopoulos for commenting on the manuscript. This work was supported in part by the EU network on Integrability, Non-perturbative effects, and Symmetry in Quantum Field theory (FMRX-CT96-0012).
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# The 𝑄̂ operator for canonical quantum gravity
## 1 Introduction
It is well known that considerable progresses have been made in non-perturbative canonical quantum gravity in the light of Ashtekar’s variables and loop variables. One of the most remarkable physical results is the evidence for a quantum discreteness of space at the Planck scale. This is due to the fact that certain operators corresponding to the measurement of area and volume have discrete spectra . Also, the kinematical Hilbert space, $``$, of the theory has been rigorously defined by completing the space of all finite linear combinations of cylindrical functions in the norm induced by Haar measure .
There is another geometrical operator $`\widehat{Q}[\omega ]`$ proposed in previous literatures , which corresponds to the integrated norm of any smooth one form $`\omega _a`$ on the 3-manifold. While the operators of area and volume have been shown to be well defined self-adjoint operators on $``$, the general properties of $`\widehat{Q}[\omega ]`$ are still unclear. We even do not know if it is well defined on $``$. The obstacles are due to two facts: First, the result of $`\widehat{Q}[\omega ]`$ operating on a cylindrical function will involve integrals over edges of the graph on which the function defined, and hence it is no more a cylindrical function in general; Second, the current effective regularization technique of smearing the triads in 2-dimension could not be directly applied to the regularization of $`\widehat{Q}[\omega ]`$, whose classical expression involves the square of the triads while there is an integral over 3-dimensional manifold at last.
As $`\widehat{Q}[\omega ]`$ operator is rather convenient for constructing certain weave states in the study of the classical approximation of the quantum theory , the present paper is devoted to study the properties of $`\widehat{Q}[\omega ]`$ on $``$ in order to lay a foundation of its applications. To bypass the above mentioned obstacles, we will use a 3-dimensional smearing function for regularization. Then, instead of acting the regulated $`\widehat{Q}[\omega ]`$ on general cylindrical functions, we will operate it on spin network states which form a complete orthonormal basis in $``$. It turns out that the operation gives a real discrete spectrum, which is in the same form as its eigenvalues on coloured loop states. Thus, $`\widehat{Q}[\omega ]`$ is a well defined symmetric operator in $``$. A further discussion shows it is also self-adjoint.
We work in the real Ashtekar formalism defined over an oriented 3-manifold $`\mathrm{\Sigma }`$ . The basic variables are real $`SU(2)`$ connections, $`A_a^i`$, as the configuration and the densitized(weight 1) triads, $`E_j^b`$, corresponding to the conjugate momentum. We use $`a,b,\mathrm{}`$ for spatial indices and $`i,j,\mathrm{}`$ for internal $`SU(2)`$ indices. The basic variables satisfy
$$\{A_a^i(x),E_j^b(y)\}=G\delta _j^i\delta _a^b\delta ^3(x,y),$$
(1)
where $`G`$ is the usual gravitational constant.
## 2 $`\widehat{Q}`$ operator and its regularization
The operator $`\widehat{Q}[\omega ]`$ is constructed to represent the classical quantity
$$Q[\omega ]=d^3x\sqrt{E_i^a(x)\omega _a(x)E^{bi}(x)\omega _b(x)},$$
(2)
where $`\omega _a`$ is any smooth 1-form on $`\mathrm{\Sigma }`$ which makes the integral meaningful and the integral is well defined since the integrand is a density of weight 1. If we know $`Q[\omega ]`$ for all smooth $`\omega _a`$, the triad $`E_i^a`$ can be reconstructed up to local $`SU(2)`$ gauge transformations. Hence, the collection of $`Q[\omega ]`$ provides a good coordinates system on the space of the triads fields.
Since $`E_i^a`$ represents the conjugate momentum of the configuration variable $`A_a^i`$, the formal expression of the corresponding momentum operator would be some functional derivative with respect to $`A_a^i`$, i.e.,
$$\widehat{E}_i^a(x)=iG\mathrm{}\frac{\delta }{\delta A_a^i(x)}.$$
(3)
This is an operator-valued distribution rather than a genuine operator, hence it has to be integrated against smearing functions in order to be well defined. Our aim is to construct a well defined operator $`\widehat{Q}[\omega ]`$ corresponding to the classical quantity $`Q[\omega ]`$. We begin with a formal expression obtained by replacing $`E_i^a`$ in Eq.(2) by the operator-valued distribution $`\widehat{E}_i^a`$, and then regulate it by 3-dimensional smearing functions.
Let $`f_ϵ(x,y)`$ be a 1-parameter family of fields on $`\mathrm{\Sigma }`$ which tends to $`\delta (x,y)`$ as $`ϵ`$ tends to zero, such that $`f_ϵ(x,y)`$ is a density of weight 1 in $`x`$ and a function in $`y`$. We then define the smeared version of $`E_i^a(x)\omega _a(x)`$ as:
$$[E_i\omega ]_f(x):=d^3yf_ϵ(x,y)E_i^a(y)\omega _a(y).$$
(4)
Hence, $`[E_i\omega ]_f(x)`$ tends to $`E_i^a(x)\omega _a(x)`$ as $`ϵ`$ tends to zero. Then $`Q[\omega ]`$ can be regulated as:
$$Q[\omega ]=\underset{ϵ0}{lim}d^3x\left([E_i\omega ]_f(x)[E^i\omega ]_f(x)\right)^{\frac{1}{2}}.$$
(5)
To go over to the quantum theory, we simply replace $`E_i^a`$ by $`\widehat{E}_i^a`$ and obtain
$$\widehat{Q}[\omega ]=\underset{ϵ0}{lim}d^3x\left([\widehat{E}_i\omega ]_f(x)[\widehat{E}^i\omega ]_f(x)\right)^{\frac{1}{2}},$$
(6)
where
$$[\widehat{E}_i\omega ]_f(x):=d^3yf_ϵ(x,y)\widehat{E}_i^a(y)\omega _a(y)=iG\mathrm{}d^3yf_ϵ(x,y)\omega _a(y)\left(\frac{\delta }{\delta A_a^i(y)}\right).$$
(7)
By operating the regulated $`\widehat{Q}[\omega ]`$ on spin network states in next section, we will show that it is a well defined symmetric operator in the kinematical Hilbert space and admits self-adjoint extensions. For technical reasons, we attach the following concreteness to the smearing function $`f_ϵ`$ for sufficiently small $`ϵ>0`$: (i) $`f_ϵ(x,y)`$ is non-negative; and (ii) for any given $`y`$, $`f_ϵ(x,y)`$ has compact support in $`x`$ which is a 3-dimensional box, $`U_ϵ`$, of coordinate height $`ϵ^\beta `$, $`1<\beta <2`$, and square horizonal section, $`S_ϵ`$, of coordinate side $`ϵ`$, and with $`y`$ as its centre. These conditions are in the same spirit as that in Ref.. More concretely, $`f_ϵ(x,y)`$ can be constructed as follows. Take any 1-dimensional non-negative function $`\theta (x)`$ of compact support $`[\frac{1}{2},\frac{1}{2}]`$ on $`R`$ such that $`𝑑x\theta (x)=1`$, and set
$$f_ϵ(x,y)=\left(\frac{1}{ϵ^{2+\beta }}\right)\theta \left(\frac{x_1y_1}{ϵ}\right)\theta \left(\frac{x_2y_2}{ϵ}\right)\theta \left(\frac{x_3y_3}{ϵ^\beta }\right).$$
(8)
## 3 Action of $`\widehat{Q}`$ on spin network basis
### 3.1 Preliminaries
It has been shown that spin networks play a key role in non-perturbative quantum gravity . Consider a graph $`\mathrm{\Gamma }`$ with $`n`$ edges $`e_I`$, $`I=1,\mathrm{},n`$, and $`m`$ vertices $`v_\alpha `$, $`\alpha =1,\mathrm{},m`$, embedded in the 3-manifold $`\mathrm{\Sigma }`$. To each $`e_I`$ we assign a non-trivial irreducible spin $`j_I`$ representation of $`SU(2)`$. This is called a colouring of the edge. Next, consider a vertex $`v_\alpha `$, say a $`K`$-valent one, i.e., there are $`K`$ edges $`e_1,\mathrm{},e_K`$ meeting at $`v_\alpha `$. Let $`_{j_1},\mathrm{},_{j_K}`$ be the Hilbert spaces of the representations, $`j_1,\mathrm{},j_K`$, associated to the $`K`$ edges. Consider the tensor product of these spaces $`_{v_\alpha }=_{j_1}\mathrm{}_{j_K}`$, and fix, once and for all, an orthonormal basis, $`N_\alpha `$, in $`_{v_\alpha }`$. This is called a colouring of the vertex. A (non-gauge invariant) spin network, $`S`$, is then defined as the embedded graph whose edges and vertices have been coloured.
The holonomy of the $`SU(2)`$ connection $`A_a^i`$ along any edge $`e_I`$ is an element of $`SU(2)`$ and can be expressed as:
$$h[A,e_I]=𝒫exp_{e_I}𝑑s\dot{e}_I(s)A_a^i(e_I(s))\tau _i,$$
(9)
where $`𝒫`$ denotes path ordering and $`\tau _i`$ are the $`SU(2)`$ generators in the fundamental representation. The (non-gauge invariant) spin network state, $`\mathrm{\Psi }_S(A)`$, based on $`S`$ is defined as:
$$\mathrm{\Psi }_S(A)=\underset{e_I\mathrm{\Gamma }}{}j_I\left(h[e_I]\right)\underset{v_\alpha \mathrm{\Gamma }}{}N_\alpha ,$$
(10)
where $`j_I\left(h[e_I]\right)`$ is the representation matrix of the holonomy $`h[e_I]`$ in the spin $`j_I`$ representation associated to the edge $`e_I`$, and the holonomy matrices are constructed with the vector $`N_\alpha `$ at each vertex $`v_\alpha `$ where the edges meet. By varying the graph, the colours of the edges, and the colours of the vertices, we obtain a family of spin network states. It turns out that these states form a complete orthonormal basis in the kinematical Hilbert space $``$ .
Since a $`SU(2)`$ gauge transformation acts on a spin network state simply by $`SU(2)`$ transforming the colouring of the vertices $`N_\alpha `$, it is easy to recover the gauge invariant<sup>1</sup><sup>1</sup>1The gauge invariance discussed in this paper is restricted to that of internal $`SU(2)`$, while the whole gauge invariance of a gravitational theory should also involve that of 4-dimensional diffeomorphism. spin network states by colouring each vertex with a $`SU(2)`$ invariant basis. These states form a complete orthonormal basis in the $`SU(2)`$ gauge invariant Hilbert space $`_0`$ .
It is obvious from Eqs. (3) and (9) that the action of $`\widehat{E}_i^a(x)`$ on a holonomy $`h[e_I]`$ yields
$$\widehat{E}_i^a(x)h[e_I]=il_p^2_{e_I}𝑑s\dot{e}_I^a(s)\delta ^3(x,e_I(s))h_I[1,s]\tau _ih_I[s,0],$$
(11)
where $`l_p=\sqrt{G\mathrm{}}`$ is the Planck length.
### 3.2 Spectrum of $`\widehat{Q}`$
We first apply the operator $`[\widehat{E}_i\omega ]_f(x)`$ defined by Eq.(7) to the spin network state $`\mathrm{\Psi }_S`$,
$`[\widehat{E}_i\omega ]_f(x)\mathrm{\Psi }_S(A)`$ $`=`$ $`il_p^2{\displaystyle \underset{I=1}{\overset{n}{}}}{\displaystyle d^3yf_ϵ(x,y)\omega _a(y)\left[\frac{\delta }{\delta A_a^i(y)}j_I(h_I)_{lm}\right]\left(\frac{\mathrm{\Psi }_S}{j_I(h_I)_{lm}}\right)}`$ (12)
$`=`$ $`il_p^2{\displaystyle \underset{I=1}{\overset{n}{}}}{\displaystyle d^3yf_ϵ(x,y)_{e_I}𝑑t\dot{e}_I^a(t)\omega _a(y)\delta ^3(y,e_I(t))j_I\left(h_I[1,t]\tau _ih_I[t,0]\right)_{lm}\left(\frac{\mathrm{\Psi }_S}{j_I(h_I)_{lm}}\right)}`$
$`=`$ $`il_p^2{\displaystyle \underset{I=1}{\overset{n}{}}}{\displaystyle _{e_I}}𝑑t\dot{e}_I^a(t)\omega _a(e_I(t))f_ϵ(x,e_I(t))Tr\left[j_I\left(h_I[1,t]\tau _ih_I[t,0]{\displaystyle \frac{}{h_I}}\right)\right]\mathrm{\Psi }_S(A),`$
where $`l`$ and $`m`$ are indices in $`_{j_I}`$ associated to $`e_I`$. Repeating the action of $`[\widehat{E}^i\omega ]_f(x)`$ on Eq.(12), we have
$`[\widehat{E}^i\omega ]_f(x)[\widehat{E}_i\omega ]_f(x)\mathrm{\Psi }_S`$ $`=`$ $`l_p^4{\displaystyle \underset{I=1}{\overset{n}{}}}{\displaystyle _{e_I}}𝑑t\dot{e}_I^a(t)\omega _a(e_I(t)){\displaystyle \underset{J=1}{\overset{n}{}}}{\displaystyle _{e_J}}𝑑s\dot{e}_J^b(s)\omega _b(e_J(s))f_ϵ(x,e_I(t))f_ϵ(x,e_J(s))`$ (13)
$`Tr\left[j_J\left(h_J[1,s]\tau ^ih_J[s,0]{\displaystyle \frac{}{h_J}}\right)\right]Tr\left[j_I\left(h_I[1,t]\tau _ih_I[t,0]{\displaystyle \frac{}{h_I}}\right)\right]\mathrm{\Psi }_S`$
$`l_p^4{\displaystyle \underset{I=1}{\overset{n}{}}}{\displaystyle _{e_I}}dt\dot{e}_I^a(t)\omega _a(e_I(t))f_ϵ(x,e_I(t))({\displaystyle _t^1}ds\dot{e}_I^b(s)\omega _b(e_I(s))f_ϵ(x,e_I(s))`$
$`Tr\left[j_I\left(h_I[1,s]\tau ^ih_I[s,t]\tau _ih_I[t,0]{\displaystyle \frac{}{h_I}}\right)\right]+{\displaystyle _0^t}𝑑s\dot{e}_I^b(s)\omega _b(e_I(s))f_ϵ(x,e_I(s))`$
$`Tr\left[j_I\left(h_I[1,t]\tau _ih_I[t,s]\tau ^ih_I[s,0]{\displaystyle \frac{}{h_I}}\right)\right])\mathrm{\Psi }_S.`$
We denote respectively the first and second terms in the right hand side of Eq.(13) as $`A`$ and $`B`$. Consider first the term $`A`$. Note that a spin network state can always be written as:
$$\mathrm{\Psi }_S(A)=j_I\left(h[e_I]\right)_{lm}\mathrm{\Psi }_{Se_I}^{lm}(A),$$
(14)
where $`\mathrm{\Psi }_{Se_I}^{lm}(A)=\frac{\mathrm{\Psi }_S}{j_I(h_I)_{lm}}`$ is independent of $`h[e_I]`$. Hence, we can choose $`ϵ`$ sufficiently small, such that the term $`A`$ vanishes unless the support $`U_ϵ`$ of the smearing function $`f_ϵ`$ contains a vertex $`v_\alpha `$ of the spin network as its centre. It then turns out
$`A`$ $`=`$ $`l_p^4{\displaystyle \underset{I,J=1}{\overset{n}{}}}{\displaystyle _{e_I}}𝑑t\dot{e}_I^a(t)\omega _a(e_I(t)){\displaystyle _{e_J}}𝑑s\dot{e}_J^b(s)\omega _b(e_J(s))[f_ϵ(x,v_{IJ})]^2`$ (15)
$`Tr\left[j_J\left(h_J[1,s]\tau _ih_J[s,0]{\displaystyle \frac{}{h_J}}\right)\right]Tr\left[j_I\left(h_I[1,t]\tau _ih_I[t,0]{\displaystyle \frac{}{h_I}}\right)\right]\mathrm{\Psi }_S`$
$`=`$ $`l_p^2{\displaystyle \underset{\alpha =1}{\overset{m}{}}}[f_ϵ(x,v_\alpha )]^2{\displaystyle \underset{I_\alpha ,J_\alpha }{}}{\displaystyle _{e_I}}𝑑t\dot{e}_I^a(t)\omega _a(e_I(t)){\displaystyle _{e_J}}𝑑s\dot{e}_J^b(s)\omega _b(e_J(s))`$
$`Tr\left[j_J\left(h_J[1,s]\tau _ih_J[s,0]{\displaystyle \frac{}{h_J}}\right)\right]Tr\left[j_I\left(h_I[1,t]\tau _ih_I[t,0]{\displaystyle \frac{}{h_I}}\right)\right]\mathrm{\Psi }_S.`$
To simplify technicalities, given a 1-form $`\omega _a`$ we choose $`f_ϵ`$ such that at each vertex, $`\omega _a`$ is a normal covector of the horizonal section $`S_ϵ`$ of $`U_ϵ`$, i.e., $`\omega _a(v_\alpha )=|\omega (v_\alpha )|(dx_3)_a`$. (Note that the special $`f_ϵ`$ is chosen here in order to obtain a succinct expression of $`A`$, see Eq.(23), while the final result that $`A`$ makes no contribution to the spectrum of $`\widehat{Q}`$ is independent of this choice.) Since $`ϵ^\beta `$ goes to zero faster than $`ϵ`$, for sufficiently small $`ϵ`$ the edge $`e_I`$ which meets the vertex $`v_\alpha `$ would cross the top or bottom of the box $`U_ϵ`$ if it is not tangent to $`S_ϵ`$ at $`v_\alpha `$. Also, it follows from $`\beta <2`$ that any edge tangential to $`S_ϵ`$ at $`v_\alpha `$ exits $`U_ϵ`$ from the side, irrespectively from its second (and higher) derivatives, and gives a vanishing contribution to $`A`$ as $`ϵ`$ goes to zero. Thus, if we first consider only the “outgoing” edges from the vertices, it turns out
$$_{e_I}𝑑t\dot{e}_I^a(t)\omega _a(e_I(t))f_ϵ(x,v_\alpha )=\frac{1}{2}\kappa _Iϵ^\beta |\omega (v_\alpha )|f_ϵ(x,v_\alpha )+O(ϵ),$$
(16)
where
$`\kappa _I:=\{\begin{array}{ccc}0,\hfill & \text{ if }e_I\text{ is tangent to }S_ϵ\hfill & \\ 1,\hfill & \text{ if }e_I\text{ lies above }S_ϵ\hfill & \\ 1,\hfill & \text{ if }e_I\text{ lies below }S_ϵ\hfill & \end{array}`$ (20)
and hence Eq.(15) becomes
$$A_{out}=\frac{1}{4}l_p^4\underset{\alpha =1}{\overset{m}{}}[f_ϵ(x,v_\alpha )ϵ^\beta |\omega (v_\alpha )|]^2\underset{I_\alpha ,J_\alpha }{}[\kappa _I\kappa _JL_I^iL_J^i+O(ϵ)]\mathrm{\Psi }_S,$$
(21)
where
$$L_I^i\mathrm{\Psi }_S:=Tr\left[j_I\left(h[e_I]\tau ^i\frac{}{h[e_I]}\right)\right]\mathrm{\Psi }_S.$$
(22)
Including the “incoming” edges to the vertices, the final expression of $`A`$ reads
$$A=\frac{1}{4}l_p^4ϵ^{2\beta }\underset{\alpha =1}{\overset{m}{}}[f_ϵ(x,v_\alpha )|\omega (v_\alpha )|]^2\underset{I_\alpha ,J_\alpha }{}[\kappa _I\kappa _JX_I^iX_J^i+O(ϵ)]\mathrm{\Psi }_S,$$
(23)
where
$`X_I^i\mathrm{\Psi }_S:=\{\begin{array}{cc}Tr\left[j_I\left(h[e_I]\tau ^i\frac{}{h[e_I]}\right)\right]\mathrm{\Psi }_S,\hfill & \text{if }e_I\text{ is outgoing}\hfill \\ Tr\left[j_I\left(\tau ^ih[e_I]\frac{}{h[e_I]}\right)\right]\mathrm{\Psi }_S,\hfill & \text{if }e_I\text{ is incoming.}\hfill \end{array}`$ (26)
Note that $`\mathrm{\Delta }_{S_ϵ,v_\alpha }=_{I_\alpha ,J_\alpha }\kappa _I\kappa _JX_I^iX_J^i`$ is the vertex operator associated with $`S_ϵ`$ and $`v_\alpha `$ in arbitrary spin representations, which has been fully investigated . A discussion similar to that in Ref. leads that the spin network state $`\mathrm{\Psi }_S`$ is an eigenvector of $`\mathrm{\Delta }_{S_ϵ,v_\alpha }`$ with eigenvalue:
$$\lambda _{S_ϵ,v_\alpha }=2j^{(d)}(j^{(d)}+1)+2j^{(u)}(j^{(u)}+1)j^{(d+u)}(j^{(d+u)}+1),$$
(27)
where $`j^{(d)}`$, $`j^{(u)}`$ and $`j^{(d+u)}`$ are half integers subject to the condition
$`j^{(d+u)}\{|j^{(d)}j^{(u)}|,|j^{(d)}j^{(u)}|+1,\mathrm{},j^{(d)}+j^{(u)}\}`$.
Now consider the second term $`B`$. For small $`ϵ`$, $`f_ϵ(x,e_I(t))f_ϵ(x,e_I(s))`$ is non-zero only for the parameters satisfying $`t=s+O(ϵ)`$, where we have
$`Tr\left[j_I\left(h_I[1,t]\tau _ih_I[t,s]\tau ^ih_I[s,0]{\displaystyle \frac{}{h_I}}\right)\right]\mathrm{\Psi }_S`$ $`=`$ $`\left(Tr\left[j_I\left(h_I[1,t]\tau _i\tau ^ih_I[t,0]{\displaystyle \frac{}{h_I}}\right)\right]+O(ϵ)\right)\mathrm{\Psi }_S`$ (28)
$`=`$ $`\left(j_I(j_I+1)Tr\left[j_I\left(h_I[e_I]{\displaystyle \frac{}{h[e_I]}}\right)\right]+O(ϵ)\right)\mathrm{\Psi }_S`$
$`=`$ $`[j_I(j_I+1)+O(ϵ)]\mathrm{\Psi }_S,`$
here $`j_I(\tau _i\tau ^i)=j_I(j_I+1)`$ is the Casimir operator of $`SU(2)`$. Substituting Eq.(28) into $`B`$, we obtain
$$B=l_p^4\underset{I=1}{\overset{n}{}}\left[_{e_I}𝑑t\dot{e}_I^a(t)\omega _a(e_I(t))f_ϵ(x,e_I(t))\right]^2[j_I(j_I+1)+O(ϵ)]\mathrm{\Psi }_S.$$
(29)
It then follows from Eqs. (23) and (29) that
$`[\widehat{E}^i\omega ]_f(x)[\widehat{E}_i\omega ]_f(x)\mathrm{\Psi }_S={\displaystyle \frac{1}{4}}ϵ^{2\beta }l_p^4{\displaystyle \underset{\alpha =1}{\overset{m}{}}}[f_ϵ(x,v_\alpha )|\omega (v_\alpha )|]^2(\lambda _{S_ϵ,v_\alpha }+O(ϵ))\mathrm{\Psi }_S`$
$`+l_p^4{\displaystyle \underset{I=1}{\overset{n}{}}}\left[{\displaystyle _{e_I}}𝑑t\dot{e}_I^a(t)\omega _a(e_I(t))f_ϵ(x,e_I(t))\right]^2[j_I(j_I+1)+O(ϵ)]\mathrm{\Psi }_S,`$ (30)
which implies that $`[\widehat{E}^i\omega ]_f(x)[\widehat{E}_i\omega ]_f(x)`$ is a well defined non-negative operator and hence has a well defined square-root. Since we have chosen $`ϵ`$ to be sufficiently small, for any given $`x\mathrm{\Sigma }`$, $`f_ϵ(x,v_\alpha )`$ is non-zero for at most one vertex, and $`f_ϵ(x,e_I(t))`$ is non-zero for at most a piece of one edge where its vertices are not included. Therefore we can take the sum over $`v_\alpha `$ and $`e_I`$ out side the square root and obtain
$`\left([\widehat{E}^i\omega ]_f(x)[\widehat{E}_i\omega ]_f(x)\right)^{\frac{1}{2}}\mathrm{\Psi }_S={\displaystyle \frac{1}{2}}ϵ^\beta l_p^2{\displaystyle \underset{\alpha =1}{\overset{m}{}}}f_ϵ(x,v_\alpha )|\omega (v_\alpha )|(\lambda _{S_ϵ,v_\alpha }+O(ϵ))^{\frac{1}{2}}\mathrm{\Psi }_S`$
$`+l_p^2{\displaystyle \underset{I=1}{\overset{n}{}}}\left|{\displaystyle _{e_I}}𝑑t\dot{e}_I^a(t)\omega _a(e_I(t))f_ϵ(x,e_I(t))\right|[j_I(j_I+1)+O(ϵ)]^{\frac{1}{2}}\mathrm{\Psi }_S.`$ (31)
Now we can remove the regulator. Taking the limit $`ϵ0`$ and integrating over $`\mathrm{\Sigma }`$, the first term in the right hand side of Eq.(3.2) vanishes due to the factor $`ϵ^\beta `$. We thus conclude that the action of $`\widehat{Q}[\omega ]`$ on spin network states yields
$$\widehat{Q}[\omega ]\mathrm{\Psi }_S(A)=l_p^2\underset{I=1}{\overset{n}{}}\left[_{e_I}𝑑t|\dot{e}_I^a(t)\omega _a(e_I(t))|\sqrt{j_I(j_I+1)}\right]\mathrm{\Psi }_S(A).$$
(32)
Therefore, spin network states are also eigenvectors of $`\widehat{Q}[\omega ]`$. The complete spectrum of $`\widehat{Q}[\omega ]`$ with respect to the spin network basis in the Hilbert space $``$ is obtained.
## 4 Discussions
The general properties of $`\widehat{Q}[\omega ]`$ operator are implied by Eq.(32). In contrast to the volume operator which acts only on vertices , $`\widehat{Q}[\omega ]`$ acts only on edges of spin networks. Hence, the spin network states based on a same graph with same colouring of the edges are all degenerate with respect to this operator. As a result, the action of $`\widehat{Q}[\omega ]`$ on the gauge invariant spin network states gives the same result as Eq.(32). In this sense, the spectrum of $`\widehat{Q}[\omega ]`$ respects the physically relevant states in $``$.
There are alternative approaches to regulate $`\widehat{Q}[\omega ]`$ and calculate its spectrum. One could also apply the blocking regularization technique of Ref., then express $`\widehat{Q}[\omega ]`$ by the loop operator $`𝒯^{ab}`$ up to $`O(ϵ)`$, whose action on spin network states is obtained from the recoupling theory . It is not difficult to check that this approach will give the same result as we have obtained. By restricting the support of the regulator, our approach reveals the inherent relation of the two approaches.
We have shown that $`\widehat{Q}[\omega ]`$ is diagonalized in the spin network basis with real eigenvalues, hence it is a well defined symmetric operator in the kinematical Hilbert space $``$. Moreover, it is obvious from Eqs. (6) and (7) that the expression of $`\widehat{Q}[\omega ]`$ is purely real, and hence it commutes with the complex conjugation. Therefore, it follows from Von Neumann’s theorem that $`\widehat{Q}[\omega ]`$ admits self-adjoint extensions on $``$. The same reasons lead that $`\widehat{Q}[\omega ]`$ is also self-adjoint on the gauge invariant Hilbert space $`_0`$.
The discrete spectrum of $`\widehat{Q}[\omega ]`$ shows a quantum discreteness of the space at the Planck scale, corresponding to the measurement of the integrated norm of any smooth one forms.
### Acknowledgements
Y. Ma would like to thank discussions with Profs. C. N. Kozameh, O. M. Moreschi, and especially O. A. Reula for his enlightening comments and suggestions, and acknowledge support from FONCYT BID 802/OC-AR PICT: 00223. Y. Ling is grateful to Abhay Ashtekar for discussion, and was supported by the NSF through grant PHY95-14240, a gift from the Jesse Phillips Foundation, and a Braddock fellowship from the Department of Physics at Pennsylvania State University.
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# Limit Theorems for Height Fluctuations in a Class of Discrete Space and Time Growth Models
## 1 Introduction
Growth processes have been extensively studied by mathematicians and physicists for many years (see, e.g., and references therein), but it was only recently that K. Johansson proved that the fluctuations of the limiting shape in a class of growth models are described by certain distribution functions first appearing in random matrix theory (RMT) . Further work by Johansson , Prähofer and Spohn and Baik and Rains strongly suggests the universal nature of these RMT distribution functions. These developments are part of the recent activity relating Robinson-Schensted-Knuth (RSK) type problems of combinatorial probability to the distribution functions of RMT, see e.g. .
In this paper we analyze a class of one-dimensional discrete space-discrete time stochastic growth models, called oriented digital boiling . Digital boiling dynamics is a cellular automaton that models an excitable medium in the presence of persistent random spontaneous excitation. Alternatively, digital boiling models represent contour (constant height) lines for one of the simplest models for growing a connected interface. It is this latter point of view we adopt here; that is, we introduce a height function $`h_t(x)`$ that characterizes the state of the system. Fig. 1 illustrates the height fluctuations in oriented digital boiling.
We shall derive various limit theorems for $`h_t(x)`$. We find four limiting regimes:
1. GUE Universal Regime: $`x\mathrm{}`$, $`t\mathrm{}`$ such that $`p_c:=1x/t`$ is fixed and $`p<p_c`$.
2. Critical Regime: $`x\mathrm{}`$, $`t\mathrm{}`$ such that $`p_c:=1x/t`$ is fixed and $`pp_c`$.
3. Deterministic Regime: $`x\mathrm{}`$, $`t\mathrm{}`$ such that $`p_c:=1x/t`$ is fixed and $`p>p_c`$.
4. Finite $`x`$ GUE Regime: Fixed $`x`$ and $`t\mathrm{}`$.
The limit theorems are stated at the beginning of §3. Here is an outline of how they are obtained. First we show that $`h_t(x)`$ satisfies a last passage property, i.e. it equals the maximum over a certain class of paths in space-time. Then applying the dual RSK algorithm , we obtain a reformulation of the problem in terms of Young tableaux. This is followed by an application of a theorem of Gessel (see also ) which gives a Toeplitz determinant representation for the distribution function for $`h_t(x)`$. An identity of Borodin and Okounkov expresses the Toeplitz determinant in terms of the Fredholm determinant of an infinite matrix. Finally we use a saddle point analysis (steepest descent) to determinine the limiting behavior of the entries, and therefore the Fredholm determinant, of the infinite matrix.
Along the way we identify<sup>1</sup><sup>1</sup>1A referee points out that this identification can be made at the very beginning. ODB with a first-passage percolation model of Seppäläinen whose limit law in the universal regime was determined by Johansson . Thus we could have used the analysis in to establish our limit law in the universal regime, or alternatively used Riemann-Hilbert methods , to investigate the Toeplitz determinant asymptotics. But the method we present is in our opinion more straightforward and technically simpler than these, and it is very general. (The Fredholm determinant is easier to handle than the Toeplitz determinant, even though they are essentially equal.) Also, our analysis permits a nice conceptual understanding of the various limiting regimes. For example, the universal regime is characterized by the coalescence of two saddle points; and the emergence of the Airy kernel is related to the well-known appearance of Airy functions in such a saddle point analysis .
Even in this simpler approach there are technical details to work out after the saddle point analysis gives us the answer. For example in the universal regime we need uniform estimates on the entries of the infinite matrix in order to show that the matrix scales in trace norm to the Airy kernel. These details are given completely only for this regime.
In Regime 4 we give an independent proof that the suitably centered and normalized $`h_t(x)`$ has a limiting distribution. The proof proceeds through the introduction of a certain Brownian motion functional. This leads to some apparently new identities for $`n`$-dimensional Brownian motion; see (4.29) below.
The initial conditions are corner initialization. Due to the fact there is no known symmetry theorem for the dual RSK algorithm , we are unable to prove limit theorems with different initial conditions, e.g. growth from a flat substrate. From work of Baik and Rains and Prähofer and Spohn , it is natural to conjecture that the limiting distribution is now of GOE symmetry and hence given by the analogous distribution function in the GOE case .
The table of contents provides a detailed description of the organization of this paper.
## 2 Growth Models and Increasing Paths
In this section we introduce three classes of discrete space and discrete time stochastic growth models. Each of these models will have an equivalent path description, but only for one of these models are we able to prove limit theorems. Nevertheless, we believe it is useful to place this “solvable” case in a larger context.
We assume that the occupied set of our growth models can be described by a height function $`h_t:\text{Z}_+\text{Z}_+\{\mathrm{},\mathrm{}\}`$, where $`\text{Z}_+`$ is the set of nonnegative integers. Here, time $`t=0,1,2,\mathrm{}`$ proceeds in discrete steps. The occupied set at time $`t`$ is thus given by
$$\eta _t=\{(x,y)\text{Z}_+\times \text{Z}_+:yh_t(x)\}.$$
In the models below we use the following one-dimensional neighborhood: $`(x+𝒩)=\{x1,x\}`$ (the oriented case) and assume corner initialization,
$$h_0(x)=\{\begin{array}{cc}\hfill 0,& \text{if}x=0,\hfill \\ \hfill \mathrm{},& \text{otherwise}.\hfill \end{array}$$
(2.1)
### 2.1 Oriented Digital Boiling
The first class of growth rules we call oriented digital boiling (ODB) .<sup>2</sup><sup>2</sup>2For spatial dimensions greater than one, visual features of this dynamics resemble bubble formation, growth and annihilation in a boiling liquid, hence the process is called digital boiling and oriented refers to the choice of neighborhood $`𝒩`$, see Fig. 2 in or Feb. 12, 1996 Recipe of . The rules for ODB are
1. $`h_t(x)h_{t+1}(x)`$ for all $`x`$ and $`t`$.
2. If $`h_t(x1)>h_t(x)`$, then $`h_{t+1}(x)=h_t(x1)`$.
3. Otherwise, then independently of the other sites and other times, $`h_{t+1}(x)=h_t(x)+1`$ with probability $`p`$. (With probability $`1p`$, we have $`h_{t+1}(x)=h_t(x)`$.)
It follows from these rules that for every $`x`$ and $`t`$, $`h_t(x1)h_t(x)+1`$.
This process can be readily visualized by imagining the growth proceeding by the addition of unit squares starting with the initial square centered at $`(1/2,1/2)`$. We denote this initial time by placing a $`0`$ in this box. At time $`t=1`$ a box is added to the right (centered at $`(3/2,1/2)`$) and with probability $`p`$ a box is added to the top of the initial box (centered at $`(1/2,1/2)`$). We place $`1`$’s in the boxes added at time $`t=1`$. The boxes that are added stochastically (Rule (3)) are shaded. An example of this process run for seven time steps is shown in Fig. 2.
#### 2.1.1 Path Description
As has been observed many times before (see, e.g., ), a most productive way to think about height processes is to introduce the (discrete) backwards lightcone of a point $`(x,t)`$. Precisely, if $`𝒮=\text{Z}_+\times \text{Z}_+`$ denotes space-time, then
$$_B(x,t)=\{(x^{},t^{})𝒮:0x^{}x,x^{}t^{}<x^{}+tx\}.$$
For those space-time points in $`_B`$ at which a box was added stochastically (according to Rule (3)), we place a $`\times `$. We call such space-time points marked. We define the length of a sequence $`\pi =\{(x_1,t_1),\mathrm{},(x_k,t_k)\}`$ of (distinct) space-time points in $`_B`$ to be $`k`$. Such a sequence $`\pi `$ is increasing if $`0x_ix_{i1}t_it_{i1}1`$ for $`i=2,\mathrm{},k`$. Let $`L(x,t)`$ equal the length of the longest increasing sequence of marked space-time points in $`_B(x,t)`$. (If $`xt`$ and $`_B(x,t)`$ contains no increasing path, then $`L(x,t):=0`$.) For the example in Fig. 2, the discrete backwards lightcone $`_B(3,7)`$ and an increasing path are shown in Fig. 3. One observes that $`h_7(3)=L(3,7)=4`$. Indeed, this is a general fact. However, before proceeding with its proof, it is useful to change slightly the point of view of the process defined by Rules (1)–(3). Let $`\mathrm{\Pi }=\mathrm{\Pi }(p)`$ be a random subset of $`𝒮`$ to which every point of $`𝒮`$ belongs with probability $`p`$. We mark the points of $`𝒮`$ that belong to $`\mathrm{\Pi }`$. (Accordingly, we call these points marked.) $`L(x,t)`$ remains the same; namely, the length of the longest increasing sequence of marked space-time points in $`_B(x,t)`$. With regard to the process, we may intuitively think that all the “coins” used in Rule (3) are thrown in advance—of course, many of these are ignored as $`x`$ at time $`t`$ may become occupied deterministically by Rule (2). Precisely, Rule (3) is replaced with
1. Otherwise,
$$h_{t+1}(x)=\{\begin{array}{cc}\hfill h_t(x)+1,& \text{if}(x,t)\mathrm{\Pi },\hfill \\ \hfill h_t(x),& \text{if}(x,t)\mathrm{\Pi }.\hfill \end{array}$$
We are now ready to prove the last passage property
Proposition . $`h_t(x)=L(x,t)`$.
Proof. We first show that our process is attractive<sup>3</sup><sup>3</sup>3For examples of the kind of exotic shapes that can occur from cellular automaton rules without this monotonicity property, see . in the following sense: Let $`\mathrm{\Pi }`$ and $`\mathrm{\Pi }^{}`$ be two sets of marked points such that $`\mathrm{\Pi }\mathrm{\Pi }^{}`$. Let $`h_t`$ evolve using $`\mathrm{\Pi }`$ and $`h_t^{}`$ using $`\mathrm{\Pi }^{}`$, then $`h_th_t^{}`$ for all $`t`$. For if this were not true, then, for some $`t`$, $`h_sh_s^{}`$, $`st`$, and $`h_{t+1}(x)>h_{t+1}^{}(x)`$ for some $`x`$. This, of course, implies that $`h_t(x)=h_t^{}(x)`$. But then $`h_{t+1}(x)=h_t(x)+1`$ either because of Rule (2); in which case, $`h_t^{}(x1)h_t(x1)>h_t(x)=h_t^{}(x)`$, so (by Rule (2)) $`h_{t+1}^{}(x)=h_t^{}(x)+1`$; or, because $`(x,t)\mathrm{\Pi }\mathrm{\Pi }^{}`$, so again $`h_{t+1}^{}(x)=h_t^{}(x)+1`$. This is a contradiction. Thus we’ve established the attractiveness of our process.
The property of attractiveness immediately implies $`h_t(x)L(x,t)`$, since any increasing path of length $`k`$ will, without the addition of other marked points, cause $`h_t(x)k`$.
We now show that $`h_t(x)L(x,t)`$. We will show, by induction on $`k`$ and $`t`$, that $`h_t(x)=k`$ implies there exists an increasing sequence of marked points of length $`k`$ in $`_B(x,t)`$. This is obviously true for either $`t=0`$ or $`k=0`$. (Note that $`h_t(x)0`$ means that $`xt`$.) Now assume the claim has been demonstrated for all $`k^{}<k`$ and $`t^{}<t`$. We can clearly assume that $`h_{t1}(x)=k1`$, or else we can use the induction hypothesis right away. Therefore, we have two possibilities.
Case 1. $`h_t(x)=h_{t1}(x)+1`$ by application of Rule (2). This means (by Rule (2)) that $`h_{t1}(x1)=k`$. Thus by the induction hypothesis, there is an increasing sequence of length $`k`$ in $`_B(x1,t1)_B(x,t)`$.
Case 2. $`h_t(x)=h_{t1}(x)+1`$ by application of Rule ($`3^{}`$). This means that $`h_{t1}(x)=k1`$ and $`(x,t1)\mathrm{\Pi }`$. By the induction hypothesis, $`_B(x,t1)`$ contains an increasing path of length $`k1`$. Adjoin the marked point $`(x,t1)`$ to the sequence. Observe that the increasing property is preserved. This completes the proof of the proposition.
We summarize this section by noting that $`h_t`$ satisfies for all $`t1`$, $`x0`$,
$$h_t(x)=\mathrm{max}\{h_{t1}(x1),h_{t1}(x)+ϵ_{x,t}\}$$
where $`ϵ_{x,t}=1`$ if $`(x,t)\mathrm{\Pi }`$ and $`0`$ otherwise. The initial conditions are (2.1). (We take $`h_t(1)=\mathrm{}`$.) Formulated this way ODB is a “stochastic dynamic programming” problem.
#### 2.1.2 The $`(0,1)`$-Matrix Description of ODB
Without changing the increasing path property, the backwards lightcone $`_B`$ of any space-time point $`(x,t)`$ can be deformed into a rectangle of size $`(tx)\times (x+1)`$. Thus the equivalent problem is to fix $`x`$ and $`t`$ and to set $`m=tx`$, $`n=x+1`$, and to consider a $`(0,1)`$-matrix $`A`$ of size $`m\times n`$. We number the rows of $`A`$ starting at the bottom of $`A`$ and the columns of $`A`$ starting at the left of $`A`$. A increasing path in $`_B`$ becomes a sequence of $`1`$’s in $`A`$ at, say, positions $`\{(i_1,j_1),\mathrm{},(i_k,j_k)\}`$ such that the $`i_{\mathrm{}}`$ ($`\mathrm{}=1,\mathrm{},k`$) are increasing and the $`j_{\mathrm{}}`$ ($`\mathrm{}=1,\mathrm{},k`$) are weakly increasing. Any such $`(0,1)`$-matrix $`A`$ of size $`m\times n`$ corresponds (bijectively) to a two-line array<sup>4</sup><sup>4</sup>4We have chosen both a nonstandard labeling of $`A`$ and a nonstandard bijection $`Aw_A`$ so that our increasing path property remains (essentially) the same under the bijections.
$$w_A=\left(\begin{array}{cccc}j_1& j_2& \mathrm{}& j_k\\ i_1& i_2& \mathrm{}& i_k\end{array}\right)$$
(2.2)
where $`j_1j_2\mathrm{}j_k`$ and if $`j_{\mathrm{}}=j_{\mathrm{}+1}`$, then $`i_{\mathrm{}}<i_{\mathrm{}+1}`$ and the pair $`\left(\genfrac{}{}{0pt}{}{j}{i}\right)`$ appears in $`w_A`$ if and only if the $`(i,j)`$ entry of $`A`$ is 1. Note that the upper numbers belong to $`\{1,2,\mathrm{},n\}`$ and the lower numbers to $`\{1,2,\mathrm{},m\}`$. For example, the matrix
$$A=\left(\begin{array}{ccccccc}0& 0& 0& 1& 0& 0& 1\\ 1& 1& 1& 1& 0& 1& 1\\ 1& 1& 0& 0& 0& 1& 0\\ 1& 0& 1& 1& 0& 1& 1\\ 0& 1& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 1\end{array}\right)$$
maps to the two-line array
$`w_A`$ $`=`$ $`\left(\begin{array}{cccccccccccccccccccc}1& 1& 1& \text{2}& 2& 2& 3& \text{3}& 3& 4& 4& 4& 5& 6& \text{6}& \text{6}& 7& 7& 7& \text{7}\\ 3& 4& 5& \text{2}& 4& 5& 2& \text{3}& 5& 3& 5& 6& 1& 3& \text{4}& \text{5}& 1& 3& 5& \text{6}\end{array}\right).`$
(Recall the convention for row labels.) As an example, a longest increasing path (of length 5) is indicated in bold typeface. We remark that one can compute the length of an increasing path by patience sorting on the bottom row of $`w_A`$ (from left to right) with the rule that a number is placed on the left most pile such that it is less than or equal to the number showing in the pile. Patience sorting on the above example results in the five piles
$$\begin{array}{ccccc}& 3& & & \\ 1& 3& 4& & \\ 1& 3& 5& & \\ 2& 3& 5& 5& \\ 2& 4& 5& 5& \\ 3& 4& 5& 6& 6.\end{array}$$
If $`N`$ denotes the number of $`1`$’s in a random $`m\times n`$ $`(0,1)`$-matrix $`A`$, then the above mappings imply that for any nonnegative integer $`h`$,
$`\text{Prob}\left(h_t(x)h\right)`$ $`=`$ $`{\displaystyle \underset{k0}{}}\text{Prob}\left(h_t(x)h|N=k\right)\text{Prob}\left(N=k\right)`$ (2.4)
$`=`$ $`{\displaystyle \underset{k=0}{\overset{mn}{}}}\left({\displaystyle \genfrac{}{}{0pt}{}{mn}{k}}\right)p^k(1p)^{mnk}\text{Prob}\left(L_{m,n,k}h\right)`$
where $`L_{m,n,k}`$ is the length of the longest increasing path in a random (0,1)-matrix $`A`$ with $`k`$ 1’s (or equivalently, in the associated $`w_A`$).
#### 2.1.3 Tableaux Description of ODB
The dual RSK algorithm is a bijection between $`(0,1)`$-matrices $`A`$ of size $`m\times n`$ and pairs $`(P,Q)`$ such that $`P^t`$ (the transpose of $`P`$) and $`Q`$ are semistandard Young tableaux (SSYTs) with $`\text{sh}(P)=\text{sh}(Q)`$ where the elements of $`P`$ are from $`\{1,2,\mathrm{},m\}`$ and the elements of $`Q`$ are from $`\{1,2,\mathrm{},n\}`$. In terms of the associated $`w_A`$, (2.2), one forms $`P`$ by successive row bumping of the second row of $`w_A`$ starting with $`i_1`$ and with the rule an element $`i`$ bumps the leftmost element $`i`$. Thus each row of $`P`$ is strictly increasing. A fundamental property of the dual RSK algorithm is that the length of the longest strictly increasing subsequence of the second row of $`w_A`$ equals the number of boxes in the first row of $`P`$.
If $`d_\lambda (M)`$ denotes the number of SSTYs of shape $`\lambda `$ with entries coming from $`\{1,2,\mathrm{},M\}`$, then the number of pairs $`(P,Q)`$ of fixed shape $`\lambda `$ in the above dual RSK algorithm is
$$d_\lambda ^{}(m)d_\lambda (n)$$
where $`\lambda ^{}`$ is the conjugate partition. (Conjugate since $`P^t`$ is a SSYT.) Since there are $`\left(\genfrac{}{}{0pt}{}{mn}{k}\right)`$ $`(0,1)`$-matrices with $`k`$ 1’s,
$$\text{Prob}\left(L_{m,n,k}h\right)=\frac{1}{\left(\genfrac{}{}{0pt}{}{mn}{k}\right)}\underset{\genfrac{}{}{0pt}{}{\lambda k}{\lambda _1h}}{}d_\lambda ^{}(m)d_\lambda (n)=\frac{1}{\left(\genfrac{}{}{0pt}{}{mn}{k}\right)}\underset{\genfrac{}{}{0pt}{}{\lambda k}{\mathrm{}(\lambda )h}}{}d_\lambda (m)d_\lambda ^{}(n).$$
And hence from (2.4)
$$\text{Prob}\left(h_t(x)h\right)=(1p)^{mn}\underset{k=0}{\overset{mn}{}}r^k\underset{\genfrac{}{}{0pt}{}{\lambda k}{\mathrm{}(\lambda )h}}{}d_\lambda (m)d_\lambda ^{}(n)$$
where $`r=p/(1p)`$. Observe that for $`|\lambda |>mn`$, $`d_\lambda (m)d_\lambda ^{}(n)=0`$. (A SSYT with entries from $`\{1,2,\mathrm{},M\}`$ can have at most $`M`$ rows.) If $`𝒫`$ denotes the set of all partitions (including the empty partition), then the above sum can be summed over all partitions without changing its value,
$$\text{Prob}\left(h_t(x)h\right)=(1p)^{mn}\underset{\genfrac{}{}{0pt}{}{\lambda 𝒫}{\mathrm{}(\lambda )h}}{}r^{|\lambda |}d_\lambda (m)d_\lambda ^{}(n).$$
(2.5)
Comparing (2.5) with Johansson’s Krawtchouck ensemble results establishes the equivalence of ODB with the Seppäläinen-Johansson model.
#### 2.1.4 Application of Gessel’s Theorem and the Borodin-Okounkov Identity
Gessel’s theorem is
$$\underset{\genfrac{}{}{0pt}{}{\lambda 𝒫}{\mathrm{}(\lambda )h}}{}r^{|\lambda |}s_\lambda (x)s_\lambda (y)=D_h(\phi )$$
where $`s_\lambda `$ are the Schur functions (see, e.g. ) and $`D_h(\phi )`$ is the $`h\times h`$ Toeplitz determinant<sup>5</sup><sup>5</sup>5If $`\varphi `$ is a function on the unit circle with Fourier coefficients $`\varphi _k`$ then $`T_n(\varphi )`$ denotes the Toeplitz matrix $`(\varphi _{ij})_{i,j=0,\mathrm{},n1}`$ and $`D_n(\varphi )`$ its determinant. with symbol
$$\phi (z)=\underset{j=1}{\overset{\mathrm{}}{}}(1x_jz)^1\underset{j=1}{\overset{\mathrm{}}{}}(1y_jrz^1)^1.$$
If we apply to both sides of this identity the automorphism $`\omega `$ (see Stanley , pg. 332), $`\omega (s_\lambda )=s_\lambda ^{}`$, to the symmetric functions in the $`x`$-variables we obtain
$$\underset{\genfrac{}{}{0pt}{}{\lambda 𝒫}{\mathrm{}(\lambda )h}}{}r^{|\lambda |}s_\lambda ^{}(x)s_\lambda (y)=D_h(\phi )$$
(2.6)
where now the symbol is
$$\phi (z)=\underset{j=1}{\overset{\mathrm{}}{}}(1+x_jz)\underset{j=1}{\overset{\mathrm{}}{}}(1y_jrz^1)^1.$$
(2.7)
Recalling the specialization $`\text{ps}_n^1`$ (see Stanley , pg. 303), we apply $`\text{ps}_n^1`$ to the $`x`$-variables and $`\text{ps}_m^1`$ to the $`y`$-variables in Gessel’s identity (2.6) and observe<sup>6</sup><sup>6</sup>6Note that $`\text{ps}_n^1s_\lambda =d_\lambda (n)`$ which follows from the combinatorial definition of the Schur function. that the resulting LHS is precisely the RHS of (2.5). Since the specialization $`\text{ps}_n^1`$ is a ring homomorphism, we may apply it directly to the symbol (2.7). Doing so we obtain
$$\text{Prob}\left(h_t(x)h\right)=(1p)^{mn}D_h(\phi )$$
(2.8)
where
$$\phi (z)=(1+z)^n(1r/z)^m.$$
(2.9)
This derivation required $`r<1`$. However, by (2.4) the left side is a rational function of $`r`$, and analytic continuation shows that (2.8) holds for all $`r0`$ if in the integral representing the Fourier coefficients of $`\phi `$ the contour has $`r`$ on the inside.
The Borodin-Okounkov identity expresses a Toeplitz determinant in terms of a Fredholm determinant of an infinite matrix which in turn is a product of two Hankel matrices. Subsequent simplifications of the proof by Basor and Widom extended the identity to block Toeplitz determinants. We now apply this identity to the Toeplitz determinant (2.8). First we find the Wiener-Hopf factorization of $`\phi (z)`$:
$$\phi (z)=\phi _+(z)\phi _{}(z)$$
where
$$\phi _+(z)=(1+z)^n,\phi _{}(z)=(1r/z)^m.$$
Define $`K_h`$ acting on $`\mathrm{}^2(\{0,\mathrm{\hspace{0.17em}1},\mathrm{}\})`$ by
$$K_h(j,k)=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}(\phi _{}/\phi _+)_{h+j+\mathrm{}+1}(\phi _+/\phi _{})_{hk\mathrm{}1}.$$
(2.10)
The Borodin-Okounkov identity is then
$$D_h(\phi )=Zdet\left(IK_h\right).$$
Since the determinant on the right tends to 1 as $`h\mathrm{}`$ as does $`\text{Prob}\left(h_t(x)h\right)`$, we have $`Z=(1p)^{mn}`$. Thus we have derived a representation of the distribution function of the random variable $`h_t(x)`$ in terms of a Fredholm determinant,
$$\text{Prob}\left(h_t(x)h\right)=det\left(IK_h\right).$$
(2.11)
This derivation also required $`r<1`$. As above, analytic continuation shows that (2.11) holds for all $`r0`$ if in the integral representing the Fourier coefficients of $`\phi _{}/\phi _+`$ the contour has $`r`$ on the inside and $`1`$ on the outside. (In fact the contour must have $`1`$ on the outside no matter what $`r`$ is.)
A somewhat different direction (and one we do not follow here) is to apply isomonodromy and Riemann-Hilbert methods directly to the Toeplitz determinant $`D_h(\phi )`$. This would result in the identification of $`D_h(\phi )`$ as a $`\tau `$-function of an integrable ODE.
### 2.2 Inhomogeneous ODB
In ODB the probability $`p`$ appearing in Rule (3) is independent of the site $`x`$. Inhomogeneous ODB replaces Rule (3), for each site $`x\text{Z}_+`$, with
$`(3_x)`$ Otherwise, then independently of the other sites and other times, $`h_{t+1}(x)=h_t(x)+1`$ with probability $`0<p_x<1`$ and $`h_{t+1}(x)=h_t(x)`$ with probability $`q_x:=1p_x`$.
Since the dual RSK algorithm is a bijection between $`(0,1)`$-matrices $`A`$ and pairs $`(P,Q)`$ such that $`P^t`$ and $`Q`$ are SSYTs with $`\text{col}(A)=\text{type}(P)`$ and $`\text{row}(A)=\text{type}(Q)`$ , we have
$$\text{Prob}\left(h_t(x)h\right)=q_0^m\mathrm{}q_x^m\underset{\genfrac{}{}{0pt}{}{\lambda 𝒫}{\mathrm{}(\lambda )h}}{}d_\lambda (m)s_\lambda ^{}(r)$$
(2.12)
where, as before, $`m=tx`$, but now $`r=(r_0,\mathrm{},r_x,0,\mathrm{})`$ with $`r_j:=p_j/q_j`$. The proof of (2.12) is straightforward and similar to the proof of the analogous result in ; therefore, we omit it. The right hand side of (2.12) clearly reduces to (2.5) in the homogeneous case.
We again apply Gessel’s theorem to obtain the Toeplitz determinant representation
$$\text{Prob}\left(h_t(x)h\right)=q_1^m\mathrm{}q_n^mD_h(\phi )$$
where<sup>7</sup><sup>7</sup>7The homogeneous case of (2.13) does not directly reduce to (2.9). It does after $`zz/r`$ which corresponds to a similarity transformation of the Toeplitz matrix.
$$\phi (z)=(11/z)^m\underset{j=0}{\overset{x}{}}(1+r_jz).$$
(2.13)
Application of the Borodin-Okounkov identity results in a Fredholm determinant representation for this distribution function. Observe that from either (2.12) or (2.13) it follows that $`\text{Prob}(h_t(x)h)`$ is a symmetric function of $`(p_0,p_1,\mathrm{},p_x)`$. This property opens the possibility for an analysis of the spin glass version of ODB which we plan to address in future work.
### 2.3 Weak ODB and Strict ODB
Here are two natural variants of the ODB. We let the “spontaneous increase” in Rule (3) apply after Rule (2) has already taken effect to get weak ODB:
1. $`h_t(x)h_{t+1}(x)`$ for all space-time points $`(x,t)`$.
2. If $`h_t(x1)>h_t(x)`$, then $`\stackrel{~}{h}_t(x)=h_t(x1)`$ else $`\stackrel{~}{h}_t(x)=h_{t1}(x)`$. (Here $`\stackrel{~}{h}_t`$ is an intermediate height function.)
3. Independently of the other sites and other times, $`h_{t+1}(x)=\stackrel{~}{h}_t(x)+1`$ with probability $`p`$. (With probability $`1p`$, $`h_{t+1}(x)=\stackrel{~}{h}_t(x)`$.)
In strict ODB we require that the left neighbor is rested<sup>8</sup><sup>8</sup>8The height at a site cannot increase at two consecutive times, i.e. it must rest for one time unit before it is allowed to increase. for the spontaneous increase. (We take $`h_t(x)=\mathrm{}`$ for $`x<0`$ which in this model implies $`h_t(0)=0`$ for every $`t`$.)
1. $`h_t(x)h_{t+1}(x)`$ for all space-time points $`(x,t)`$.
2. If $`h_t(x1)>h_t(x)`$, then $`h_{t+1}(x)=h_t(x1)`$.
3. Otherwise, if $`x1`$ is rested at time $`t`$, $`h_t(x1)=h_t(x)`$ then independently of other sites and times, $`h_{t+1}(x)=h_t(x)+1`$ with probability $`p`$ ($`h_{t+1}(x)=h_t(x)`$ with probability $`1p`$.)
In a similar way one shows
* In weak ODB, $`h_t(x)`$ equals, in distribution, the longest sequence $`(i_{\mathrm{}},j_{\mathrm{}})`$ of positions in a random $`(0,1)`$-matrix of size $`m\times n`$ ($`m=tx+1`$, $`n=x+1`$) which have entry $`1`$ such that $`i_{\mathrm{}}`$ are $`j_{\mathrm{}}`$ are both weakly increasing. (The lower left corner of the matrix is fixed to be a 0.)
* In strict ODB, $`h_t(x)`$ equals, in distribution, the longest sequence $`(i_{\mathrm{}},j_{\mathrm{}})`$ of positions in a random $`(0,1)`$-matrix of size $`m\times n`$ ($`m=tx`$, $`n=x`$) which have entry $`1`$ such that $`i_{\mathrm{}}`$ are $`j_{\mathrm{}}`$ are both strictly increasing.
## 3 Limit Theorems
In this section we derive limit theorems for the distribution function $`\text{Prob}\left(h_t(x)h\right)`$ for ODB. Our starting point will be the Fredholm determinant representation (2.11). This distribution function is a function of four variables, $`x`$, $`t`$, $`h`$ and $`p`$; and accordingly, there are several asymptotic regimes:
1. GUE Universal Regime: Let $`x\mathrm{}`$, $`t\mathrm{}`$ such that $`p_c:=1x/t<1`$ is fixed. For fixed $`p<p_c`$ define
$$c_1:=2p_cpp+2\sqrt{pp_c(1p)(1p_c)},$$
(3.1)
$$c_2:=(p_c(1p_c))^{1/6}(p(1p))^{1/2}\left[\left(1+\sqrt{\frac{(1p)(1p_c)}{pp_c}}\right)\left(\sqrt{\frac{p_c}{1p_c}}\sqrt{\frac{p}{1p}}\right)\right]^{2/3}.$$
(3.2)
We will show that
$$\text{Prob}\left(\frac{h_t(x)c_1t}{c_2t^{1/3}}<s\right)F_2(s)$$
where
$$F_2(s)=det\left(IK_{\text{Airy}}\right)=\mathrm{exp}\left(_s^{\mathrm{}}(xs)q(x)^2𝑑x\right).$$
(3.3)
Here $`K_{\text{Airy}}`$ is the operator with Airy kernel acting on $`L^2((s,\mathrm{}))`$ (see (3.5) below) and $`q`$ is the (unique) solution of the Painlevé II equation
$$q^{\prime \prime }=sq+2q^3$$
with boundary condition $`q(s)\text{Ai}(s)`$ as $`s\mathrm{}`$. The limiting shape, $`c_1`$, and the normalization constant, $`c_2`$, as functions of $`x/t`$ are shown in Fig. 4 for $`p=1/2`$. The probability density, $`f_2=dF_2/ds`$, is shown in Fig. 5.
2. Critical Regime: Let $`x\mathrm{}`$, $`t\mathrm{}`$ such that
$$x=(1p)t+o(\sqrt{t}).$$
For fixed $`\mathrm{\Delta }\text{Z}_+`$ we will show that
$$\text{Prob}\left(h_t(x)(tx)\mathrm{\Delta }\right)$$
converges to a $`\mathrm{\Delta }\times \mathrm{\Delta }`$ determinant. One can think of this as
$$p=p_c+o(\frac{1}{\sqrt{t}}).$$
3. Deterministic Regime: For $`x\mathrm{}`$, $`t\mathrm{}`$ and fixed $`p>p_c`$, we will show that
$$\text{Prob}\left(h_t(x)=p_ct\right)1.$$
4. Finite $`x`$ GUE Regime: Fix $`x`$ and let $`t\mathrm{}`$, then we will show that
$$\text{Prob}\left(\frac{h_t(x)pt}{(p(1p)t)^{1/2}}<s\right)$$
converges to the distribution of the largest eigenvalue in the GUE of $`(x+1)\times (x+1)`$ hermitian matrices, denoted below by $`F_{x+1}^{GUE}`$.
### 3.1 GUE Universal Regime
It is convenient to use the variables $`m=tx`$ and $`n=x+1`$ rather than $`x`$ and $`t`$ and to translate back to the space-time variables at the end. We assume $`p<p_c:=m/(n+m)`$. (This is asymptotically $`1x/t`$ as defined above.) Further, when there is no chance of confusion, we denote the random variable $`h_t(x)`$ by $`H`$. (We reserve lower case $`h`$ to denote the values of $`H`$.) Set $`h=cm+sm^{1/3}`$, where $`c`$ will be determined shortly, and $`\alpha =n/m`$. (In this notation the condition $`p<p_c`$ is $`\alpha r<1`$.) For any $`v`$ the matrix $`((v)^{kj}K_h(j,k))`$ has the same Fredholm determinant (the determinant of $`I`$ minus the matrix) as $`(K_h(j,k))`$. We shall show that for a particular $`v`$ and a certain constant $`b>0`$ this matrix scales to a kernel with the same Fredholm determinant as
$$K_{\mathrm{Airy}}(s/v(3b)^{1/3}+x,s/v(3b)^{1/3}+y),$$
(3.4)
on $`(0,\mathrm{})`$, where
$$K_{\mathrm{Airy}}(s+x,s+y)=_0^{\mathrm{}}\mathrm{Ai}(t+s+x)\mathrm{Ai}(t+s+y)𝑑t.$$
(3.5)
This gives
$$\underset{m\mathrm{}}{lim}\mathrm{Prob}\left(\frac{Hcm}{m^{1/3}}s\right)=F_2(s/v(3b)^{1/3}).$$
(3.6)
Here is what we mean by scaling. Any matrix $`(M(j,k))`$ acting on $`\mathrm{}^2(\text{Z}_+)`$ has the same Fredholm determinant as the kernel $`M([x],[y])`$ on $`L^2(0,\mathrm{})`$ and this in turn has the same Fredholm determinant as $`M_m(x,y)=m^{1/3}M([m^{1/3}x],[m^{1/3}y])`$. If this kernel has the limit $`k(x,y)`$ we say that the matrix $`(M(j,k))`$ has, after the scaling $`jm^{1/3}x,km^{1/3}y`$, the limit $`k(x,y)`$. If $`M_m(x,y)`$ converges to $`k(x,y)`$ in trace norm then the Fredholm determinant of $`(M(j,k))`$ converges to that of $`k(x,y)`$. And if $`(M(j,k))`$ were the product of two matrices each having scaling limits in Hilbert-Schmidt norm (under the same scaling, of course), then the Fredholm determinant of the product converges to the Fredholm determinant of the product of the limits. This is what we shall show in our case.
There is a slightly awkward notational problem. Since $`h`$ is always an integer and $`h=cm+sm^{1/3}`$, the quantity $`s`$ as it appears here and the analysis which follows is not completely arbitrary. What we actually show is that if $`h`$ and $`m`$ tend to infinity, and $`s`$ is defined in terms of them by the formula $`h=cm+sm^{1/3}`$, then
$$\mathrm{Prob}\left(Hh\right)F_2\left(s/v(3b)^{1/3}\right)0$$
(3.7)
uniformly for $`s`$ lying in a bounded set. From this we easily deduce (3.6) for fixed $`s`$, which now has a different meaning. These observations are important when one tries to estimate errors. It can be shown that the difference in (3.7) is $`O(m^{2/3})`$. But the difference between the right side of (3.6) and the probability on the left can only be expected to be $`O(m^{1/3})`$. The reason is that if the quantity $`s^{}`$ is defined by $`cm+s^{}m^{1/3}=[cm+sm^{1/3}]`$ then the probability is within $`O(m^{2/3})`$ of $`F_2(s^{}/v(3b)^{1/3})`$, but $`ss^{}`$ is very likely of the order $`m^{1/3}`$.
#### 3.1.1 The Saddle Point Method
The matrix $`(K_h(j,k))`$ is the product of two matrices, the matrix on the right having $`j,k`$ entry $`(\phi _+/\phi _{})_{hjk1}`$ and the one on the left having $`j,k`$ entry $`(\phi _{}/\phi _+)_{h+j+k+1}`$. Notice that the first vanishes if $`h+j+k+1>m`$ so we may assume that all our indices $`j`$ and $`k`$ satisfy $`h+j+k<m`$. We have
$`(\phi _+/\phi _{})_{hjk1}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle (1+z)^n(zr)^mz^{m+h+j+k}𝑑z}`$
$`=`$ $`(1)^{h+j+k}{\displaystyle \frac{1}{2\pi i}}{\displaystyle (1+z)^n(rz)^m(z)^{m+h+j+k}𝑑z},`$
and a similar formula holds for $`(\phi _{}/\phi _+)_{h+j+k+1}`$. If we set
$$\psi (z)=(1+z)^n(rz)^m(z)^{(1c)m}$$
then
$$(1)^{h+j+k}(\phi _+/\phi _{})_{hjk1}=\frac{1}{2\pi i}\psi (z)(z)^{sm^{1/3}+j+k}𝑑z$$
and
$$(1)^{h+j+k}(\phi _{}/\phi _+)_{h+j+k+1}=\frac{1}{2\pi i}\psi (z)^1(z)^{sm^{1/3}jk2}𝑑z.$$
The contours for the first integral surrounds 0 while the contour for the second integral has $`r`$ on the inside and $`1`$ on the outside. The restriction $`h+j+k<m`$ is the same as $`sm^{1/3}+j+k<(1c)m`$. If we make the replacements $`jm^{1/3}x,km^{1/3}y`$ these become
$$\frac{1}{2\pi i}\psi (z)(z)^{m^{1/3}(s+x+y)}𝑑z,\frac{1}{2\pi i}\psi (z)^1(z)^{m^{1/3}(s+x+y+2m^{1/3})}𝑑z.$$
Our restrictions become $`s+x+y<(1c)m^{2/3}`$. For convenience we replace $`s+x+y`$ by $`x`$, and we want to evaluate
$$\frac{1}{2\pi i}\psi (z)(z)^{m^{1/3}x}𝑑z,\frac{1}{2\pi i}\psi (z)^1(z)^{m^{1/3}x2}𝑑z$$
(3.9)
asymptotically. Our restriction is now $`x<(1c)m^{2/3}`$.
To do a steepest descent we have to find the zeros of
$$\frac{d}{dz}\mathrm{log}\psi (z)=\frac{n}{1+z}+\frac{m}{zr}\frac{(1c)m}{z},$$
or equivalently the zeros of
$$(c+\alpha )z^2+(c+rcr\alpha r)z+r(1c).$$
(Recall that $`\alpha =n/m`$.) The discriminant of this quadratic equals zero when
$$c=\frac{1}{1+r}\left(2\sqrt{\alpha r}+(1\alpha )r\right).$$
(3.10)
This is the value of $`c`$ we take.<sup>9</sup><sup>9</sup>9If there were two critical points, or if we took the negative square root in (3.10), the Fredholm determinant would tend exponentially to either zero or one. It is only for this value of $`c`$ that we get a nontrivial limit. The critical probability is the condition $`c=1`$, i.e. $`p_c=m/(m+n)`$. The single zero of the quadratic is then at $`u=v`$ where
$$v=\frac{(1r)c+(1\alpha )r}{2(c+\alpha )}=\frac{1\sqrt{\alpha r}}{1+\sqrt{\alpha /r}}.$$
Note that $`0<v<1`$ since $`0<p<p_c`$. (It is because $`u<0`$ that we used powers of $`z`$ rather than $`z`$ in the definition of $`\psi `$.) We write
$`6b:={\displaystyle \frac{1}{m}}{\displaystyle \frac{d^3}{dz^3}}\mathrm{log}\psi (z)|_{z=u}`$ $`=`$ $`{\displaystyle \frac{2\alpha }{(1+u)^3}}+{\displaystyle \frac{2}{(ur)^3}}{\displaystyle \frac{2(1c)}{u^3}}`$
$`=`$ $`{\displaystyle \frac{2\left(\sqrt{\alpha }+\sqrt{r}\right)^5}{r\sqrt{\alpha }(1+r)^3(1\sqrt{\alpha r})}}.`$
The quantity $`b`$ is positive since $`\alpha r<1`$. In the neighborhood of $`z=u`$,
$$\psi (z)\psi (u)e^{mb(zu)^3}.$$
(3.11)
The steepest descent curves will come into $`u`$ at angles $`\pm \pi /3`$ and $`\pm 2\pi /3`$. Call the former $`C^+`$ and the latter $`C^{}`$. For the integral involving $`\psi (z)`$ we want $`|\psi (z)|`$ to have a maximum at that point of the curve and for the integral involving $`\psi (z)^1`$ we want $`|\psi (z)|`$ to have a minimim there. Since $`b>0`$ the curve for $`\psi (z)`$ must be $`C^+`$ and the curve for $`\psi (z)^1`$ must be $`C^{}`$. Both contours will be described downward near $`u`$. The curve $`C^+`$ will loop around the origin and close at $`r`$, the upper and lower parts making an angle there depending on $`c`$ while $`C^{}`$ will loop around on both sides and go to infinity with slopes depending on $`c`$. (That $`C^\pm `$ have these forms follows from the fact that the contours cannot cross and, since the only critical point is at $`z=u`$, the contours can end only where $`\psi `$, respectively $`\psi ^1`$, is zero.) The steepest descent curves are shown in Fig. 6.
Proceeding formally now, consider the $`\psi (z)`$ integral and make the substitution $`zu+z=v+z`$. Then the old $`z`$ becomes the new $`v(1z/v)ve^{z/v}`$, and recall (3.11). If we make these replacements in the integral we get
$$\psi (u)v^{m^{1/3}x}\frac{1}{2\pi i}_{\mathrm{}e^{i\pi /3}}^{\mathrm{}e^{i\pi /3}}e^{mbz^3m^{1/3}xz/v}𝑑z.$$
The contour can be deformed to the imaginary axis since we only pass through regions where $`\mathrm{}z^3`$ is negative. If we then set $`z=i\zeta /m^{1/3}`$ the above becomes
$$\psi (u)m^{1/3}v^{m^{1/3}x}\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}e^{ib\zeta ^3+ix\zeta /v}𝑑\zeta =\frac{\psi (u)m^{1/3}v^{m^{1/3}x}}{(3b)^{1/3}}\mathrm{Ai}(x/v(3b)^{1/3}).$$
If we recall that $`x`$ was a replacement for $`s+x+y`$ we see that the matrix with $`j,k`$ entry
$$(1)^{h+j+k}\psi (u)^1v^{m^{1/3}sjk}(\phi _+/\phi _{})_{hjk1}$$
has the scaling limit
$$\frac{1}{(3b)^{1/3}}\mathrm{Ai}((s+x+y)/v(3b)^{1/3}).$$
Similarly the matrix with $`j,k`$ entry
$$(1)^{h+j+k}\psi (u)v^{m^{1/3}s+j+k}(\phi _{}/\phi _+)_{h+j+k+1}$$
has $`1/v^2`$ times exactly same scaling limit. Hence the scaling limit of $`u^{kj}K_h(j,k))`$, which has the same Fredholm determinant as $`(K_h(j,k))`$, is the product of these scaling limits,
$$\frac{1}{v^2(3b)^{2/3}}_0^{\mathrm{}}\mathrm{Ai}((s+x+t)/v(3b)^{1/3})\mathrm{Ai}((s+t+y)/v(3b)^{1/3})𝑑t$$
$$=\frac{1}{v(3b)^{1/3}}K_{\mathrm{Airy}}((s+x)/v(3b)^{1/3},(s+y)/v(3b)^{1/3}).$$
And, as promised,<sup>10</sup><sup>10</sup>10The time constant $`c_1=p_cc`$ and the normalization constant $`c_2=p_c^{1/3}v(3b)^{1/3}`$. A computation then gives (3.1) and (3.2). this kernel has the same Fredholm determinant as
$$K_{\mathrm{Airy}}(s/v(3b)^{1/3}+x,s/v(3b)^{1/3}+y).$$
#### 3.1.2 Convergence Proof
Now for the justification. We have to obtain not only the pointwise limit, but uniform estimates to establish convergence of the operators in trace norm. We first obtain asymptotics under the assumption that $`x`$ lies in a bounded set. (Notice that $`xs`$ always.) We begin with
$$\frac{1}{2\pi i}_{C^+}\psi (z)z^{m^{1/3}x}𝑑z,$$
(3.12)
and denote by $`C_\epsilon ^+`$ the portion of $`C^+`$ which lies within $`\epsilon `$ of the critical point $`u_c^{}`$.
Lemma 1. If in (3.12) we integrate only over $`C_\epsilon ^+`$ the error incurred is $`O(|\psi (u)|e^{\delta m})`$ for some $`\delta >0`$.
Proof. Define
$$\sigma (z)=\frac{1}{m}\mathrm{log}\psi (z)=\alpha \mathrm{log}(1+z)+\mathrm{log}(rz)+(c1)\mathrm{log}(z).$$
Its maximum on $`C^+`$ (it is real-valued there) is $`\sigma (u)`$ and it is stricly less than this on the complement of $`C_\epsilon `$ in $`C^+`$. Therefore $`\psi (z)/\psi (u)=O(e^{\delta m})`$ for some $`\delta >0`$ on the complement while $`z^{m^{1/3}x}=e^{O(m^{1/3})}`$. This gives the statement of the lemma.
Lemma 2. We have as $`m\mathrm{}`$
$$m^{1/3}\psi (u)^1v^{m^{1/3}x}\frac{1}{2\pi i}\psi (z)z^{m^{1/3}x}𝑑z=(3b)^{1/3}\mathrm{Ai}(x/v(3b)^{1/3})+O(m^{1/3})$$
uniformly for bounded $`x`$.
Proof. Near $`z=u`$
$$\sigma (z)=\sigma (u)+b(zu)^3+O((zu)^4)),$$
$$\mathrm{log}(z)=\mathrm{log}v+\frac{1}{u}(zu)+O((zu)^2)=\mathrm{log}v\frac{1}{v}(zu)+O((zu)^2).$$
Hence, using Lemma 1, we have
$$\frac{1}{2\pi i}_{C^+}\psi (z)z^{m^{1/3}x}𝑑z$$
$$=O(|\psi (u)|e^{\delta m})+\psi (u)v^{m^{1/3}x}\frac{1}{2\pi i}_{C_\epsilon ^+}e^{mb(zu)^3m^{1/3}x(zu)/v+O(m(zu)^4+m^{1/3}x(zu)^2)}𝑑z.$$
We show that removing the $`O`$ term in the exponential in the integrand leads to an error $`O(m^{2/3})`$ in the integral. This error equals
$$_{C_\epsilon ^+}e^{mb(zu)^3m^{1/3}x(zu)/v}\left(e^{O(m(zu)^4+m^{1/3}x(zu)^2)}1\right)𝑑z$$
$$=_{C_\epsilon ^+}e^{mb(zu)^3m^{1/3}x(zu)/v+O(m(zu)^4+m^{1/3}x(zu)^2)}O(m(zu)^4+m^{1/3}x(zu)^2)𝑑z.$$
Now the exponential has the form
$$e^{mb(1+\eta _1)(zu)^3m^{1/3}x(1+\eta _2)(zu)/v},$$
where the $`\eta _i`$ can be made arbitrarily small by taking $`\epsilon `$ small enough. If we make the substitution $`zu=\zeta /m^{1/3}`$ the error becomes
$$m^{2/3}e^{b(1+\eta _1)\zeta ^3x(1+\eta _2)\zeta /v}O(\zeta ^4+x\zeta ^2)𝑑\zeta .$$
The integral is now taken over a long contour lying in thin angles around the rays $`|\mathrm{arg}\zeta |=\pi /3`$, with ends having absolute value at least a constant times $`m`$. This integral is clearly bounded, uniformly in $`m`$ for bounded $`x`$.
Therefore with the stated error we may remove the $`O`$ terms from the exponential in the original integral. Then we make the same substitution. The integrand is exponentially small at the ends of the resulting contour. Therefore if we complete it so that it goes to infinity in the two directions $`\pm \pi /3`$ the error incurred will be exponentially small.
We have shown that
$$m^{1/3}\psi (u)^1v^{m^{1/3}x}\frac{1}{2\pi i}\psi (z)z^{m^{1/3}x}𝑑z=\frac{1}{2\pi i}_{\mathrm{}e^{i\pi /3}}^{\mathrm{}e^{i\pi /3}}e^{b\zeta ^3x\zeta /v}𝑑\zeta +O(m^{1/3}).$$
If we deform the contour to the imaginary axis and make the substitution $`\zeta i\zeta `$ then the last integral, with its factor, becomes
$$\frac{1}{2\pi }_{\mathrm{}}^{\mathrm{}}e^{ib\zeta ^3+ix\zeta /v}𝑑\zeta =(3b)^{1/3}\mathrm{Ai}(x/v(3b)^{1/3}).$$
This proves the lemma.
The second integral in (3.9) is similar.
Lemma 3. We have as $`m\mathrm{}`$
$$m^{1/3}\psi (u)v^{m^{1/3}x}\frac{1}{2\pi i}_C^{}\psi (z)^1z^{m^{1/3}x}𝑑z=(3b)^{1/3}\mathrm{Ai}(x/v(3b)^{1/3})+O(m^{1/3})$$
uniformly for bounded $`x`$.
Proof. The derivation is essentially the same. The exponentials are replaced by their negatives and the directions $`\pm \pi /3`$ are replaced by $`\pm 2\pi /3`$. The fact that $`C^{}`$ is unbounded causes no difficulty since the integrand now behaves at infinity like a large negative power of $`z`$. We get the same Airy function in the end, as we have already seen.
Now for the tricky part. We need estimates that are uniform for all $`x`$ and where the error term contains a factor which is very small for large $`x`$. In fact we shall show that the statements of Lemmas 2 and 3 hold, uniformly for all $`x`$, when the error terms are replaced by $`m^{1/3}e^x`$. (The $`x`$ in the exponential can be improved to a constant times $`x^{3/2}`$ but that makes no difference.) To do this we have to be more careful and use the steepest descent curves for the full integrands in (3.9), not just for the factors $`\psi ^{\pm 1}`$. We consider in detail only the first integral in (3.9); as before, the second is treated analogously.
Set
$$\psi (z,c^{})=(1+z)^n(rz)^m(z)^{(1c^{})m}=\psi (z)(z)^{(c^{}c)m}.$$
We are interested in the asymptotics of
$$I(c^{})=\frac{1}{2\pi i}\psi (z,c^{})𝑑z$$
(3.13)
when $`c^{}c=m^{2/3}x`$. Our condition on $`x`$ says that $`c^{}<1`$ and, in view of what we already know, we may assume $`x`$ is positive and bounded away from zero, so $`c^{}>c`$.
We let $`C`$ be the steepest descent curve for $`\psi (z,c^{})`$. This curve now passes vertically through one of the critical points of $`\psi (z,c^{})`$. For $`c^{}>c`$ there are two critical points
$$u_c^{}^\pm =\frac{(1r)c^{}(1\alpha )r\pm \sqrt{\left((1+r)c^{}+(\alpha 1)r\right)^24\alpha r}}{2(\alpha +c^{})},$$
which are real and satisfy
$$1<u_c^{}^{}<v<u_c^{}^+<0.$$
To determine which critical point our curve passes through we consider the function
$$\sigma (z,c^{})=\frac{1}{m}\mathrm{log}\psi (z,c^{})=\sigma (z)+(c^{}c)\mathrm{log}(z).$$
The critical points $`u_c^{}^\pm `$ are the zeros of $`\sigma _z(u_c^{},c^{})`$. (Subscripts here and below denote derivatives in the usual way.) We use the fact that $`u_c^{}^\pm `$ are smooth functions of $`\gamma =\sqrt{c^{}c}`$ and compute, recalling that $`\sigma _z(u,c)=\sigma _{zz}(u,c)=0`$ and observing that $`dc^{}/d\gamma =0`$ when $`\gamma =0`$,
$$\frac{d}{d\gamma }\sigma _{zz}(u_c^{}^\pm ,c^{})|_{\gamma =0}=\sigma _{zzz}(u,c)\frac{du_c^{}^\pm }{d\gamma }|_{\gamma =0}.$$
(3.14)
The first factor on the right is positive (we denoted it by $`6b`$), while
$$\frac{du_c^{}^\pm }{d\gamma }|_{\gamma =0}=\pm \beta $$
(3.15)
where
$$\beta =\frac{(\alpha r)^{1/4}(1+r)^{3/2}}{(\sqrt{\alpha }+\sqrt{r})^2}.$$
Since $`\sigma _{zz}(u_c^{}^\pm ,c^{})=0`$ when $`\gamma =0`$ we deduce that for $`c^{}`$ close to, but greater than, $`c`$ we have
$$\sigma _{zz}(u_c^{}^+,c^{})>0,\sigma _{zz}(u_c^{}^{},c^{})<0.$$
These inequalities hold for all $`c^{}`$ since the second derivative can be zero only when $`c^{}=c`$. This shows that the steepest descent curve $`C`$ for $`\psi (z,c^{})`$ passes through $`u_c^{}^+`$, because on the curve $`|\psi (z,c^{})|`$ has a maximum at the critical point. (Similarly the steepest descent curve for $`\psi (z,c^{})^1`$ passes through $`u_c^{}^{}`$.) To make the notation less awkward we write $`u_c^{}`$ instead of $`u_c^{}^+`$. First, we have the analogues of Lemmas 2 and 3.
Lemma 4. Given $`\epsilon >0`$ there exists a $`\delta >0`$ such that $`I(c^{})=O(|\psi (u,c^{})|e^{\delta m})`$ if $`c^{}c>\epsilon `$.
Proof. The function $`\sigma (z,c^{})`$ is decreasing for $`u<z<u_c^{}`$ since it decreases near and to the left of $`u_c^{}`$ and has no critical point in this interval. Hence $`\sigma (u_c^{},c^{})<\sigma (u,c^{})`$ so $`\sigma (u_c^{},c^{})\sigma (u,c^{})`$ is negative and bounded away from zero for $`c^{}>c+\epsilon `$. Since the maximum of $`|\psi (z,c^{})|`$ on $`C`$ is at $`z=u_c^{}`$ the statement follows.
In view of Lemma 4 we may assume in what follows that $`c^{}c`$ is as small as we please. We denote by $`C_\epsilon `$ the portion of $`C`$ which lies within $`\epsilon `$ of the critical point $`u_c^{}`$.
Lemma 5. If in the integral (3.13), in which we integrate over $`C`$, we integrate only over $`C_\epsilon `$ the error incurred is $`O(|\psi (u,c^{})|e^{\delta m})`$ for some $`\delta >0`$.
Proof. The maximum of $`|\psi (z,c^{})|^{1/m}`$ on $`C`$ occurs at $`u_c^{}`$ and it is stricly smaller on the complement of $`C_\epsilon `$ in $`C`$ than it is at $`u_c^{}`$. Therefore the integral in question is $`O(|\psi (u_c^{},c^{})|e^{\delta m})`$ for some $`\delta >0`$. Since $`\sigma (u_c^{},c^{})<\sigma (u,c^{})`$, as we saw in the proof of the last lemma, this one is established.
Because of Lemmas 4 and 5 we need only compute the behavior of $`\psi (z,c^{})`$, or equivalently $`\sigma (z,c^{})`$, for $`z`$ near $`u_c^{}`$. Recall that $`u=u_c=v`$.
Lemma 6. We have
$`(i)\sigma (u_c^{},c^{})`$ $`=`$ $`\sigma (u)+(c^{}c)\mathrm{log}v{\displaystyle \frac{2}{3}}{\displaystyle \frac{\beta }{v}}(c^{}c)^{3/2}+O\left((c^{}c)^2\right);`$
$`(ii)\sigma _z(u_c^{},c^{})`$ $`=`$ $`0;`$
$`(iii)\sigma _{zz}(u_c^{},c^{})`$ $`=`$ $`6b\beta \sqrt{c^{}c}+O(c^{}c);`$
$`(iv)\sigma _{zzz}(u_c^{},c^{})`$ $`=`$ $`6b+O(\sqrt{c^{}c}).`$
Proof. From (3.15) and the fact that $`u_c^{}`$ is a smooth function of $`\gamma `$ (or directly) we see that
$$u_c^{}u=\beta \sqrt{c^{}c}+O(c^{}c).$$
(3.16)
Consequently, since $`u=v`$,
$$\frac{u_c^{}}{u}=1\frac{\beta }{v}\sqrt{c^{}c}++O(c^{}c).$$
(3.17)
Now since $`\sigma _z(u_c^{},c^{})=0`$ we have
$$\frac{d}{dc^{}}\sigma (u_c^{},c^{})=\frac{}{c^{}}\sigma (z,c^{})|_{z=u_c^{}}=\mathrm{log}(u_c^{})=\mathrm{log}v+\mathrm{log}\frac{u_c^{}}{u}.$$
Integrating with respect to $`c^{}`$ from $`c`$ to $`c^{}`$ and using (3.17) we obtain (i). Of course (ii) is immediate. As for (iii) and (iv), these follow from (3.14) and (3.15) and the fact that $`u_c^{}`$ is a smooth functions of $`\gamma `$.
Lemma 7. The conclusions of Lemmas 2 and 3 hold uniformly for all $`x`$ when the error terms are replaced by $`O(e^{\delta m})+O(m^{1/3}e^x)`$ for some $`\delta >0`$.
Proof. We consider (3.12), which is $`I(c^{})`$ with $`c^{}c=m^{2/3}x`$. Putting together Lemmas 5 and 6 we deduce that
$$I(c^{})=O(|\psi (u)v^{(c^{}c)m}|e^{\delta m})+\psi (u)v^{(c^{}c)m}e^{\frac{2\beta }{3v}(c^{}c)^{3/2}m}\times $$
$$\frac{1}{2\pi i}_{C_\epsilon }e^{mb(zu_c^{})^3+3mb\beta \sqrt{c^{}c}(zu_c^{})^2+O(m[(c^{}c)^2+(c^{}c)|zu_c^{}|^2+\sqrt{c^{}c}|zu_c^{}|^3+|zu_c^{}|^4])}𝑑z.$$
If $`c^{}c=m^{2/3}x`$ the exponential factor equals $`e^{\frac{2\beta }{3v}x^{3/2}}`$ while the integral equals
$$\frac{1}{2\pi i}_{C_\epsilon }e^{mb(zu_c^{})^3+3m^{2/3}bx^{1/2}\beta (zu_c^{})^2+O(m^{1/3}x^2+m^{1/3}x|zu_c^{}|^2+m^{2/3}x^{1/2}|zu_c^{}|^3+m|zu_c^{}|^4])}𝑑z.$$
Now $`C_\epsilon `$, rather than looking like two rays near the critical point, looks like one branch of a hyperbola.
Note that by Lemma 4 we may assume that $`c^{}c=m^{2/3}x`$ is as small as desired. It follows that the exponent, without the $`O(m^{1/3}x^2)`$ term, can be written
$$m(b+\eta _1)(zu_c^{})^3+3m^{2/3}(b+\eta _2)x^{1/2}\beta (zu_c^{})^2,$$
where, if $`\epsilon `$ is chosen small enough, the $`\eta _i`$ can be made as small as desired. Upon making the variable change $`zu_c^{}=\zeta /m^{1/3}`$ the integral becomes
$$\frac{m^{1/3}}{2\pi i}e^{(b+\eta _1)\zeta ^3+3(b+\eta _2)x^{1/2}\beta \zeta ^2}𝑑\zeta ,$$
taken over a long contour in the right half-plane on which $`|\mathrm{arg}\zeta |>\pi /3\eta `$, with another small $`\eta `$. The integral here is uniformly bounded.
To take care of the term $`O(m^{1/3}x^2)`$ in the exponential in the original integral, observe that if $`m^{2/3}x`$ is small enough then $`m^{1/3}x^2`$ will be at most a small constant times $`x^{3/2}`$ and so
$$e^{\frac{2\beta }{3v}x^{3/2}}\left(e^{O(m^{1/3}x^2)}1\right)=O(m^{1/3}x^2e^{\frac{\beta }{2v}x^{3/2}})=O(m^{1/3}e^x).$$
Thus removing the term from the exponential leads to an eventual error $`O(m^{2/3}e^x)`$. That removing the other $`O`$ terms from the exponential leads to the same error is seen as it was in the proof of Lemma 2—the substitution in the integral representating the error results in an extra factor $`m^{1/3}`$ and there is the exponential factor $`e^{\frac{2\beta }{3v}x^{3/2}}`$ outside the integral.
After removing all the $`O`$ terms and making the variable change $`zu_c^{}=\zeta /m^{1/3}`$ the integral becomes
$$\frac{m^{1/3}}{2\pi i}e^{b\zeta ^3+3bx^{1/2}\beta \zeta ^2}𝑑\zeta ,$$
taken over a long contour in the right half-plane on which $`|\mathrm{arg}\zeta |>\pi /3\eta `$. Completing the contour so that it goes to infinity in the directions $`\mathrm{arg}\zeta =\pm \pi /3`$ leads to an exponentially small error. It follows that (the first part of) the lemma holds with the negative of the Airy function in the statement replaced by
$$e^{\frac{2\beta }{3v}x^{3/2}}\frac{1}{2\pi i}_{\mathrm{}e^{i\pi /3}}^{\mathrm{}e^{i\pi /3}}e^{b\zeta ^3+3bx^{1/2}\beta \zeta ^2}𝑑\zeta .$$
If we complete the cube and make the substitution $`\zeta \zeta \beta x^{1/2}`$ this becomes, upon noting that $`3b\beta ^2=1/v`$,
$$_{\mathrm{}e^{i\pi /3}}^{\mathrm{}e^{i\pi /3}}e^{b\zeta ^3x\zeta /v}𝑑\zeta =(3b)^{1/3}\mathrm{Ai}(x/(v(3b)^{1/3})).$$
The second part of the lemma is analogous, just as the proof of Lemma 3 was analogous to the proof of Lemma 2.
We have now shown that if we set $`j=m^{1/3}x,k=m^{1/3}y`$ then
$$(1)^{h+j+k}m^{1/3}\psi (u)v^{m^{1/3}s+j+k}(\phi _+/\phi _{})_{hjk1}(3b)^{1/3}\mathrm{Ai}((s+x+y)/v(3b)^{1/3}),$$
and the difference between the two is $`O(m^{1/3}e^{(x+y)})+O(e^{\delta m})`$. It follows easily from this that if we denote the matrix on the left, without the factor $`m^{1/3}`$, by $`(M(j,k))`$ and the kernel on the right by $`A(x,y)`$ then the kernel $`m^{1/3}M([m^{1/3}x],[m^{1/3}y])`$ converges in Hilbert-Schmidt norm to the kernel $`A(x,y)`$ on $`(0,\mathrm{})`$. (Recall that $`j`$ and $`k`$ are at most $`O(m)`$. Therefore the error term $`O(e^{\delta m})`$ can only contribute an exponentially small error to the norm and so can be ignored. Similarly we can let our indices $`j`$ and $`k`$ run to infinity.) Thus, under the scaling $`jm^{1/3}x,km^{1/3}y`$ the matrices with $`j,k`$ entry
$$(1)^{h+j+k}\psi (u)v^{m^{1/3}s+j+k}(\phi _+/\phi _{})_{hjk1}$$
scale in Hilbert-Schmidt norm to the kernel $`A(x,y)`$. Similarly so do the matrices with $`j,k`$ entry
$$(1)^{h+j+k}\psi (u)^1v^{(n^{1/3}s+j+k)}(\phi _{}/\phi _+)_{h+j+k+1}.$$
Therefore the product of the matrices scale in trace norm to the (operator) square of the kernel, which is the Airy kernel (3.4). This or completeness the justification.
### 3.2 Critical Regime: $`pp_c`$
When $`p=p_c`$ ($`\alpha r=1`$),<sup>11</sup><sup>11</sup>11See the remark at end of this section. the analysis of the previous section must be modified. We set $`h=m\mathrm{\Delta }h`$ ($`\mathrm{\Delta }h=0,1,2,\mathrm{}`$) and introduce the new $`\psi `$
$$\psi =(1+z)^n(zr)^m$$
and the corresponding new $`\sigma `$
$$\sigma (z)=\frac{1}{m}\mathrm{log}\psi =\alpha \mathrm{log}(1+z)+\mathrm{log}(zr).$$
The saddle point now occurs at $`z=0`$ with $`\sigma ^{\prime \prime }(0)=\alpha (1+\alpha )`$. Thus in the neighborhood of $`z=0`$
$$\psi (z)(1)^mr^me^{m\alpha (1+\alpha )z^2/2}.$$
(3.18)
Since $`(\phi _+/\phi _{})_{hkj1}`$ vanishes for $`h+j+k+1>m`$, we can again assume $`h+j+k<m`$ which becomes the condition $`j+k<\mathrm{\Delta }h`$. As before our starting point is the integral expression
$$(\phi _+/\phi _{})_{hjk1}=\frac{1}{2\pi i}\psi (z)z^{m+h+j+k}𝑑z$$
where the contour is a circle centered at 0 with radius $`\rho <1`$. Taking this $`\rho `$ sufficiently small so that we may use the approximation (3.18) on the integrand, we obtain after making the change of variables
$$\zeta =\left(\frac{m\alpha (1+\alpha )}{2}\right)^{1/2}z=z/S,$$
$`(\phi _+/\phi _{})_{hjk1}`$ $``$ $`(1)^mr^mS^{j+k\mathrm{\Delta }h+1}{\displaystyle \frac{1}{2\pi i}}{\displaystyle e^{\zeta ^2}\zeta ^{j+k\mathrm{\Delta }h}𝑑\zeta }`$
$`=`$ $`\{\begin{array}{cc}(1)^mr^mS^{j+k\mathrm{\Delta }h+1}\frac{(1)^L}{L!}& \text{if}\mathrm{\Delta }hjk1=2L=0,2,4\mathrm{}\\ 0& \text{if}\mathrm{\Delta }hjk1=\text{odd integer}.\end{array}`$
Our second integral is
$$(\phi _{}/\phi _+)_{h+j+k+1}=\frac{1}{2\pi i}\psi (z)^1z^{mhjk2}𝑑z$$
where the contour has $`1`$ on the outside and $`r`$ on the inside. We deform the contour to the imaginary axis going from $`i\mathrm{}`$ to $`i\mathrm{}`$ with an infinitesimal indentation going around 0 to the left. The part of the contour lying in the right half plane is exponentially small because of the factor $`(1+z)^n`$ and can therefore be neglected. For the integral along the imaginary axis we can replace $`\psi `$ by (3.18) with an error that is exponentially small. Thus the above integral is asymptotically equal to
$$(1)^m\frac{r^m}{2\pi i}_i\mathrm{}^i\mathrm{}e^{z^2/S^2}z^{\mathrm{\Delta }hjk2}𝑑z,$$
which in turn equals
$$(1)^mi^{\mathrm{\Delta }hjk1}r^mS^{\mathrm{\Delta }hjk1}\frac{1}{2\pi i}_{\mathrm{}}^{\mathrm{}}e^{\zeta ^2}\zeta ^{\mathrm{\Delta }hjk2}𝑑\zeta $$
where there is an indentation above $`\zeta =0`$. If we now substitute $`\zeta =\sqrt{t}`$, the above integral becomes
$$(1)^mi^{\mathrm{\Delta }hjk1}r^mS^{\mathrm{\Delta }hjk1}\frac{1}{4\pi i}_{\mathrm{}}^{0^+}e^tt^{(\mathrm{\Delta }hjk1)/21}𝑑t.$$
The contour starts at $`+\mathrm{}`$, loops around 0 in the positive direction and then returns to $`+\mathrm{}`$. This last integral is Hankel’s integral representation of the $`\mathrm{\Gamma }`$ function. Thus
$$(\phi _{}/\phi _+)_{h+j+k+1}\frac{(1)^{h+j+k+1}}{2\pi }r^mS^{\mathrm{\Delta }hjk1}\mathrm{sin}\left(\frac{\pi }{2}(\mathrm{\Delta }hjk1)\right)\mathrm{\Gamma }\left(\frac{\mathrm{\Delta }hjk1}{2}\right).$$
We now use these two asymptotic expressions along with the condition $`j+k<\mathrm{\Delta }h`$ in (2.10) to obtain (after a short calculation)
$$K_h(j,k)\frac{(S)^{kj}}{2\pi }\underset{\mathrm{}=0}{\overset{[\frac{\mathrm{\Delta }hk1}{2}]}{}}\frac{1}{\mathrm{}!}\mathrm{sin}\frac{\pi }{2}(kj)\mathrm{\Gamma }\left(\mathrm{}+\frac{kj}{2}\right).$$
(3.20)
When $`\mathrm{}+(kj)/2`$ is a nonpositive integer, the product of the sine and gamma functions is replaced by
$$\frac{(1)^{\mathrm{}}\pi }{\left(\frac{jk}{2}\mathrm{}\right)!}.$$
The factor $`(S)^{kj}`$ may be dropped when computing the determinant $`det(IK_h)`$ since it does not change its value. We evaluate this determinant and display the results for $`\mathrm{\Delta }h9`$ in Table 1.
Remark. Since $`m`$ and $`n`$ are integers it is extremely unlikely that $`p=p_c=m/(m+n)`$. If $`p`$ is irrational this never occurs. However the preceding analysis shows that if $`\alpha r=1+o(m^{1/2})`$ rather than 1 then in the integrals one gets extra factors $`(1+z)^{o(m^{1/2})}`$. Then after the substitution $`z=S\zeta `$ this drops out since $`S=O(m^{1/2})`$. The upshot is that the asymptotics hold for any $`p`$ when $`m`$ and $`n`$ go to infinity in such a way that $`m/(m+n)=p+o(m^{1/2})`$.
### 3.3 Deterministic Regime: $`p>p_c`$
#### 3.3.1 Large Deviations Approach
Assume that $`p>p_c`$. Then there exists an $`ϵ>0`$ so that $`n/m`$ approaches $`(1+ϵ)(1/p1)`$. To simplify the statements, we will just assume that $`n=(1+ϵ)(1/p1)m`$.
Imagine the random $`m\times n`$ matrix $`A`$ from §2.1 as the lower left corner of an infinite matrix of 0’s and 1’s, created by the independent coin flips. Fix a position $`(i,j)`$ ($`i,j1`$) in this infinite random matrix. Define $`J`$ as the column index of the first entry, from left to right, with a 1 on the row above $`(i,j)`$ and in the columns larger or equal $`j`$. Then define $`\xi _{(i,j)}=Jj`$. In the example given, $`\xi _{(3,1)}=0`$ and $`\xi _{(5,1)}=3`$.
Now create a sequence of i.i.d. random variables $`\xi _1,\xi _2,\mathrm{}`$, as follows. Let $`\xi _1`$ equal the column index minus one of the first 1 on the first row. Then let $`\xi _2=\xi _{(1,1+\xi _1)}`$, $`\xi _3=\xi _{(2,1+\xi _1+\xi _2)},\mathrm{}`$. The basic observation is that, since we are always taking the best positioned 1 on the next line, we have equality of the two events
$$\{\text{there is an increasing path of length }m\text{ in }A\}=\{\xi _1+\mathrm{}+\xi _m<n\}.$$
Therefore, we need to show that
$$\text{Prob}(\xi _1+\mathrm{}+\xi _mn)$$
goes to 0 exponentially as $`m\mathrm{}`$. However, $`\text{Prob}(\xi _1=i)=p(1p)^i`$, $`i=0,1,\mathrm{}`$ and so $`E(\xi _1)=1/p1`$. By elementary large deviations (e.g. §1.9 in ),
$$m^1\mathrm{log}P(\xi _1+\mathrm{}+\xi _mn)\gamma (ϵ)$$
where an elementary calculation shows
$$\gamma (ϵ)=(1/p1)(1+ϵ)\mathrm{log}(1+ϵ)p^1(1+ϵϵp)\mathrm{log}(1+ϵϵp)$$
which is positive whenever $`ϵ>0`$.
#### 3.3.2 Saddle Point Approach
For completeness, we show how the saddle point method gives the same result. Thus we show that when $`p>p_c`$ (or $`\alpha r>1`$)
$$det(IK_h)0$$
exponentially as $`m\mathrm{}`$ even when $`h=m1`$, thus establishing assertion 1(c) in §3 with exponential approach to the limit.
As we saw at the beginning in the last section we need only consider the entries $`K_h(j,k)`$ when $`h+j+k<m`$, which in the present situation means $`j=k=0`$. Our claim is therefore that $`K_{m1}(0,0)1`$ exponentially as $`m\mathrm{}`$. The first integral to consider is
$$(\phi _+/\phi _{})_m=\frac{1}{2\pi i}\psi (z)z^1𝑑z=\psi (0)=(r)^m.$$
The second integral is
$$(\phi _{}/\phi _+)_m=\frac{1}{2\pi i}\psi (z)^1z^1𝑑z.$$
Recall that the contour here surrounds 0 and has $`1`$ on the outside, $`r`$ on the inside. The critical point for steepest descent is at $`z=u`$ where
$$\frac{\alpha }{1+u}+\frac{1}{ur}=0,u=\frac{\alpha r1}{\alpha +1}.$$
The steepest descent curve will pass vertically through this point and go to $`\mathrm{}`$ in two directions. But notice that since $`u`$ is positive, in order to deform our original contour to this one we have to pass through $`z=0`$. The residue of the integrand there equals $`(r)^m`$ and so
$$(\phi _{}/\phi _+)_m=(r)^m+\frac{1}{2\pi i}\psi (z)^1z^1𝑑z,$$
where now the integral is taken over the steepest descent curve. This integral is asymptotically a constant times $`m^{1/2}`$ times the value of the integrand at $`z=u`$, and this value equals $`(1)^m`$ times
$$\left(\frac{\alpha (r+1)}{\alpha +1}\right)^{\alpha m}\left(\frac{r+1}{\alpha +1}\right)^m.$$
Our claim is therefore equivalent to the statement that this is exponentially smaller than $`r^m`$, which in turn is equivalent to the inequality
$$(r+1)^{\alpha +1}\frac{\alpha ^\alpha }{(\alpha +1)^{\alpha +1}}>r.$$
It is an elementary exercise that this is true for all $`r0`$ except for $`r=1/\alpha `$, when equality holds. But in our case $`r>1/\alpha `$ so the inequality holds.
### 3.4 Finite GUE Regime: Fixed $`x`$ and $`t\mathrm{}`$
#### 3.4.1 Saddle Point Calculation
We return to (3.1.1) and this time set
$$h=\frac{r}{1+r}m+sm^{1/2}=pm+sm^{1/2},$$
and make the substitutions $`jxm^{1/2},kym^{1/2}`$ to write the integral (3.1.1) as
$$\frac{1}{2\pi i}(1+z)^n(rz)^m(z)^{m/(1+r)}(z)^{(s+x+y)m^{1/2}}𝑑z.$$
(3.21)
Now we set
$$\psi (z)=(rz)^m(z)^{m/(1+r)},$$
which is the main part of the integrand. There is a single critical point, $`z=1`$, and at this point $`d^2/dz^2\mathrm{log}\psi (z)`$ is equal to
$$m\frac{r}{(1+r)^2}=mp(1p).$$
This is positive and so the steepest descent curve is vertical at the critical point; it goes around the origin and closes at $`z=r`$. The main contribution to the integral comes from the immediate neighborhood of the critical point. If we make the variable change
$$z=1+\frac{\zeta }{\sqrt{m}}$$
and take into account the other factors in the integrand we see the integral is asymptotically
$$\frac{(r+1)^m}{2\pi i}_i\mathrm{}^i\mathrm{}\left(\frac{\zeta }{\sqrt{m}}\right)^ne^{\frac{1}{2}p(1p)\zeta ^2(s+x+y)\zeta }\frac{d\zeta }{\sqrt{m}}.$$
(3.22)
Now
$$\frac{1}{2\pi i}_i\mathrm{}^i\mathrm{}e^{a\zeta ^2b\zeta }𝑑\zeta =\frac{e^{b^2/4a}}{2\sqrt{a\pi }}$$
and so
$$\frac{1}{2\pi i}_i\mathrm{}^i\mathrm{}\zeta ^ne^{a\zeta ^2b\zeta }𝑑\zeta =\frac{(1)^n}{2\sqrt{a\pi }}\frac{d^n}{db^n}e^{b^2/4a}=\frac{1}{\sqrt{\pi }(2\sqrt{a})^{n+1}}e^{b^2/4a}H_n\left(\frac{b}{2\sqrt{a}}\right).$$
($`H_n`$ are the Hermite polynomials.) Hence our first integral (3.21) is asymptotically equal to $`(r+1)^m/\sqrt{m}^{n+1}`$ times this expression with
$$a=\frac{1}{2}p(1p),b=s+x+y.$$
Thus we have shown that the matrix with $`j,k`$ entry
$$(1)^{hjk}(r+1)^mm^{n/2}(\phi _+/\phi _{})_{hjk1}$$
scales to the operator on $`(0,\mathrm{})`$ with kernel
$$\frac{1}{\sqrt{\pi }(2\sqrt{a})^{n+1}}e^{(s+x+y)^2/4a}H_n\left(\frac{s+x+y}{2\sqrt{a}}\right),$$
with $`a`$ as given above.
Next, with the same substitutions in the integral,
$$(\phi _{}/\phi _+)_{h+j+k+1}=\frac{1}{2\pi i}(1+z)^n(zr)^mz^{mhjk2}𝑑z$$
$$=(1)^{hjk}\frac{1}{2\pi i}(1+z)^n(rz)^m(z)^{m/(1+r)}(z)^{(s+x+y)m^{1/2}2}𝑑z.$$
The contour here encloses $`0`$ and $`r`$ and has $`1`$ on the outside. The steepest descent curve should go through the critical point $`1`$ horizontally. We deform the given contour to a curve starting at $`\mathrm{}+0i`$, going above the the real axis, looping around $`z=1`$ clockwise, then back below the real axis to $`\mathrm{}0i`$. The original contour can be deformed to this because the integrand is small at $`\mathrm{}`$. The main contribution is again in the neighborhood of $`z=1`$. Making the same variable change as before leads to an integral which is asymptotically
$$\frac{(r+1)^m}{2\pi i}\left(\frac{\zeta }{\sqrt{m}}\right)^ne^{\frac{1}{2}p(1p)\zeta ^2+(s+x+y)\zeta }\frac{d\zeta }{\sqrt{m}},$$
(3.23)
where now the contour is a circle going around $`\zeta =0`$ counterclockwise. Using now the fact
$$\frac{1}{2\pi i}\zeta ^ne^{a\zeta ^2+b\zeta }𝑑\zeta =\frac{a^{(n1)/2}}{(n1)!}e^{b^2/4a}\frac{d^{n1}}{d\zeta ^{n1}}e^{(\zeta b/2\sqrt{a})^2}|_{\zeta =0}$$
$$=\frac{a^{(n1)/2}}{(n1)!}H_{n1}\left(\frac{b}{2\sqrt{a}}\right)$$
we find that the matrix with $`j,k`$ entry
$$(1)^{h+j+k}(r+1)^mm^{n/2}(\phi _{}/\phi _+)_{h+j+k+1}$$
scales to the operator on $`(0,\mathrm{})`$ with kernel
$$\frac{a^{(n1)/2}}{(n1)!}H_{n1}\left(\frac{s+x+y}{2\sqrt{a}}\right).$$
Combining, we see that the product of the two matrices (aside from a factor $`(1)^{jk}`$, which does not affect the determinant) has scaling limit the operator with kernel
$$\frac{1}{\sqrt{\pi }2^{n+1}a(n1)!}_0^{\mathrm{}}e^{(s+x+z)^2/4a}H_n\left(\frac{s+x+z}{2\sqrt{a}}\right)H_{n1}\left(\frac{s+z+y}{2\sqrt{a}}\right)𝑑z.$$
Instead of a direct evaluation of this last integral, we will not evaluate our $`\zeta `$ integrals (3.22) and (3.23), but rather consider them as integrals with variables $`\zeta _1`$ and $`\zeta _2`$, combine and integrate with respect to $`z`$. We see that the scaled kernel for the product is
$$\frac{1}{4\pi ^2}_0^{\mathrm{}}\left(\frac{\zeta _1}{\zeta _2}\right)^ne^{a(\zeta _1^2\zeta _2^2)(s+x+z)\zeta _1+(s+z+y)\zeta _2}𝑑z𝑑\zeta _1𝑑\zeta _2,$$
where the $`\zeta _1`$ contour is a vertical line described upward and the $`\zeta _2`$ contour goes around 0 counterclockwise. If the vertical line is to the right of the circle we can integrate first with respect to $`z`$, yielding
$$\frac{1}{4\pi ^2}\left(\frac{\zeta _1}{\zeta _2}\right)^ne^{a(\zeta _1^2\zeta _2^2)(s+x)\zeta _1+(s+y)\zeta _2}\frac{d\zeta _1d\zeta _2}{\zeta _1\zeta _2}.$$
Let’s call this $`L_n(x,y)`$. This is 0 when $`n=0`$, and
$$L_k(x,y)L_{k1}(x,y)=\frac{1}{4\pi ^2}\frac{\zeta _1^{k1}}{\zeta _2^k}e^{a(\zeta _1^2\zeta _2^2)(s+x)\zeta _1+(s+y)\zeta _2}𝑑\zeta _1𝑑\zeta _2.$$
This integral is a product and we can use the computations we did above to see that it equals
$$\frac{1}{2^k\sqrt{\pi a}(k1)!}e^{(s+x)^2/4a}H_{k1}\left(\frac{s+x}{2\sqrt{a}}\right)H_{k1}\left(\frac{s+y}{2\sqrt{a}}\right).$$
If $`\phi _k`$ are the oscillator wave functions<sup>12</sup><sup>12</sup>12The oscillator wave functions are $`\phi _k(x):=e^{x^2/2}H_k(x)/\sqrt{2^kk!\pi ^{1/2}}`$ and form an orthonormal basis for $`L^2((0,\mathrm{}))`$. The Hermite kernel is $`K_{H,n}(x,y):=_{k=0}^{n1}\phi _k(x)\phi _k(y)`$. then this equals
$$\frac{1}{2\sqrt{a}}\phi _{k1}\left(\frac{s+x}{2\sqrt{a}}\right)\phi _{k1}\left(\frac{s+y}{2\sqrt{a}}\right)$$
times the factor
$$e^{(s+x)^2/8a}e^{(s+y)^2/8a}.$$
It follows that if $`K_{H,n}`$ is the Hermite kernel then
$$L_n(x,y)=e^{(s+x)^2/8a}\frac{1}{2\sqrt{a}}K_{H,n}(\frac{s+x}{2\sqrt{a}},\frac{s+y}{2\sqrt{a}})e^{(s+y)^2/8a}.$$
(3.24)
We deduce that
$$\underset{m\mathrm{}}{lim}\mathrm{Prob}(Hpm+sm^{1/2})$$
is equal to the Fredholm determinant of
$$\frac{1}{2\sqrt{a}}K_{H,n}(\frac{s+x}{2\sqrt{a}},\frac{s+y}{2\sqrt{a}})$$
over $`(0,\mathrm{})`$, or equivalently the Fredholm determinant of $`K_{H,n}(x,y)`$ over $`(s/2\sqrt{a},\mathrm{})`$. It is notationally convenient to introduce
$$\sigma ^2:=2a=p(1p)$$
and to define
$$F_n^{GUE}(s):=\underset{m\mathrm{}}{lim}\text{Prob}\left(\frac{Hpm}{\sigma \sqrt{m}}s\right).$$
This equals the Fredholm determinant of $`K_{H,n}`$ over $`(s/\sqrt{2},\mathrm{})`$ and is equal to the distribution of the largest eigenvalue in the finite $`n`$ GUE.<sup>13</sup><sup>13</sup>13Our normalization of $`F_n^{GUE}`$ differs from the usual one by a factor of $`\sqrt{2}`$, i.e. the usual normalization is the Fredholm determinant of the Hermite kernel over $`(s,\mathrm{})`$.
#### 3.4.2 Moments of $`F_n^{GUE}`$
From the theory of random matrices, e.g. , we know that the distribution function $`det(IK_{H,n})`$ has an alternative representation as an $`n\times n`$ determinant. Explicitly,
$$F_n^{GUE}(s)=det\left(\delta _{i,j}_{s/\sqrt{2}}^{\mathrm{}}\phi _i(x)\phi _j(x)𝑑x\right)_{0i,jn1}$$
where $`\phi _j`$ are the oscillator functions previously introduced. This last representation implies that the $`F_n^{GUE}`$ are expressible in terms of elementary functions and the error function with increasing complexity for increasing values of $`n`$. In the simplest case, $`n=1`$, $`F_1^{GUE}`$ is the standard normal; a result easily anticipated from the original formulation of the growth model. The next simplest case is $`n=2`$,
$$F_2^{GUE}(s)=\frac{1}{4}\frac{1}{2\pi }e^{s^2}\frac{1}{2^{3/2}\sqrt{\pi }}se^{s^2/2}+\frac{1}{2}(1\frac{1}{\sqrt{2\pi }}se^{s^2/2})\text{erf}(s/\sqrt{2})+\frac{1}{4}\text{erf}(s/\sqrt{2})^2.$$
The moments of $`F_n^{GUE}`$ are, of course,
$$\mu _j(n):=_{\mathrm{}}^{\mathrm{}}s^jf_n^{GUE}(s)𝑑s,j=1,2,\mathrm{}$$
where $`f_n^{GUE}=dF_n^{GUE}/ds`$. For $`1n5`$ we have,
First Moments:
$`\mu _1(1)`$ $`=`$ $`0,`$
$`\mu _1(2)`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{\pi }}}1.128379,`$
$`\mu _1(3)`$ $`=`$ $`{\displaystyle \frac{27}{8\sqrt{\pi }}}1.904140,`$
$`\mu _1(4)`$ $`=`$ $`{\displaystyle \frac{7}{48\sqrt{2}\pi ^{3/2}}}+{\displaystyle \frac{475}{128\sqrt{\pi }}}+{\displaystyle \frac{475}{64\pi ^{3/2}}}\mathrm{arcsin}(1/3)2.528113,`$
$`\mu _1(5)`$ $`=`$ $`{\displaystyle \frac{13715}{4096\sqrt{\pi }}}{\displaystyle \frac{16975}{41472\sqrt{2}\pi ^{3/2}}}+{\displaystyle \frac{41145}{2048\pi ^{3/2}}}\mathrm{arcsin}(1/3)3.063268.`$
Second Moments:
$`\mu _2(1)`$ $`=`$ $`1,`$
$`\mu _2(2)`$ $`=`$ $`2,`$
$`\mu _2(3)`$ $`=`$ $`3+{\displaystyle \frac{9\sqrt{3}}{4\pi }}4.240490,`$
$`\mu _2(4)`$ $`=`$ $`4+{\displaystyle \frac{16}{\sqrt{3}\pi }}6.940420,`$
$`\mu _2(5)`$ $`=`$ $`5{\displaystyle \frac{155\sqrt{5}}{864\pi ^2}}+{\displaystyle \frac{2495}{108\sqrt{3}\pi }}+{\displaystyle \frac{499}{54\sqrt{3}\pi ^2}}\mathrm{arcsin}(1/4)1.977575.`$
Third Moments:
$`\mu _3(1)`$ $`=`$ $`0,`$
$`\mu _3(2)`$ $`=`$ $`{\displaystyle \frac{7}{\sqrt{\pi }}}3.949327,`$
$`\mu _3(3)`$ $`=`$ $`{\displaystyle \frac{297}{16\sqrt{\pi }}}10.472769,`$
$`\mu _3(4)`$ $`=`$ $`{\displaystyle \frac{333}{32\sqrt{2}\pi ^{3/2}}}+{\displaystyle \frac{7109}{256\pi ^{1/2}}}+{\displaystyle \frac{7109}{128\pi ^{3/2}}}\mathrm{arcsin}(1/3)20.378309,`$
$`\mu _3(5)`$ $`=`$ $`{\displaystyle \frac{2595475}{82944\sqrt{2}\pi ^{3/2}}}+{\displaystyle \frac{259385}{8192\sqrt{\pi }}}+{\displaystyle \frac{778155}{4096\pi ^{3/2}}}33.432221.`$
Fourth Moments:
$`\mu _4(1)`$ $`=`$ $`3,`$
$`\mu _4(2)`$ $`=`$ $`9,`$
$`\mu _4(3)`$ $`=`$ $`19+{\displaystyle \frac{33\sqrt{3}}{2\pi }}28.096927,`$
$`\mu _4(4)`$ $`=`$ $`33+{\displaystyle \frac{496}{3\sqrt{3}\pi }}63.384348,`$
$`\mu _4(5)`$ $`=`$ $`51+{\displaystyle \frac{7475\sqrt{5}}{1296\pi ^2}}+{\displaystyle \frac{99575}{324\sqrt{3}\pi }}+{\displaystyle \frac{99575}{162\sqrt{3}\pi ^2}}\mathrm{arcsin}(1/4)117.872208.`$
Let $`H_n^{\mathrm{}}`$ denote the weak limit $`m\mathrm{}`$, $`n`$ fixed, of
$$\frac{Hpm}{\sigma \sqrt{m}},$$
and $`H^{\mathrm{}}`$ the weak limit $`m\mathrm{}`$, $`n\mathrm{}`$, $`\alpha =n/m`$ fixed, of
$$\frac{1}{v(3b)^{1/3}m^{1/3}}\left(Hcm\right).$$
(Thus the distribution functions of $`H_n^{\mathrm{}}`$ and $`H^{\mathrm{}}`$ are $`F_n^{GUE}`$ and $`F_2`$, respectively.) For $`\alpha 0`$, $`c=p+2\sigma \sqrt{\alpha }+O(\alpha )`$, $`(3b)^{1/3}\sigma /\alpha ^{1/6}`$, and $`v1`$. Proceeding heuristically,
$`H`$ $``$ $`cm+m^{1/3}v(3b)^{1/3}H^{\mathrm{}}`$
$``$ $`pm+2\sigma \sqrt{\alpha }m+m^{1/3}\sigma \alpha ^{1/6}H^{\mathrm{}}`$
$``$ $`pm+\sigma m^{1/2}\left\{2\sqrt{n}+{\displaystyle \frac{H^{\mathrm{}}}{n^{1/6}}}\right\}.`$
Thus we expect
$$H_n^{\mathrm{}}2\sqrt{n}+\frac{H^{\mathrm{}}}{n^{1/6}},$$
and hence
$`E(H_n^{\mathrm{}})`$ $``$ $`2\sqrt{n}+{\displaystyle \frac{E(H^{\mathrm{}})}{n^{1/6}}},E(H^{\mathrm{}})=1.77109\mathrm{},`$ (3.25)
$`\text{Var}(H_n^{\mathrm{}})`$ $``$ $`{\displaystyle \frac{\text{Var}(H^{\mathrm{}})}{n^{1/3}}},\text{Var}(H^{\mathrm{}})=0.8132\mathrm{}.`$ (3.26)
For $`2n9`$, these approximations are compared with the exact moments in Table 2. We also compute the skewness and the excess kurtosis<sup>14</sup><sup>14</sup>14The skewness of a random variable $`X`$ is $`E\left((\frac{X\mu }{\sigma })^3\right)`$ and the excess kurtosis is $`E\left((\frac{X\mu }{\sigma })^4\right)3`$. Here $`\mu =E(X)`$ and $`\sigma ^2=\text{Var}(X)`$. of $`F_n^{GUE}`$.
The densities $`f_n^{GUE}`$, $`1n7`$, are graphed in Fig. 7.
## 4 Brownian Motion Representation in the Finite $`x`$ GUE Regime
Let $`B(t)=(B_0(t),\mathrm{},B_x(t))`$ be the $`(x+1)`$-dimensional Brownian motion. Let $`F`$ be the following functional on continuous functions $`f=(f_0,\mathrm{},f_x)`$ from $`[0,1]`$ to $`\text{R}^{x+1}`$, which satisfy $`f(0)=0`$,
$$F(f):=\mathrm{max}\left\{f_0(t_0)+f_1(t_1)f_1(t_0)+\mathrm{}+f_x(t_x)f_x(t_{x1}):0t_0t_1\mathrm{}t_x=1\right\}.$$
Note that $`F`$ is continuous in the $`L^{\mathrm{}}`$ metric. Finally, let
$$M_x:=F(B).$$
Theorem. For $`x\text{Z}_+`$, and $`t\mathrm{}`$, we have
$$\frac{h_t(x)pt}{\sigma \sqrt{t}}\stackrel{d}{}M_x,$$
where $`\sigma ^2=p(1p)`$.
Proof. First, the path representation tells us that we change $`h_t(x)`$ by at most a constant if we only obey the increasing property within the same column. That is, we change $`L(x,t)`$ to $`L^{}(x,t)`$, where $`L^{}(x,t)`$ is the longest path $`(x_i,t_i),i=1,\mathrm{},k`$, of marked points such that $`0t_it_{i1}1`$ if $`x_i=x_{i1}`$, while $`0x_ix_{i1}t_it_{i1}1`$ if $`x_ix_{i1}`$. Thus, the first observation is
$$|L(x,t)L^{}(x,t)|x.$$
Let $`S_k^i`$ equal the length of the longest increasing sequence of points $`(t,i)`$, $`0tk`$. Then
$$L^{}(x,t)=\mathrm{max}\{S_{k_1}^0+S_{k_2}^1S_{k_1}^1+\mathrm{}+S_{k_x}^xS_{k_{x1}}^x:0k_1k_2\mathrm{}k_x=tx\}.$$
(4.27)
Now for every fixed $`i`$, $`S_k^i`$ is independent of $`S_k^j`$ for $`ji`$. Let $`X_k^i`$ equal the indicator of the event that $`(i,k)`$ is a marked point. Of course, $`S_k^i=_{\mathrm{}=1}^kX_{\mathrm{}}^i`$.
Let $`S^i(\tau )`$, $`\tau \text{R}_+`$, equal $`S_k^i`$ when $`\tau =k`$ and be obtained by linear interpolation off the integers. Moreover, let $`\stackrel{~}{S}^i`$ be the centered versions $`\stackrel{~}{S}^i(\tau )=S^i(\tau )p\tau `$, and $`\stackrel{~}{S}(\tau )=(\stackrel{~}{S}_0(\tau ),\mathrm{},\stackrel{~}{S}_x(\tau ))`$. For $`0\tau 1`$, define
$$X_t(\tau ):=\frac{\stackrel{~}{S}(t\tau )}{\sigma \sqrt{t}},$$
then the standard invariance principle (see, e.g. , Ch. 7.) implies that $`X_t`$ converges as $`t\mathrm{}`$ in distribution to the $`(x+1)`$-dimensional Brownian motion $`B`$.
Now define
$$L^{\prime \prime }(x,t)=\mathrm{max}\{\stackrel{~}{S}^0(t_0)+\stackrel{~}{S}^1(t_1)\stackrel{~}{S}^1(t_0)+\mathrm{}+\stackrel{~}{S}^x(t_x)\stackrel{~}{S}^x(t_{x1}):0t_0t_1\mathrm{}t_x=t\},$$
(4.28)
then
$$\left|L^{}(x,t)ptL^{\prime \prime }(x,t)\right|5x.$$
(The linear interpolation gives an error of at most four at each $`t_i`$ and we incur an additional $`x`$ by replacing $`tx`$ by $`t`$.) Note now that (by making a substitution $`t_i^{}=t_i/t`$)
$$\frac{L^{\prime \prime }(x,t)}{\sigma \sqrt{t}}=F\left(X_t\right).$$
It follows (by continuity of $`F`$) that the theorem holds with $`L^{\prime \prime }(x,t)`$ in place of $`h_t(x)`$, but this is clearly enough.
Remark 1. We should note that this theorem clearly holds in more general circumstances. For example, we could make every $`\times `$ count an independent random number of jumps. We would get the same theorem, with the only assumption that the said random number has finite variance.
Remark 2. The theorem also holds for random words over an alphabet with $`n`$ letters , except that the $`x+1`$ Brownian motions are not independent, but they have to sum to 0, so the covariances $`\mathrm{\Gamma }_{ij}`$ equals $`(n1)/n^2`$ when $`i=j`$ and $`1/n^2`$ otherwise. In the case of two equiprobable letters, the limiting distribution of the centered and normalized length of the longest weakly increasing subsequence in a random word is equal to the distribution of the random variable
$`X`$ $`=`$ $`\underset{0t1}{\mathrm{max}}\left(B_0(t)+B_1(1)B_1(t)\right)`$
$`=`$ $`2\underset{0t1}{\mathrm{max}}\left(B_0(t)\right)B_0(1)`$
$`=`$ $`2MN.`$
($`M`$ denotes the random variable $`\mathrm{max}_{0t1}B_0(t)`$ and $`N`$ denotes the random variable $`B_0(1)`$.) From the reflection principle it follows (see, e.g. pg. 395 in ) that the joint density of $`(M,N)`$ is
$$f_{M,N}(m,n)=\sqrt{\frac{2}{\pi }}(2mn)e^{(2mn)^2/2},\text{for}m0,mn.$$
Thus the density of $`X`$ equals<sup>15</sup><sup>15</sup>15C. Grinstead, in unpublished notes, also found a random walk interpretation of the two-letter random word problem and used this to determine the limiting distribution in this case.
$$f_X(x)=_0^xf_{M,N}(m,2mx)𝑑m=\sqrt{\frac{2}{\pi }}x^2e^{x^2/2}.$$
Remark 3. Limiting distribution of the centered and normalized $`h_t(1)`$: Here we have
$$M_1=\underset{0t1}{\mathrm{max}}(B_1(t)+(B_2(1)B_2(t)))=\underset{0t1}{\mathrm{max}}(B_1(t)B_2(t))+B_2(1)=M+N.$$
(The random variables $`M`$ and $`N`$ are defined by the last equality; and therefore, are not to be confused with the random variables of the previous remark.) Note that $`N`$ is standard normal. Since $`(B_1B_2)/\sqrt{2}`$ is the standard Brownian motion, $`M`$ equals, in distribution, $`\sqrt{2}|N|`$ again by the reflection principle. Even though $`M`$ and $`N`$ are not independent, $`E(M_1)=\sqrt{2}E(|N|)=2/\sqrt{\pi }`$. Moreover, the conditional distribution of $`N`$ given the entire path of $`W:=B_1(t)B_2(t)`$, $`0t1`$, depends only on its final point $`W(1)`$. Given this final point $`S`$ equals $`s`$, the distribution is normal with mean $`s/2`$ and variance $`1/2`$. That is, if $`_t`$ is the Brownian filtration for $`W`$, $`S=W(1)`$, then
$$\text{Prob}\left(Ndn|_1\right)=\text{Prob}(Ndn|S=s)=\frac{1}{\pi }e^{(x+s/2)^2}dn.$$
This makes it immediately possible to compute the second moment of $`M_1`$, since
$$E(MN)=E\left(E(MN|_1)\right)=E\left(ME(N|_1)\right)=E(MS)/2=E\left((M/\sqrt{2})(S/\sqrt{2})\right)=1/2,$$
by a straightforward computation with the joint density above. Therefore $`E(M_1^2)=E(M^2)+E(N^2)+2E(MN)=2`$.
In this way, the density of $`M_1`$ is
$`f_{M_1}(x)`$ $`=`$ $`E\left(\text{Prob}\left(M_1=x|M,S\right)\right)`$
$`=`$ $`{\displaystyle }f_{M,S}(m,s)\text{Prob}(M_1=x|M=m,S=s)dmds`$
$`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑m{\displaystyle _{\mathrm{}}^m}{\displaystyle \frac{1}{2}}f_{M,N}(m/\sqrt{2},n/\sqrt{2}){\displaystyle \frac{1}{\sqrt{\pi }}}e^{(xm+n/2)^2}𝑑n.`$
An explicit evaluation shows this last integral equals, as it must, $`f_2^{GUE}(x)`$.
Remark 4. Since the distribution function of $`M_x`$ equals $`F_{x+1}^{GUE}`$, it follows from RMT that we have the alternative representation
$$\text{Prob}\left(M_xs\right)=c_n_{\mathrm{}}^s\mathrm{}_{\mathrm{}}^s\mathrm{\Delta }(x)^2e^{\frac{1}{2}{\scriptscriptstyle x_j^2}}𝑑x_1\mathrm{}𝑑x_n$$
(4.29)
where
$$\mathrm{\Delta }(x)=\mathrm{\Delta }(x_1,\mathrm{},x_n)=\underset{1i<jn}{}(x_ix_j)$$
is the Vandermonde determinant, $`c_n^1=1!2!\mathrm{}n!(2\pi )^{n/2}`$, and $`n=x+12`$. In the context of Brownian motion, can one directly prove (4.29)?
Remark 5. For connections between Brownian motion exit times and random matrices, see Grabiner .
Remark 6. The Brownian motion functional $`M_x`$ has appeared previously in Glynn and Whitt , and consequently (4.29) provides an exact formula for the limiting distribution of the departure time of the first $`(x+1)`$ customers from $`n`$ single server queues. Glynn and Whitt also consider the case when $`x=t^a`$, $`0<a<1`$, and prove what would, in our setting, be the following limit theorem
$$\underset{t\mathrm{}}{lim}\frac{h_t(x)pt}{\sigma \sqrt{tx}}=\alpha :=\underset{x\mathrm{}}{lim}\frac{M_x}{\sqrt{x}},$$
with both limits in probability. They conjectured that $`\alpha =2`$, and this was later proved by Seppäläinen via a hydrodynamic limit for simple exclusion. We note that our paper proves that $`\alpha =2`$ as well, by a completely different route. Namely, one only needs to apply the result (see, e.g. ) that the largest eigenvalue in the finite $`n`$ GUE scales as $`2\sqrt{n}`$.<sup>16</sup><sup>16</sup>16We thank Timo Seppäläinen for bringing this connection with queuing theory to our attention.
Acknowledgments
This work was supported, in part, by the National Science Foundation through grants DMS–9703923, DMS–9802122 and DMS–9732687. In addition, the first author was supported in part by the Republic of Slovenia’s Ministry of Science, grant number J1–8542–0101–97. It is our pleasure to acknowledge Iain Johnstone, Bruno Nachtergaele, Timo Seppäläinen and Richard Stanley for helpful comments. Finally, we wish to thank both referees for their helpful comments and suggestions.
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# What are the Confining Field Configurations of Strong-Coupling Lattice Gauge Theory?
## Abstract:
Starting from the strong-coupling SU(2) Wilson action in $`D=3`$ dimensions, we derive an effective, semi-local action on a lattice of spacing $`L`$ times the spacing of the original lattice. It is shown that beyond the adjoint color-screening distance, i.e. for $`L5`$, thin center vortices are stable saddlepoints of the corresponding effective action. Since the entropy of these stable objects exceeds their energy, center vortices percolate throughout the lattice, and confine color charge in half-integer representations of the SU(2) gauge group. This result contradicts the folklore that confinement in strong-coupling lattice gauge theory, for $`D>2`$ dimensions, is simply due to plaquette disorder, as is the case in $`D=2`$ dimensions. It also demonstrates explicitly how the emergence and stability of center vortices is related to the existence of color screening by gluon fields.
Confinement, Lattice Gauge Field Theories, Solitons Monopoles and Instantons
Quark confinement is commonly attributed to the influence of some special class of gauge field configurations, which dominate the QCD vacuum at large scales. Because of their high probability, these dominant configurations would most naturally correspond to the saddlepoints of an infrared effective action, derived at large scales by integrating over high-frequency modes. In strong-coupling lattice gauge theory there are methods available which enable us to compute the QCD spectrum and string tension analytically, and the same methods could also be applied to extract a long-range effective action. An interesting question, then, is what type of saddlepoint configurations are actually found at strong lattice couplings; it is likely that the answer would also shed some light on QCD in the continuum limit.
The classical Euclidean action of pure SU(N) gauge theory is stationary, or nearly so, at multi-instanton configurations. In quantized lattice gauge theory, however, we can imagine performing renormalization-group (RG) transformations so as to obtain an effective action at some scale $`R`$. For scales $`R`$ well below the confinement scale, the main effect of the RG transformations will simply be the running of the lattice coupling constant. At larger scales, however, so-called irrelevant operators can become important in the effective action. As a consequence, at these larger scales, the effective theory may have non-trivial saddlepoints which are something other than instantons.
There are good reasons to believe that at sufficiently large scales, these non-trivial saddlepoints are center vortices. On the theoretical side, we note that the asymptotic string tension between static color charges in SU(N) gauge theory depends only on the $`N`$-ality of the color charge representation. Although this fact is deduced rather trivially from the possibility of color-screening by gluon fields, it has some profound implications for the infrared structure of the QCD vacuum. Consider, for example, Wilson loops $`W_j(C)`$ in SU(2) gauge theory, where $`j=0,\frac{1}{2},1,\frac{3}{2},\mathrm{}`$ labels the group representation. Wilson loop expectation values can be viewed as a probe of vacuum fluctuations in the *absence* of external sources (think of evaluating a spacelike loop in the Hamiltonian formulation), and large Wilson loops are presumed to become “disordered,” i.e. have an area-law falloff, due to averaging over certain types of large-scale fluctuations which dominate the vacuum state. Whatever the nature of these confining fluctuations, they must have the highly non-trivial property of disordering only the $`j=`$ half-integer loops, but not the $`j=`$ integer loops. Center vortices are the only configurations we know of that have this property. Dual-superconductor models, in which all multiples of abelian electric charge (identified in an abelian-projection gauge) are confined by the dual Meissner effect, do not seem satisfactory. In these models, the potential between charged objects is roughly proportional to the electric charge. But this charge dependence cannot be correct, since gluons carrying two units of electric charge are available to screen multiply charged sources, and numerical simulations indicate that only odd multiples of the abelian charge (non-zero $`N`$-ality) are confined, while even multiples of abelian charge (zero $`N`$-ality) are screened . Center vortices seem to be the natural way of accounting, in terms of dominant field configurations, for this very fundamental distinction between zero and non-zero $`N`$-ality, which is evident even in the abelian projection.
On the numerical side, there is now abundant evidence in favor of the vortex theory of confinement , which we will not attempt to review here. The present situation can just be summarized as follows: There exists a method (known as “center projection”) for locating center vortices on thermalized lattices; the rationale undelying this method is explained in ref. . By locating the vortices, their effects on gauge-invariant observables such as Wilson loops, Polyakov lines, topological charge, etc., can be studied in detail. The numerical evidence indicates that fluctuations in vortex linking number are the origin of the asymptotic string tension of Wilson loops. The free energy of a vortex, inserted into a finite lattice via twisted boundary conditions, has also been computed, and has been shown to fall off exponentially with the lattice size at just the rate predicted by the vortex theory .
Vortices presumably have a finite thickness comparable to the adjoint string-breaking length, at about $`1.25`$ fm , where the crossover from Casimir scaling to $`N`$-ality confinement occurs (cf. the discussion in ref. ). Independent estimates of the vortex thickness are based on measurements of “vortex-limited” Wilson loops in ref. , and on the vortex free energy in finite volumes . Both of these estimates put the vortex thickness at a little over one fermi. Beyond this scale, the presence of vortex sheets in the vacuum should be very evident. A reasonable conjecture is that if the appropriate effective action could be determined at this scale, it would be found to have stable saddlepoints corresponding to vortex configurations, which percolate through the lattice according to the usual energy-entropy arguments. Unfortunately, the calculation of long-range effective actions is very difficult even numerically, via Monte Carlo RG methods, and at such large scales the problem is quite intractable by perturbative (e.g. one-loop) methods.<sup>1</sup><sup>1</sup>1On the other hand, there do exist some intriguing results at one-loop that should be noted. Diakonov has computed a one-loop effective potential for magnetic flux tubes, and his result indicates that the potential is minimized for magnetic flux in the center of the gauge group. There is also the old, but still provocative, Copenhagen vacuum picture, which is again based on one-loop considerations .
We therefore turn our attention, in this article, to strong-coupling lattice gauge theory, where analytic methods can be brought to bear at arbitrarily large scales. In the strongly coupled theory in $`D>2`$ dimensions, we have a linear static potential for all color charge representations up to a screening scale of about four lattice spacings (for the SU(2) Wilson action). Beyond that scale, the string tension depends only on the $`N`$-ality of the representation, just as in the continuum theory. Thus, if our conjecture is correct, and if the $`N`$-ality dependence implies a vortex mechanism, the effective strong-coupling action at a scale beyond four lattice spacings should have saddlepoints which are stable center vortices.
There is, however, some folklore to the effect that confinement in strong coupling lattice gauge theory, in $`D>2`$ dimensions, is essentially the same as in $`D=2`$ dimensions, where the mechanism is simply plaquette disorder. If that were so, then vortices (or any other topological objects) are irrelevant at strong couplings in any dimension. This folklore is quite misleading, according to an argument presented in ref. , which we now review.
Consider SU(2) lattice gauge theory at strong-coupling, and denote by $`U(C)`$ the product of link variables around loop $`C`$. Let the minimal area of a planar loop be decomposed into a set of smaller areas, bounded by loops $`\{C_i\}`$, as shown in Fig. 1. The area-law of a Wilson loop is thought to be due to “magnetic disorder”, in which the gauge field strength fluctuates independently in sub-areas of the minimal surface of the loop. If this is so, then the holonomies $`\{U(C_i)\}`$ should be (nearly) uncorrelated, for large areas and small $`\beta `$. The test for such independent fluctuation in the subareas $`A(C_i)`$ is whether
$$<\underset{i}{}F[U(C_i)]>\stackrel{\mathrm{?}}{=}\underset{i}{}<F[U(C_i)]>$$
(1)
for any class function
$$F[g]=\underset{j0}{}f_j\chi _j[g]$$
(2)
In fact, in $`D=2`$ dimensions, it is easy to show that this equality is satisfied exactly. However, for dimensions $`D>2`$, evaluating the left- and right-hand sides of (1) one finds for the exponential falloff on each side
$$e^{4\sigma 𝒫(C)}\underset{i}{}\frac{1}{3}f_1\underset{i}{}f_1e^{4\sigma 𝒫(C_i)}$$
(3)
where the inequality holds for perimeters $`𝒫(C)_i𝒫(C_i)`$. The conclusion is that the holonomies $`U(C_i)`$ do *not* fluctuate independently, even at strong-coupling, for $`D>2`$.
If the loop holonomies in sub-areas of the loop are correlated, then where is the magnetic disorder required to give an area-law falloff for Wilson loops? The question is resolved by extracting a center element from the holonomies
$$z[U(C_i)]=\text{signTr}[U(C_i)]Z_2$$
(4)
and asking if these center elements fluctuate independently; i.e
$$<\underset{i}{}z[U(C_i)]>\stackrel{\mathrm{?}}{=}\underset{i}{}<z[U(C_i)]>$$
(5)
In fact, it is easy to show that they do:
$$e^{\sigma A(C)}\underset{i}{}\frac{3}{4\pi }=\underset{i}{}\frac{3}{4\pi }e^{\sigma A(C_i)}$$
(6)
Thus, *magnetic disorder is center disorder* in $`D>2`$ dimensions, at least at strong couplings. Confining configurations must disorder the center elements $`z`$, but not the coset elements, of SU(2) holonomies $`U(C_i)`$. Again, the only configurations known to have this property are center vortices.
We now return to our conjecture that vortices are stable saddlepoints of a long-range effective action. Actually, there are various ways of integrating out the smaller-scale fluctuations, to obtain an effective action at a larger scale. One simple approach is to superimpose, on a lattice of spacing $`a`$ with link variables denoted $`U`$, a lattice of spacing $`La`$ with links denoted $`V`$. An effective action for the lattice with the larger spacing can then be obtained from
$$\mathrm{exp}\left[S_{eff}[V]\right]=DU\underset{l^{}}{}\delta [V_l^{}^{}(UU..U)_l^{}I]e^{S_W[U]}$$
(7)
where $`(UU..U)_l^{}`$ is the product of $`U`$-link variables along the link $`l^{}`$ of the V-lattice, and $`S_W`$ is the Wilson action. Obviously, all observables computed on the V-lattice with $`S_{eff}`$ will agree with the corresponding quantity computed on the U-lattice using $`S_W`$.
It is trivial to compute $`S_{eff}`$ in $`D=2`$ dimensions, and the result is
$`\mathrm{exp}\left[S_{eff}[V]\right]`$ (8)
$``$ $`{\displaystyle \underset{P^{}}{}}{\displaystyle \underset{j}{}}(2j+1)\left(I_{2j+1}(\beta )\right)^{L^2}\chi _j[V(P^{})]`$
$`=`$ $`\mathrm{exp}\left[{\displaystyle \underset{P^{}}{}}\mathrm{log}\left(1+{\displaystyle \underset{j=\frac{1}{2},1,\frac{3}{2}}{}}(2j+1)\left({\displaystyle \frac{I_{2j+1}(\beta )}{I_1(\beta )}}\right)^{L^2}\chi _j[V(P^{})]\right)+\text{const.}\right]`$
$``$ $`\mathrm{exp}\left[2\left({\displaystyle \frac{\beta }{4}}\right)^{L^2}{\displaystyle \underset{P^{}}{}}\chi _{1/2}[V(P^{})]+\text{const.}\right]`$
where $`V(P^{})`$ is the product of V-links around the plaquette $`P^{}`$. One might imagine that the action (8) is also a good approximation in $`D>2`$ dimensions, at least at strong couplings, since this action gives the correct string tension for fundamental representation Wilson loops. But a quick calculation of higher representation loops shows that (8) cannot even be approximately correct for large $`L`$. A loop in the adjoint representation, for example, calculated on the U-lattice with the Wilson action, is easily seen to have an asymptotic perimeter law falloff
$$W_{adj}(C)\mathrm{exp}[\mu 𝒫(C)]$$
(9)
where
$$\mu =4\mathrm{log}\left(\frac{\beta }{4}\right)$$
(10)
is the “gluelump” mass (gluon bound to a static adjoint color charge), and $`𝒫(C)`$ is the loop perimeter in units of the U-lattice spacing. However, carrying out the same calculation with the effective action (8), one finds instead (with $`𝒫(C)`$ again in U-lattice units) the erroneous result
$$\mu =4L\mathrm{log}\left(\frac{\beta }{4}\right)\text{(wrong)}$$
(11)
with an $`L`$-dependent gluelump mass. The mismatch is not resolved by including a few more contours in the effective action, so long as the coupling associated with each contour is of order $`(\beta /4)^A`$, where $`A`$ is the minimal area of the contour in U-lattice units.
In fact, what happens in $`D>2`$ dimensions is that the effective action contains Wilson loops of all sizes in $`j=`$ integer representations, and these loops are only suppressed by perimeter-law coefficients. This is easily seen by bringing down a “tube” of plaquettes from $`\mathrm{exp}(S_W)`$ in eq. (7), such that the tube borders a rectangular contour $`C`$ on the V-lattice, as shown in Fig. 2. Integrating over all U-links in the tube, except those which lie on contour $`C`$, will yield contributions to $`S_{eff}`$ such as
$`\mathrm{exp}\left[S_{eff}[V]\right]`$ $``$ $`{\displaystyle }DU_{lC}{\displaystyle \underset{l^{}C}{}}\delta [V_l^{}^{}(UU..U)_l^{}I]\left({\displaystyle \frac{\beta }{4}}\right)^{4(𝒫(C)4)}\left(\chi _{\frac{1}{2}}[U(C)]\right)^2`$ (12)
$``$ $`\left({\displaystyle \frac{\beta }{4}}\right)^{4(𝒫(C)4)}(\chi _1[V(C)]+\text{const.})`$
This shows that $`S_{eff}[V]`$ contains adjoint (and, in general, integer) representation loops with perimeter-falloff coefficients. Such non-local terms introduce non-local correlations among $`SU(2)/Z_2`$ coset elements in loop holonomies $`\{U(C_i)\}`$. Truncating $`S_{eff}[V]`$ by removing these large loops will yield erroneous results for any Wilson loop in representation $`j>\frac{1}{2}`$.
Our aim is to modify the prescription (7) for the effective action, in such a way that at least the *leading* contribution to any Wilson loop on the V-lattice is obtained from a local action. The strategy for doing this is to prevent the formation of closed tube diagrams, of the form shown in Fig. 2, bordering contours on the V-lattice. This can be accomplished by not integrating, in eq. (7), over the U-links in a cube of volume $`2^D`$ around each site on the V-lattice.<sup>2</sup><sup>2</sup>2Non-local terms in the effective action will still arise from closed tubes which go around the 2-cubes. These, however, are associated with sub-leading contributions to color screening; the leading contributions in $`\beta `$, arising from diagrams like Fig. 2, will now be generated by local terms. U-links belonging to these 2-cubes will be denoted $`\stackrel{~}{U}_l`$. In order to ease the task of illustration, we will work in $`D=3`$ dimensions, although the extension to higher dimensions should be straightforward. We then have
$`Z`$ $`=`$ $`{\displaystyle DV\underset{l2cubes}{}d\stackrel{~}{U}_l\mathrm{exp}\left[\stackrel{~}{S}_L[V,\stackrel{~}{U}]\right]}`$ (13)
$`=`$ $`{\displaystyle }DV{\displaystyle }{\displaystyle \underset{l2cubes}{}}d\stackrel{~}{U}_l\left\{{\displaystyle }{\displaystyle \underset{l^{\prime \prime }2cubes}{}}dU_{l^{\prime \prime }}{\displaystyle \underset{l^{}}{}}\delta [V_l^{}^{}(UU..U)_l^{}I]e^{S_W[U]}\right\}`$
where the long-range action $`\stackrel{~}{S}_L[V,\stackrel{~}{U}]`$ depends on the $`V`$-link variables, and on the $`\stackrel{~}{U}`$-links in 2-cubes around sites of the V-lattice, as shown in Fig. 3.
Now introduce, in each 2-cube, a set of plaquette variables $`\{h_{ij},g_{ij}\}`$ which are Wilson lines beginning and ending at the center of the 2-cube, and running around one of the plaquettes in the cube. The $`h_{ij}`$ lines run around plaquettes on the faces of the 2-cube, and the $`g_{ij}`$ lines run around plaquettes on the interior of the 2-cubes. To fix notation: The orientation of the $`g_{ij}`$ lines is taken to be counterclockwise in either the $`xy`$-plane ($`i=1`$), the $`yz`$-plane ($`i=2`$), or the $`zx`$-plane ($`i=3`$). The second index ($`j=14`$) distinguishes between the four interior plaquettes in a given plane with the convention shown in Fig. 4, where $`(x_a,x_b)=(xy),(yz),(zx)`$. Each $`h_{ij}`$ line begins at the center of the 2-cube, runs to a center of one of the faces of the 2-cube, goes around one of the plaquettes on the face, and returns to the center of the 2-cube. The orientation around the $`h`$-plaquettes is defined by a right-hand rule: the thumb points in an outward direction normal to the 2-cube. The first index $`i=1,2,3`$ refers to a face of the 2-cube in the $`xy`$, $`yz`$, $`zx`$-planes, one lattice spacing away from the center of the 2-cube, in the negative $`z,x,y`$ directions, respectively. Index values $`i=4,5,6`$ refer to faces in the $`xy`$, $`yz`$, $`zx`$-planes one lattice spacing away from the center in the positive $`z,x,y`$ directions, respectively. These conventions are illustrated in Fig. 5.
We then integrate over the $`U`$ links which do not belong to the 2-cubes. Keeping, for each type of contribution, only terms of leading order in $`\beta `$, the result is approximately
$`Z`$ $``$ $`{\displaystyle }DVD\stackrel{~}{U}\mathrm{exp}[{\displaystyle \frac{\beta }{2}}{\displaystyle }(\text{Tr}[h]+\text{Tr}[g])`$ (14)
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^{4(L2)}{\displaystyle \underset{l^{}}{}}f_l^{}^{ijkl}\text{Tr}[h_{ij}^{}V_l^{}h_{kl}^{}V_l^{}^{}]`$
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^{L^24}{\displaystyle \underset{P^{}}{}}\text{Tr}[VgVgV^{}g^{}V^{}g^{}]]`$
Coefficients $`f_l^{}^{ijkl}=1`$ if plaquette variables $`h_{ij}`$ and $`h_{kl}`$ on nearest-neighbor 2-cubes can be joined by a cylinder of plaquettes, bordering link $`l^{}`$, of length $`L2`$ U-lattice spacings. Otherwise, $`f_l^{}^{ijkl}=0`$.
The next step is to change integration variables from links $`\stackrel{~}{U}`$ to plaquettes $`g,h`$. This change of variables on the lattice was worked out many years ago by Batrouni , and the result is simply to introduce a Bianchi constraint into the integration measure<sup>3</sup><sup>3</sup>3There are 36 $`g,h`$ plaquette variables on the 2-cube, and 8 Bianchi constraints, leaving 28 independent group-valued variables. Similarly, up to 26 out of 54 link variables on the 2-cube can be gauge fixed to the identity, again leaving 28 independent group-valued variables.
$`Z`$ $``$ $`{\displaystyle DVDhDg\underset{2cubesK}{}\underset{cK}{}\delta [\text{Bianchi}(c(K))]}`$ (15)
$`\mathrm{exp}[{\displaystyle \frac{\beta }{2}}{\displaystyle }(\text{Tr}[h]+\text{Tr}[g])`$
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^{4(L2)}{\displaystyle \underset{l^{}}{}}f_l^{}^{ijkl}\text{Tr}[h_{ij}^{}V_l^{}h_{kl}^{}V_l^{}^{}]`$
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^{L^24}{\displaystyle \underset{P^{}}{}}\text{Tr}[VgVgV^{}g^{}V^{}g^{}]]`$
Each 2-cube $`K`$ contains eight unit sub-cubes; the index $`c`$ in eq. (15) labels these subcubes. The $`\delta `$-function constraints force a certain product of three $`g`$ and three $`h`$ variables on each subcube to be the unit matrix. For example, the Bianchi constraint for the unit sub-cube containing $`h_{11}`$ is
$$\text{Bianchi}=h_{11}g_{23}h_{62}g_{11}^{}h_{53}g_{32}I=0$$
(16)
Now expand the $`\delta `$-functions in group characters
$$\delta [\text{Bianchi}]=\underset{j=0,\frac{1}{2},1,..}{}(2j+1)\chi _j[hghghg]$$
(17)
and integrate out the $`g`$-variables in eq. (15). Displaying only terms of low order in both $`h`$ and $`\beta `$, the result is
$`Z`$ $``$ $`{\displaystyle }DVDh{\displaystyle \underset{2cubesK}{}}\{1+2\left({\displaystyle \frac{\beta }{4}}\right)^3{\displaystyle \underset{cK}{}}\chi _{\frac{1}{2}}[(hhh)_c]`$ (18)
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^4{\displaystyle \underset{\stackrel{adjacent}{c_1c_2K}}{}}\chi _{\frac{1}{2}}[(hhh)_{c_1}(hhh)_{c_2}]+\mathrm{}\}`$
$`\times `$ $`\mathrm{exp}[{\displaystyle \frac{\beta }{2}}{\displaystyle }\text{Tr}[h]+2\left({\displaystyle \frac{\beta }{4}}\right)^{4(L2)}{\displaystyle \underset{l^{}}{}}f_l^{}^{ijkl}\text{Tr}[h_{ij}^{}V_l^{}h_{kl}^{}V_l^{}^{}]`$
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^{L^2}{\displaystyle \underset{P^{}}{}}\text{Tr}[VVV^{}V^{}]]`$
$``$ $`{\displaystyle DVDh\mathrm{exp}\left[S_L[V,h]\right]}`$
At this stage, the action $`S_L[V,h]`$ resembles an adjoint-Higgs Lagrangian, albeit of an unconventional form: There is an SU(2) gauge field $`V_\mu `$ coupled to 24 unitary matrix-valued “matter” fields $`h_{ij}`$ transforming in the adjoint representation. These matter fields, in turn, can be subdivided into gauge-singlet fields $`h_{ij,0}`$, and unit-modulus triplet fields $`\stackrel{}{e}_{ij}`$, where
$$h_{ij}=h_{ij,0}I+i\sqrt{1h_{ij,0}^2}\stackrel{}{e}_{ij}\stackrel{}{\sigma }$$
(19)
and $`\stackrel{}{e}\stackrel{}{e}=1`$. The unimodular $`\stackrel{}{e}_{ij}`$ degrees of freedom play a role analogous to Higgs fields. We know from the Elitzur theorem that the expectation values of these fields must vanish in the absence of gauge fixing, and cannot be viewed as order parameters. Since the coupling of the $`\stackrel{}{e}_{ij}`$ matter fields to the $`V_\mu `$ gauge field is very weak at large $`L`$, as compared to the self-couplings of the $`\stackrel{}{e}_{ij}`$ fields to each other on the 2-cubes, their expectation values depend primarily on these self-couplings, and on the complete removal of gauge redundancy in the $`e`$-fields through the choice of a unitary gauge.
We fix to a maximal unitary gauge by first transforming one of the 24 unimodular “Higgs” fields $`\stackrel{}{e}_{ij}`$ on each 2-cube, denoted $`\stackrel{}{e}_A`$, to point in the (color) 3-direction, i.e.
$$e_{A1}=e_{A2}=0,e_{A3}=1$$
(20)
This leaves a remnant U(1) symmetry, but (20) is not yet a maximal unitary gauge fixing. We then pick one other Higgs variable on each 2-cube, denoted $`\stackrel{}{e}_B`$, and use the remaining gauge freedom to fix
$$e_{B2}=0,e_{B1}=\mathrm{sin}(\theta _B)0$$
(21)
leaving a remnant $`Z_2`$ symmetry.
In the unitary gauge \[20-21\], the functional integral becomes
$$Z_{ug}=\underset{n}{}dV(n)\underset{ij}{}dh_{ij}(n)\mathrm{\Delta }(h_A,h_B)\delta (e_{A1})\delta (e_{A2})\delta (e_{B2})\mathrm{exp}\left[S_L[V,h]\right]$$
(22)
where
$`\mathrm{\Delta }^1(h_A,h_B)`$ $`=`$ $`{\displaystyle d^3\alpha \delta (e_{A1}^\alpha )\delta (e_{A2}^\alpha )\delta (e_{B2}^\alpha )}`$ (23)
$`=`$ $`{\displaystyle d^3\alpha \delta (ϵ_{1jk}\alpha _j\delta _{3k})\delta (ϵ_{2jk}\alpha _j\delta _{3k})\delta (ϵ_{2jk}\alpha _j(e_{B1}\delta _{1k}+e_{B3}\delta _{3k}))}`$
$`=`$ $`e_{B1}^1`$
From the measure
$$𝑑h=𝑑h_0d^3e\sqrt{1h_0^2}\delta (e^21)$$
(24)
we find
$$𝑑h_A𝑑h_B\mathrm{\Delta }(h_A,h_B)\delta (e_{A1})\delta (e_{A2})\delta (e_{B2})=𝑑h_{A0}𝑑h_{B0}𝑑\theta _B\sqrt{1h_{A0}^2}\sqrt{1h_{B0}^2}\mathrm{sin}\theta _B$$
(25)
Let us take, e.g., $`h_A=h_{11},h_B=h_{44}`$. Then
$`Z_{ug}`$ $`=`$ $`{\displaystyle \underset{n}{}dV(n)𝑑h_{11,0}𝑑h_{44,0}𝑑\theta _{44}\sqrt{1h_{11,0}^2}\sqrt{1h_{44,0}^2}\mathrm{sin}\theta _{44}}`$ (26)
$`{\displaystyle \underset{ij(11),(44)}{}}dh_{ij}(n)\mathrm{exp}\left[S_L[V,h]\right]`$
The final step is to integrate over the remaining $`h`$ degrees of freedom in this maximal unitary gauge. Defining $`h`$-expectation values
$`F[h]_h=`$ (27)
$`{\displaystyle \frac{1}{𝒵_h}}{\displaystyle 𝑑h_{11,0}𝑑h_{44,0}𝑑\theta _{44}\sqrt{1h_{11,0}^2}\sqrt{1h_{44,0}^2}\mathrm{sin}\theta _{44}\underset{ij(11),(44)}{}dh_{ij}(n)}`$
$`{\displaystyle \underset{2cubesK}{}}\left\{1+2\left({\displaystyle \frac{\beta }{4}}\right)^3{\displaystyle \underset{cK}{}}\chi _{\frac{1}{2}}[(hhh)_c]+2\left({\displaystyle \frac{\beta }{4}}\right)^4{\displaystyle \underset{\stackrel{adjacent}{c_1c_2K}}{}}\chi _{\frac{1}{2}}[(hhh)_{c_1}(hhh)_{c_2}]+\mathrm{}\right\}`$
$`\times \mathrm{exp}\left[{\displaystyle \frac{\beta }{2}}{\displaystyle \text{Tr}[h]}\right]F[h]`$
and
$`S_{eff}[V,h]`$ $`=`$ $`2\left({\displaystyle \frac{\beta }{4}}\right)^{4(L2)}{\displaystyle \underset{l^{}}{}}f_l^{}^{ijkl}\text{Tr}[h_{ij}^{}V_l^{}h_{kl}^{}V_l^{}^{}]`$ (28)
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^{L^2}{\displaystyle \underset{P^{}}{}}\text{Tr}[VVV^{}V^{}]`$
we have
$`Z_{ug}`$ $`=`$ $`𝒵_h{\displaystyle DV\mathrm{exp}\left[S_{eff}[V,h]\right]_h}`$ (29)
$`=`$ $`𝒵_h{\displaystyle DV\mathrm{exp}\left[S_{eff}[V]\right]}`$
where
$$S_{eff}[V]=S_{eff}[V,h_h]+\text{higher-order contributions}$$
(30)
The higher-order contributions consist of next-nearest neighbor (and more distant) couplings between $`h_h`$ terms on different 2-cubes, as well as closed loops in $`j=1`$ and higher representations. The large closed loops again introduce certain non-local correlations among the $`SU(2)/Z_2`$ elements of loop holonomies $`\{U(C_i)\}`$. But $`S_{eff}[V]`$ also contains local one-link terms, which provide by far the largest contribution to the $`Z_2`$-invariant part of the action.
For the purpose of determining saddlepoint configurations of $`S_{eff}[V]`$ we may neglect the higher-order, non-local terms in the action, so that
$`S_{eff}[V]`$ $``$ $`S_{link}[V,h_h]+S_{plaq}[V]`$ (31)
$`=`$ $`2\left({\displaystyle \frac{\beta }{4}}\right)^{4(L2)}{\displaystyle \underset{l^{}}{}}f_l^{}^{ijkl}\text{Tr}\left[h_{ij}^{}_hV_l^{}h_{kl}^{}_hV_l^{}^{}\right]`$
$`+2\left({\displaystyle \frac{\beta }{4}}\right)^{L^2}{\displaystyle \underset{P^{}}{}}\text{Tr}[VVV^{}V^{}]`$
and, for this particular gauge choice, we find
$$\begin{array}{ccc}h_{11}_h=\frac{\beta }{4}I+i\frac{8}{3\pi }\sigma _3\hfill & ,& h_{12}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{8}{3\pi }\sigma _3\hfill \\ h_{13}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{8}{3\pi }\sigma _3\hfill & ,& h_{14}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{2}{3}\sigma _1\hfill \\ h_{41}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{8}{3\pi }\sigma _3\hfill & ,& h_{42}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{2}{3}\sigma _1\hfill \\ h_{43}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{2}{3}\sigma _1\hfill & ,& h_{44}_h=\frac{\beta }{4}I+i\frac{2}{3}\sigma _1\hfill \\ h_{21}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{2}{3}\sigma _1\hfill & ,& h_{22}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^4\frac{2}{3}\sigma _1\hfill \\ h_{23}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{8}{3\pi }\sigma _3\hfill & ,& h_{24}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{2}{3}\sigma _1\hfill \\ h_{51}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{8}{3\pi }\sigma _3\hfill & ,& h_{52}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{2}{3}\sigma _1\hfill \\ h_{53}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^4\frac{8}{3\pi }\sigma _3\hfill & ,& h_{54}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{8}{3\pi }\sigma _3\hfill \\ h_{31}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{2}{3}\sigma _1\hfill & ,& h_{32}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^4\frac{2}{3}\sigma _1\hfill \\ h_{33}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{2}{3}\sigma _1\hfill & ,& h_{34}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{8}{3\pi }\sigma _3\hfill \\ h_{61}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{8}{3\pi }\sigma _3\hfill & ,& h_{62}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{2}{3}\sigma _1\hfill \\ h_{63}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^4\frac{8}{3\pi }\sigma _3\hfill & ,& h_{64}_h=\frac{\beta }{4}Ii\left(\frac{\beta }{4}\right)^8\frac{8}{3\pi }\sigma _3\hfill \end{array}$$
(32)
Inserting (32) into $`S_{link}`$ we have, to leading order in $`\beta `$,
$`S_{link}[V,h_h]=2\left({\displaystyle \frac{\beta }{4}}\right)^{4(L2)}\times {\displaystyle \underset{n}{}}`$ (33)
$`\{\left({\displaystyle \frac{8}{3\pi }}\right)^2\left({\displaystyle \frac{\beta }{4}}\right)^8\text{Tr}[\sigma _3V_z(n)\sigma _3V_z^{}(n)]+\left({\displaystyle \frac{2}{3}}\right)^2\left({\displaystyle \frac{\beta }{4}}\right)^8\text{Tr}[\sigma _1V_z(n)\sigma _1V_z^{}(n)]`$
$`\left({\displaystyle \frac{8}{3\pi }}\right)^2\left({\displaystyle \frac{\beta }{4}}\right)^{12}\text{Tr}[\sigma _3V_y(n)\sigma _3V_y^{}(n)]\left({\displaystyle \frac{2}{3}}\right)^2\left({\displaystyle \frac{\beta }{4}}\right)^{12}\text{Tr}[\sigma _1V_y(n)\sigma _1V_y^{}(n)]`$
$`\left({\displaystyle \frac{8}{3\pi }}\right)^2\left({\displaystyle \frac{\beta }{4}}\right)^{12}\text{Tr}[\sigma _3V_x(n)\sigma _3V_x^{}(n)]\left({\displaystyle \frac{2}{3}}\right)^2\left({\displaystyle \frac{\beta }{4}}\right)^{12}\text{Tr}[\sigma _1V_x(n)\sigma _1V_x^{}(n)]`$
$`+\text{const.}\}`$
Each term in $`S_{link}`$ is proportional to a component of a $`V`$-link variable in the adjoint representation, and is insensitive to the center degrees of freedom. The spatial asymmetry of $`S_{link}`$ in (33) is, of course, due to the particular unitary gauge choice.<sup>4</sup><sup>4</sup>4Despite this asymmetry, the expectation value of any Wilson loop on the V-lattice, evaluated in the full (gauge-dependent) effective action $`S_{eff}[V]`$ defined in eq. (29), is necessarily independent of the gauge choice. This should be clear from the construction, where a gauge-invariant action $`S_L[V,h]`$ is gauge-fixed, followed by integration over the remaining $`h`$ degrees of freedom.
We now look for saddlepoints of $`S_{eff}`$. It can be seen from inspection of (33) that $`S_{link}`$ is maximized by
$$V_x(n)=i\sigma _2Z_x(n),V_y(n)=i\sigma _2Z_y(n),V_z(n)=Z_z(n)$$
(34)
where the $`Z_\mu (n)=\pm I`$ are center elements. The $`S_{plaq}`$ term is also maximized if the $`Z_\mu `$ link variables are gauge-equivalent to the identity under the remnant $`Z_2`$ gauge symmetry. With this choice the effective action $`S_{eff}S_{link}+S_{plaq}`$ is maximized, and the configuration (34) is the ground state, gauge equivalent to the identity. The fact that this ground state is unique, up to a $`Z_2`$ gauge transformation, is again due to the maximal unitary gauge fixing.
Now consider the configuration
$`V_y(\stackrel{}{n})`$ $`=`$ $`\{\begin{array}{cc}\hfill i\sigma _2& n_12,n_2=1\hfill \\ \hfill +i\sigma _2& \text{otherwise}\hfill \end{array}`$ (37)
$`V_x(\stackrel{}{n})`$ $`=`$ $`i\sigma _2`$
$`V_z(\stackrel{}{n})`$ $`=`$ $`I`$ (38)
This configuration is a center vortex, one lattice spacing thick, running in the $`z`$-direction. It is not hard to see that this configuration, like the trivial ground state, is also a saddlepoint of $`S_{eff}`$. In the first place, (38) is a global maximum of $`S_{link}`$, since the thin vortex configuration (38) differs from link variables in the ground state only by center elements, to which $`S_{link}`$ is insensitive. Secondly, this configuration is also a stationary point of the plaquette action $`S_{plaq}`$ . This is because a plaquette at one of its extremal values $`\frac{1}{2}\text{Tr}[VVV^{}V^{}]=\pm 1`$ varies at most quadratically, and is therefore stationary, with respect to a small variation $`\delta V`$ of any link respecting the unitarity constraint $`(V+\delta V)(V+\delta V)^{}=I`$. All of the plaquettes formed from (38) are extremal. So the center vortex (38) is certainly a stationary configuration of $`S_{eff}`$, the remaining question is whether it is stable; i.e. whether the vortex is a local *maximum* of $`S_{eff}`$, in which case it is a stable saddlepoint.
The stability issue is settled by looking at the eigenvalues of
$$\frac{\delta ^2S_{eff}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}=\frac{\delta ^2S_{link}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}+\frac{\delta ^2S_{plaq}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}$$
(39)
where, from the coefficients shown in eqs. (31) and (33), we see that
$$\frac{\delta ^2S_{link}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}\left(\frac{\beta }{4}\right)^{4(L2)+12},\frac{\delta ^2S_{plaq}}{\delta V_\mu (n_1)\delta V_\nu (n_2)}\left(\frac{\beta }{4}\right)^{L^2}$$
(40)
The crucial observation is that for $`\beta /41`$ and
$$4(L2)+12<L^2$$
(41)
the contribution of $`S_{plaq}`$ to the stability matrix (and therefore to the eigenvalues of the stability matrix) is negligible compared to the contribution of $`S_{link}`$, which has only stable modes. Thus from the fact that the thin center vortex (38) is both a stationary point of $`S_{eff}`$, and a stable saddlepoint of $`S_{link}`$, we can conclude that the vortex is also a stable saddlepoint of the full effective action $`S_{eff}`$ when condition (41) is satisfied.
Condition (41) is satisfied for $`L5`$ lattice spacings. It is probably no coincidence that this is also where the adjoint string breaks in strong-coupling lattice gauge theory (as can be easily verified from looking at correlations of Polyakov lines in the adjoint representation). It has been known for a long time that, at intermediate distance scales and weak couplings, the static quark-antiquark potential is roughly (and maybe even accurately ) proportional to the quadratic Casimir of quark color representation. In ref. this phenomenon was dubbed “Casimir scaling,” and the problems it poses for monopole and vortex theories was discussed. In ref. we have argued that the problems with respect to the vortex theory can in principle be resolved by taking into account the finite thickness of the vortex, which should be comparable to the adjoint string-breaking length. A result of the analysis carried out above is that there are stable center vortices, one lattice spacing thick on the V-lattice, corresponding to $`L5`$ on the original U-lattice. This gives an estimate for the vortex thickness of $`L=4`$ lattice spacings in U-lattice units. This vortex thickness happens to be exactly the length where adjoint string-breaking occurs in the strong-coupling Wilson action, and is therefore consistent with the reasoning in ref. .
By inspection of $`S_{eff}`$, we see that the action of a center vortex configuration in D=3 dimensions (ignoring any further quantum corrections) is
$$8\left(\frac{\beta }{4}\right)^{L^2}\times \text{vortex length}$$
(42)
on the V-lattice. On the other hand, the entropy of a linelike object (essentially the entropy of a random walk) is a constant of $`O(1)`$ times the line length. Thus, center vortices are stable saddlepoints of the long range effective action whose entropy per unit length greatly exceeds their action per unit length. These objects therefore percolate throughout the lattice volume, and confine color charge in any half-integer group representation.
This picture has been derived in a particular unitary gauge. A natural question is whether the saddlepoints of the effective action would be qualitatively different had we chosen to gauge-fix, instead of $`\stackrel{}{e}_{11}`$ and $`\stackrel{}{e}_{44}`$, some other plaquette variables (or combination of plaquette variables) on the 2-cubes. Although an analysis of all possible unitary gauge choices is beyond us at present, it is easy to see that center vortices must be stable saddlepoints in a very large class of gauges. We first note that, by definition, any maximal unitary gauge must completely determine the minimal action configuration of the $`V_\mu `$ fields, up to residual $`Z_2`$ gauge transformations. Then a sufficient condition for center vortex stability is simply that the classical ground state of $`S_{eff}[V]`$ has the form of a pure gauge $`V_\mu (x)=g(x)g^1(x+\widehat{\mu })`$, and that this is also a maximum of the center-invariant $`S_{link}`$ part of the effective action. In that case, the effective action must have stable thin vortex solutions at large $`L`$. This is because the stable fluctuation modes around a thin vortex, associated with the center-insensitive $`S_{link}`$ term, will overwhelm (at sufficiently large $`L`$) any unstable modes associated with $`S_{plaq}`$. The vortex action at a given $`L`$ will always have the value shown in eq. (42) above, so the entropy of the configuration will exceed the action at strong couplings.
Our findings for the strong-coupling theory do not, however, prove that vortices also dominate the vacuum at weaker couplings; the strong and weak coupling regimes are separated by a roughening phase transition, and this transition prevents a simple extrapolation from one regime to the other. The result is, nonetheless, significant for continuum physics in two ways: First, it supports the very general argument that if the asymptotic quark potential is sensitive only to $`N`$-ality, then the confining field configurations must be center vortices. Secondly, it provides an explicit illustration of how center vortices are stabilized, via color-screening (center-invariant) terms in the long-range effective action.
Although strong-coupling methods only apply at strong couplings, the general approach we have advocated here should extend, at least in principle, to weak-coupling lattice gauge theory. The central idea is that if we want to extract a *local* long-range effective action from the Wilson action, then it is necessary to include composite operators, transforming like adjoint matter fields, in the derivation. With this motivation we define gluelump operators
$$G_M[x;U]=\underset{C_x}{}a_M(C_x)\underset{lC_x}{}U_l$$
(43)
which, coupled to a static adjoint source at site $`x`$, create gauge-invariant eigenstates of the appropriate transfer matrix. The $`C_x`$ are paths on the lattice beginning and ending at $`x`$.<sup>5</sup><sup>5</sup>5At strong couplings, the only relevant paths $`C_x`$ run around single plaquettes. The index $`M`$ specifies the time (i.e. worldline) direction of the static source, and any other (e.g. spin) degeneracies. The transformation from a pure gauge theory on a fine lattice, to a theory of gauge fields $`V_\mu `$ coupled to gluelump fields $`H_M`$ on a coarse lattice, could then be accomplished as follows:
$$\mathrm{exp}\left[S_{eff}[V,H]\right]=DU\underset{l^{}}{}\delta \left(V_l^{}Q_l^{}[U]\right)\underset{M,x^{}}{}\delta \left(H_M(x^{})G_M[x^{};U]\right)e^{S_W[U]}$$
(44)
where $`x^{},l^{}`$ denote sites and links on the coarse lattice. The expression $`Q_l^{}[U]`$ represents a suitable “fat link” function, i.e. a superposition of Wilson lines on the fine lattice which run between sites bounding link $`l^{}`$ on the coarse lattice. The constraints imposed by the delta-functions can be softened by replacing delta-functions with exponentials, as in the “perfect action” approach . The end result of this procedure will be an effective long-range action consisting of gauge fields coupled to a set of adjoint Higgs-like fields. Possibly this scheme can be implemented numerically at moderately weak couplings, along the lines of the Monte Carlo renormalization-group.
In the strong-coupling analysis carried out in this article, we have seen how color-screening, center-invariant “Higgs” terms predominate in the long-range effective action beyond the adjoint string-breaking scale, and stabilize center vortex configurations. The entropy of these configurations exceeds the cost in action, and vortex configurations percolate throughout the lattice. We think it likely that these important features are not specific to strong couplings, and also characterize the effective action of lattice QCD in the continuum limit.
###### Acknowledgments.
Our research is supported in part by Fonds zur Förderung der Wissenschaftlichen Forschung P13997-PHY (M.F.), the U.S. Department of Energy under Grant No. DE-FG03-92ER40711 (J.G.), the “Action Austria-Slovakia: Cooperation in Science and Education” (Project No. 30s12) and the Slovak Grant Agency for Science, Grant No. 2/7119/2000 (Š.O.).
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# A multifractal random walk
## Abstract
We introduce a class of multifractal processes, referred to as Multifractal Random Walks (MRWs). To our knowledge, it is the first multifractal processes with continuous dilation invariance properties and stationary increments. MRWs are very attractive alternative processes to classical cascade-like multifractal models since they do not involve any particular scale ratio. The MRWs are indexed by few parameters that are shown to control in a very direct way the multifractal spectrum and the correlation structure of the increments. We briefly explain how, in the same way, one can build stationary multifractal processes or positive random measures.
Multifractal models have been used to account for scale invariance properties of various objects in very different domains ranging from the energy dissipation or the velocity field in turbulent flows to financial data. The scale invariance properties of a deterministic fractal function $`f(t)`$ are generally characterized by the exponents $`\zeta _q`$ which govern the power law scaling of the absolute moments of its fluctuations, i.e.,
$$m(q,l)=K_ql^{\zeta _q},$$
(1)
where, for instance, one can choose $`m(q,l)=_t|f(t+l)f(t)|^q`$. When the exponents $`\zeta _q`$ are linear in $`q`$, a single scaling exponent $`H`$ is involved. One has $`\zeta _q=qH`$ and $`f(t)`$ is said to be monofractal. If the function $`\zeta _q`$ is no longer linear in $`q`$, $`f(t)`$ is said to be multifractal. In the case of a stochastic process $`X(t)`$ with stationary increments, these definitions are naturally extended using
$$m(q,l)=E\left(|\delta _lX(t)|^q\right)=E\left(|X(t+l)X(t)|^q\right),$$
(2)
where $`E`$ stands for the expectation. Some very popular monofractal stochastic processes are the so-called self-similar processes bTaqq . They are defined as processes $`X(t)`$ which have stationary increments and which verify (in law)
$$\delta _{\lambda l}X(t)=\lambda ^H\delta _lX(t),l,\lambda >0.$$
(3)
Widely used examples of such processes are fractional Brownian motions (fBm) and Levy walks. One reason for their success is that, as it is generally the case in experimental time-series, they do not involve any particular scale ratio (i.e., there is no constraint on $`l`$ or $`\lambda `$ in Eq. (3)). In the same spirit, one can try to build multifractal processes which do not involve any particular scale ratio. A common approach originally proposed by several authors in the field of fully developed turbulence nov ; sl ; FP ; dg ; castaing , has been to describe such processes in terms of differential equations, in the scale domain, describing the cascading process that rules how the fluctuations evolves when going from coarse to fine scales. One can state that the fluctuations at scales $`l`$ and $`\lambda l`$ ($`\lambda <1`$) are related (for fixed $`t`$) through the infinitesimal ($`\lambda =1\eta `$ with $`\eta <<1`$) cascading rule
$$\delta _{\lambda l}X(t)=W_\lambda \delta _lX(t)$$
(4)
where $`W_\lambda `$ is a stochastic variable which depends only on $`\lambda `$. Let us note that this latter equation can be simply seen as a generalization of Eq. (3) with $`H`$ being stochastic. Since Eq. (4) can be iterated, it implicitely imposes the random variable $`W_\lambda `$ to have a log infinitely divisible law. However, according to our knowledge, nobody has succeeded in building effectively such processes yet, mainly because of the peculiar constraints in the time-scale half-plane. The integral equation corresponding to this infinitesimal approach has been proposed by Castaing et al. castaing . It relates the probability density function (pdf) $`P_l(\delta X)`$ of $`\delta _lX`$ to the pdf $`G_{\lambda l,l}`$ of $`\mathrm{ln}W_\lambda `$ :
$$P_{\lambda l}(\delta X)=G_{\lambda l,l}(u)e^uP_l(e^u\delta X)𝑑u.$$
(5)
The self-similarity kernel $`G_{\lambda l,l}`$ satisfies the same iterative rule as $`W_\lambda `$ which implies that its Fourier transform is of the form $`\widehat{G}_{\lambda l,l}(k)=\widehat{G}^\lambda (k)`$. Thus one can easily show that the $`q`$ order absolute moments at scale $`l`$ scales like
$$m(q,l)=\widehat{G}_{l,L}(iq)m(q,L)=m(q,L)\left(\frac{l}{L}\right)^{F(iq)},$$
(6)
where $`F=\mathrm{ln}\widehat{G}`$ refers to the cumulant generating function of $`\mathrm{ln}W`$ castaing ; amr . Thus, identifying this latter equation with Eq. (1), one finds $`\zeta (q)=F(iq)`$.
In the case of self-similar processes of exponent $`H`$, one easily gets that the kernel is a dirac function $`G_{l,L}(u)=\delta (uH\mathrm{ln}(l/L))`$ and consequently $`\zeta _q=qH`$. The simplest non-linear (i.e., multifractal) case is the so-called log-normal model that corresponds to a parabolic $`\zeta _q`$ and a Gaussian kernel. Let us note that if $`\zeta _q`$ is non linear, a simple concavity argument shows that Eq. (1) cannot hold for all $`l`$ in $`]0,+\mathrm{}[`$. Multiplicative cascading processes man62 ; kp ; ms ; hent ; wcasc consist in writing Eq. (4) starting from some “coarse” scale $`l=L`$ and then iterating it towards finer scales using an arbitrary fixed scale ratio (e.g., $`\lambda =1/2`$). Such processes can be contructed rigorously using, for instance, orthonormal wavelet bases wcasc . However, these processes have fundamental drawbacks: they do not lead to stationary increments and they do not have continuous dilation invariance properties. Indeed, they involve a particular arbitrary scale ratio, i.e., Eq (1) holds only for the discrete scales $`l_n=\lambda ^nL`$.
The goal of this paper is to build a multifractal process $`X(t)`$, referred to as a Multifractal Random Walk (MRW), with stationary increments and such that Eq. (1) holds for all $`lL`$. We first build a discretized version $`X_{\mathrm{\Delta }t}(t=k\mathrm{\Delta }t)`$ of this process. Let us note that the theoretical issue whether the limit process $`X(t)=lim_{\mathrm{\Delta }t0}X_{\mathrm{\Delta }t}(t)`$ is well defined will be addressed in a forthcoming paper. In this paper, we explain how it is built and prove that different quantities ($`q`$ order moments, increment correlation,…) converge, when $`\mathrm{\Delta }t0`$.
Writing Eq. (4) at the smallest scale suggests that a good candidate might be such that $`\delta _{\mathrm{\Delta }t}X_{\mathrm{\Delta }t}(k\mathrm{\Delta }t)=ϵ_{\mathrm{\Delta }t}[k]W_{\mathrm{\Delta }t}[k]`$ where $`ϵ_{\mathrm{\Delta }t}`$ is a Gaussian variable and $`W_{\mathrm{\Delta }t}[k]=e^{\omega _{\mathrm{\Delta }t}[k]}`$ is a log normal variable, i.e.,
$$X_{\mathrm{\Delta }t}(t)=\underset{k=1}{\overset{t/\mathrm{\Delta }t}{}}\delta _{\mathrm{\Delta }t}X_{\mathrm{\Delta }t}(t)=\underset{k=1}{\overset{t/\mathrm{\Delta }t}{}}ϵ_{\mathrm{\Delta }t}[k]e^{\omega _{\mathrm{\Delta }t}[k]},$$
(7)
with $`X_{\mathrm{\Delta }t}(0)=0`$ and $`t=k\mathrm{\Delta }t`$. Moreover, we will choose $`ϵ_{\mathrm{\Delta }t}`$ and $`\omega _{\mathrm{\Delta }t}`$ to be decorrelated and $`ϵ_{\mathrm{\Delta }t}`$ to be a white noise of variance $`\sigma ^2\mathrm{\Delta }t`$. Obviously, we need to correlate the $`\omega _{\mathrm{\Delta }t}[k]`$’s otherwise $`X_{\mathrm{\Delta }t}`$ would simply converge towards a Brownian motion. Since, in the case of cascade-like processes, it has been shown wcasc ; ams ; abmm that the covariance of the logarithm of the increments decreases logarithmically, it seems natural to choose
$$Cov(\omega _{\mathrm{\Delta }t}[k_1],\omega _{\mathrm{\Delta }t}[k_2])=\lambda ^2\mathrm{ln}\rho _{\mathrm{\Delta }t}[|k_1k_2|],$$
(8)
with
$$\rho _{\mathrm{\Delta }t}[k]=\{\begin{array}{cc}\frac{L}{(|k|+1)\mathrm{\Delta }t}\hfill & \text{for}|k|L/\mathrm{\Delta }t1\hfill \\ 1\hfill & \text{otherwise}\hfill \end{array},$$
(9)
i.e., the $`\omega _{\mathrm{\Delta }t}`$ are correlated up to a distance of $`L`$ and their variance $`\lambda ^2\mathrm{ln}(L/\mathrm{\Delta }t)`$ goes to $`+\mathrm{}`$ when $`\mathrm{\Delta }t`$ goes to 0. For the variance of $`X_{\mathrm{\Delta }t}`$ to converge, a quick computation shows that we need to choose
$$E\left(\omega _{\mathrm{\Delta }t}[k]\right)=rVar\left(\omega _{\mathrm{\Delta }t}[k]\right)=r\lambda ^2\mathrm{ln}(L/\mathrm{\Delta }t),$$
(10)
with $`r=1`$ (this value will be changed later) for which we find $`Var(X(t))=\sigma ^2t`$.
Let us compute the moments of the increments of the MRW $`X(t)`$. Using the definition of $`X_{\mathrm{\Delta }t}(t)`$ one gets
$`E(X_{\mathrm{\Delta }t}(t_1)\mathrm{}X_{\mathrm{\Delta }t}(t_m))={\displaystyle \underset{k_1=1}{\overset{t_1/\mathrm{\Delta }t}{}}}\mathrm{}{\displaystyle \underset{k_m=1}{\overset{t_m/\mathrm{\Delta }t}{}}}`$
$`E(ϵ_{\mathrm{\Delta }t}[k_1]\mathrm{}ϵ_{\mathrm{\Delta }t}[k_m])E\left(e^{\omega _{\mathrm{\Delta }t}[k_1]+\mathrm{}+\omega _{\mathrm{\Delta }t}[k_m]}\right).`$
Since $`ϵ_{\mathrm{\Delta }t}`$ is a 0 mean Gaussian process, this expression is 0 if $`m`$ is odd. Let $`m=2p`$. Since the $`ϵ_{\delta t}[k]`$’s are $`\delta `$-correlated Gaussian variables, one shows that the previous expression reduces to
$$\frac{\sigma ^{2p}}{2^pp!}\underset{𝒮S_{2p}}{}\underset{k_1=1}{\overset{\left(t_{𝒮(1)}t_{𝒮(2)}\right)/\mathrm{\Delta }t}{}}\mathrm{}\mathrm{}.\underset{k_p=1}{\overset{\left(t_{𝒮(2p1)}t_{𝒮(2p)}\right)/\mathrm{\Delta }t}{}}E\left(e^{2_{j=1}^p\omega _{\mathrm{\Delta }t}[k_j]}\right)\mathrm{\Delta }t^p,$$
where $`ab`$ refers to the minimum of $`a`$ and $`b`$ and $`S_{2p}`$ to the set of the permutations on $`\{1,\mathrm{},2p\}`$. On the other hand, we have $`E\left(e^{2_{j=1}^p\omega _{\mathrm{\Delta }t}[k_j]}\right)=_{i<j}\rho [k_ik_j]^{4\lambda ^2}.`$ Then, when $`\mathrm{\Delta }t0`$, the general expression of the moments is
$`E(X(t_1)\mathrm{}X(t_{2p}))={\displaystyle \frac{\sigma ^{2p}}{2^pp!}}{\displaystyle \underset{𝒮S_{2p}}{}}{\displaystyle _0^{t_{𝒮(1)}t_{𝒮(2)}}}𝑑u_1`$ (11)
$`\mathrm{}{\displaystyle _0^{t_{𝒮(2p1)}t_{𝒮(2p)}}}𝑑u_p{\displaystyle \underset{i<j}{}}\rho (u_iu_j)^{4\lambda ^2},`$
where $`\rho (t)=lim_{\mathrm{\Delta }t0}\rho _{\mathrm{\Delta }t}[t/\mathrm{\Delta }t]`$. In the special case $`t_1=t_2=\mathrm{}=t_{2p}=l`$, a simple scaling argument leads to the continuous dilation invariance property
$$m(2p,l)=K_{2p}\left(\frac{l}{L}\right)^{p2p(p1)\lambda ^2},lL,$$
(12)
where we have denoted the prefactor
$$K_{2p}=L^p\sigma ^{2p}(2p1)!!_0^1𝑑u_1\mathrm{}_0^1𝑑u_p\underset{i<j}{}|u_iu_j|^{4\lambda ^2}.$$
By analytical continuation, we thus obtain the following $`\zeta _q`$ spectrum
$$\zeta _q=(qq(q2)\lambda ^2)/2.$$
(13)
We have illustrated this scaling behavior in fig. 1. Thus, the MRW $`X(t)`$ is a multifractal process with stationary increments and with continuous dilation invariance properties up to the scale $`L`$. Let us note that above this scale ($`l>>L`$), one gets from Eq. (11) that $`\zeta _q=q/2`$, i.e., the process scales like a simple Brownian motion, as if $`\omega `$ was not correlated, though, of course, $`X(t)`$ is not Gaussian. Indeed, $`K_{2p}`$ is nothing but the moment of order $`2p`$ of the random variable $`X(L)`$ and is infinite for large $`p`$’s (depending on $`\lambda `$). Actually, one can show that $`\zeta _{2p}0K_{2p}=+\mathrm{}`$. Consequently, the pdf of $`X(L)`$ has fat tails. As illustrated in fig. 2, Eq. (5) accounts very well for the evolution of the pdf of the increments. One shows that the smaller the scale $`l`$, the fatter the tails of the pdf of $`\delta _lX(t)`$.
Let us study the correlation structure of the increments of $`X(t)`$. Since $`\zeta _2=1`$, one can prove that they are decorrelated (though not independant). Let
$$C_{2p}(l,\tau )=<|\delta _\tau X(l)|^{2p}|\delta _\tau X(0)|^{2p}>,$$
(14)
with $`\tau <l`$. Using the same kind of arguments as above, one can show that
$`C_{2p}(l,\tau )=(\sigma ^{2p}(2p1)!!)^2{\displaystyle _l^{l+\tau }}𝑑u_1\mathrm{}{\displaystyle _l^{l+\tau }}𝑑u_p`$
$`{\displaystyle _0^\tau }𝑑u_{p+1}\mathrm{}{\displaystyle _0^\tau }𝑑u_{2p}{\displaystyle \underset{1i<j2p}{}}\rho (u_iu_j)^{4\lambda ^2}.`$ (15)
A straightforward argument then shows that
$$K_{2p}^2\frac{(\tau /L)^{2\zeta _{2p}}}{((l+\tau )/L)^{4\lambda ^2p^2}}C_{2p}(l,\tau )K_{2p}^2\frac{(\tau /L)^{2\zeta _{2p}}}{((l\tau )/L)^{4\lambda ^2p^2}},$$
and consequently for $`\tau <<l`$ fixed, using analytical continuation one expects $`C_q(l,\tau )`$ to scale like
$$C_q(l,\tau )K_q^2\left(\frac{\tau }{L}\right)^{2\zeta _q}\left(\frac{l}{L}\right)^{łambda^2q^2}.$$
(16)
This behavior is illustrated in fig. 3.
¿From the behavior of $`C_q`$ when $`q0`$, we can obtain using Eq. (16) that the covariance of the logarithm of the increments at scale $`\tau `$ and lag $`l`$ behaves (for $`\tau <<l`$) like
$$C^{(\mathrm{ln})}(l,\tau )\lambda ^2\mathrm{ln}\left(\frac{l}{L}\right).$$
(17)
Thus, this correlation reflects the correlation of the $`\omega _{\mathrm{\Delta }t}`$ process and is the same as observed in Refs wcasc ; ams ; abmm for the cascade models. This behavior is checked in fig. 4.
Finally, let us note that, one can built MRWs with correlated increments by just replacing the white noise $`ϵ_{\mathrm{\Delta }t}`$ by a fractional Gaussian noise (fGn)
$$ϵ_{\mathrm{\Delta }t}^{(H)}[k]=B_H((k+1)\mathrm{\Delta }t)B_H(k\mathrm{\Delta }t),$$
(18)
where $`B_H(t)`$ is a fBm with the scaling exponent $`H`$ and of variance $`\sigma ^2t^{2H}`$, and choosing $`r=1/2`$ in Eq. (10). Then, one can show (after tedious but straightforward computations) that the spectrum of the MRW $`X^{(H)}(t)`$ is
$$\zeta _q^{(H)}=qHq(q1)\lambda ^2/2,$$
(19)
($`\zeta _q^{(H)}=qH`$ at scales $`>>L`$) and consequently the MRW has correlated increments. Such a construction is illustrated in fig. 1 with $`H=2/3`$. Since $`H>1/2`$ it leads to a process which is more regular than the one previously built.
To summarize, we have built the MRWs, a class of multifractal processes, with stationary increments and continuous dilation invariance. ¿From a theoretical point of view, MRW can be seen as the simplest model that contains the main ingredients for multifractality, namely, logarithmic decaying correlation of the logarithms of amplitude fluctuations. Moreover, they involve very few parameters, mainly, the correlation length $`L`$, the intermittency parameter $`\lambda ^2`$, the variance $`\sigma ^2`$ and the roughness exponent $`H`$. They all can be easily estimated using the $`\zeta _q`$ spectrum and the increment correlation structure. We do believe that they should be very helpful in all the fields where multiscaling is observed. MRWs have already been proved successful for modelling financial data fmuzy . In this framework, we have shown that one can easily build multivariate MRWs. Actually, the construction of MRWs, as presented in this paper, can be used as a general framework in which one can easily build other classes of processes in order to match some specific experimental situations. For instance, a stationary MRW can be obtained by just adding a friction $`\gamma >0`$, i.e., $`X_{\mathrm{\Delta }_t}[k]=(1\gamma )X_{\mathrm{\Delta }_t}[k1]+ϵ_{\mathrm{\Delta }t}[k]e^{\omega _{\mathrm{\Delta }t}[k]}`$. One can build a strictly increasing MRW (and consequently a stochastic positive multifractal measure) by just setting $`ϵ_{\mathrm{\Delta }t}=\mathrm{\Delta }t`$ in Eq. (7) and use it as a multifractal time for subordinating a monofractal process (such as an fBm). One can also use other laws than the (log-)normal for $`ϵ`$ and/or $`\omega `$. Another interesting point concerns the problem of the existence of a limit ($`\mathrm{\Delta }t0`$) stochastic process and on the development of a new stochastic calculus associated to this process. All these prospects will be addressed in forthcoming studies.
We acknowledge Alain Arneodo for interesting discussions.
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# GEOMETRIC CONDITIONS ON THE TYPE OF MATTER DETERMINING THE FLAT BEHAVIOR OF THE ROTATIONAL CURVES IN GALAXIES
## I Introduction
One of the most important achievements of the present cosmology is doubtless the complete identifications and the accounting of the amounts of the different types of matter and energy which are present in the Universe (for an excellent review see for example ). Essentially, the present components of the Universe are composed by matter and vacuum energy $`\mathrm{\Omega }_0=\mathrm{\Omega }_M+\mathrm{\Omega }_\mathrm{\Lambda }`$ . Furthermore, there is a very good evidence that the Universe is flat. That evidence comes either from the theory, where the most accepted model for the early Universe is inflation, as well as from observational data, which implies that $`\mathrm{\Omega }_0=1\pm 0.12`$ (see for example ). The mass of the galaxy clusters is perhaps the most reliable way for determining the matter component $`\mathrm{\Omega }_M`$. Observations indicate that $`\mathrm{\Omega }_M0.3`$ , however, the main visible components of $`\mathrm{\Omega }_M`$, baryons, neutrinos, form a very small fraction of $`\mathrm{\Omega }_M`$. Observations indicate that stars and dust (baryons) represents something like $`5\%`$ of the whole matter of the Universe. In other words, $`\mathrm{\Omega }_M\mathrm{\Omega }_b+\mathrm{\Omega }_{DM}0.05+\mathrm{\Omega }_{DM}`$, where $`\mathrm{\Omega }_{DM}`$ represents the dark matter part of the matter contributions which has to have a value of $`\mathrm{\Omega }_{DM}0.25`$. Notice that the given value of the amount of baryonic matter is in concordance with the limits imposed by nucleosynthesis (see for example ). A greater amount of baryonic matter could change the predicted values of primordial H and <sup>4</sup>He in the standard model of cosmology which coincide very well with observations.
The existence of dark matter in the Universe has been firmly established by astronomical observations at very different length-scales, ranging from the local galaxies to clusters of galaxies. The standard way to notice this need for dark matter comes within the framework of mechanics: A large fraction of the mass needed to produce the observed dynamical effects in all these very different systems, is not seen. This puzzle has stimulated the exploration of several proposals, and very imaginative explanations have been put forward, from exotic matter like supersymmetric particles to non relativistic modifications of Newtonian dynamics and non-linear general relativistic theories . Above this, it is believed that this dark matter is such that interacts very weakly with ordinary matter, which makes it very difficult to detect by other means other than by their gravitational effects on the baryonic matter as it is well stablished for the cold dark matter scenario . The bottom line is that one of the most important components of the density of the Universe, the dark matter, still eludes us, and the question remains open: Which is the nature of the dark matter component?
At the galactic scale, the problem is clearly posed: The measurements of rotation curves (RC) in galaxies show that the coplanar orbital motion of gas in the outer parts of these galaxies keeps a more or less constant velocity up to several luminous radii . The discrepancy arises when one applies the usual Newtonian dynamics to the observed luminous matter and gas, since then, the circular velocity should decrease as we move outwards. The most widely accepted explanation is that there exists a spherical halo of dark matter, its nature being unknown, which surrounds the galaxy and account for the missing mass needed to produce the flat behavior of the RC.
The main goal of this work is to study the dark matter problem in spiral galaxies, using a fully relativistic approach, continuing with the work started by Matos and Guzmán , where they made a preliminary dynamical analysis in the context of spiral galaxies. We apply a deductive method: Starting from a reasonable general space-time, we deduce, in terms of arbitrary metric coefficients the expression for the tangential velocity of test objects following a circular stable geodesic motion in the equatorial plane. We impose next the condition, observed in hundreds of galaxies, that such tangential velocity be radii independent, and obtain a constraint equation among the metric coefficients, Eq.(28). Arriving in this way to an iff condition: In an axisymmetric static space-time, the tangential velocity of test particles at the equatorial plane is radii independent iff the metric coefficients satisfied Eq.(28). Furthermore, for the static case, the constraint equation can be solved, leaving the space-time essentially with only one independent metric coefficient, thus determining very narrowly the type of space-time which can have a geometry such that the tangential velocity of the geodesics of test particles moving in stable circular equatorial orbits, be radii independent. With the geometry thus fixed, we compute the Einstein tensor and equate it to and arbitrary stress energy tensor, in order to determined the type of energy-matter which could produce such a geometry. We are able to deduce a constraint equation among the components of the stress energy tensor, Eq.(50). Using that constraint equation, Eq.(50), we test several well known types of matter, which have been proposed as dark matter candidates and are able to impose further restrictions on most of them as possible candidates for the dark matter in the region where it is observed the mentioned behavior on the tangential velocity. It is stimulating that the gravitational physics has been developed to such degree in which we can actually follow Sherlock Holmes’ maxim: “…when you have excluded the impossible, whatever remains, however improbable, must be the truth” .
The work is composed as follows: In section 2 we determine the geometry of axisymmetric static space times allowing a tangential velocity with a magnitude independent of the radius. In section 3, we work with the Einstein equations for the geometry thus determined and for a general stress energy tensor, obtaining a constraint for the components of such stress energy tensor, and we test four types of matter into that constraint, being able to restrict most of them, among which is the perfect fluid. Finally, in the conclusion we discuss our results and propose a coupling of types of matter, which could be the one we are looking for in order to determine the nature of the dark matter. We also present this analysis for the spherical static space time in appendix V and for the axisymmetric stationary case in appendix V.
## II Form of the line element
In this section we study the conditions which the flatness of the tangential velocity of the RC, imposes on the metric coefficients. We want to stress the fact that the results presented in this section are independent of type of energy-matter tensor present in the space-time and curving it. It is a purely geometric analysis.
As mentioned above, observational data show that the galaxies are composed by almost 90% of dark matter. Thus we can suppose that luminous matter does not contribute in a determining way to the total energy density of the galaxy, at least in the region where the flatness of the RC is observed. Consequently, we consider that the dark matter will be the main contributor to the dynamics, and we will treat the observed luminous matter as a test fluid, that’s it, in this approximation we will neglect the contribution of the luminous matter to the curvature, i. e., to the dynamics. Also, it is reasonable to suppose that the halo of dark matter is symmetric with respect to the rotation axis of the galaxy, thus we take the symmetry of the space-time as axially symmetric. Furthermore, the observations allow us to take the space-time as stationary as well. Thus, the most general reasonable space-time which we can study is an axisymmetric stationary one. The line-element of such space-time, given in the Papapetrou form is:
$$ds^2=e^{2\psi }(dt+\omega d\phi )^2+e^{2\psi }[e^{2\gamma }(d\rho ^2+dz^2)+\mu ^2d\phi ^2],$$
(1)
where $`\psi ,\omega ,\gamma `$, and $`\mu `$, are functions of $`(\rho ,z)`$.
We will derive the geodesic equations in the equatorial plane, that is for $`z=\dot{z}=0`$, where dot stands for the derivative with respect to the proper time, $`\tau `$. Then, we will study the constrains imposed for circular geodesics on the energy and angular momentum of particles in such orbits, and obtain an expression for the tangential velocity for the particles moving along those geodesics, described in terms of the metric coefficients. Finally, we impose the condition that such tangential velocity be radii independent and derive the restriction that then has to be satisfied among the gravitational coefficients. In this section we derive the expressions for a static space time and in the Appendix V we do so for the stationary one. We split the presentation in this way not only for the sake of clarity, but from the fact that the static case is also quite realistic, considering that the observed velocity of the stars orbiting a galaxy in the region of interest is quite non relativistic (of the order of 230 Km per s), thus we can infer that the space-time is not very rapidly rotating.
The Lagrangian for a test particle traveling on the static space time ($`\omega =0`$) described by (1) is given by:
$$2=e^{2\psi }\dot{t}^2+e^{2\psi }[e^{2\gamma }(\dot{\rho }^2+\dot{z}^2)+\mu ^2\dot{\phi }^2],$$
(2)
thus, the associated canonical momenta, $`p_{x^a}=\frac{}{\dot{x^a}}`$, are:
$`p_t=E`$ $`=`$ $`e^{2\psi }\dot{t},`$ (3)
$`p_\phi =L`$ $`=`$ $`\mu ^2e^{2\psi }\dot{\phi },`$ (4)
$`p_\rho `$ $`=`$ $`e^{2(\psi \gamma )}\dot{\rho },`$ (5)
$`p_z`$ $`=`$ $`e^{2(\psi \gamma )}\dot{z},`$ (6)
where $`E`$, and $`L`$, are constants of motion for each geodesic, a fact which comes from the symmetries of the space-time analyzed. As there is no explicit dependence on time, $`t`$, the Hamiltonian, $`=p_a\dot{x^a}`$, is another conserved quantity, which we normalized to be equal to minus one half for time-like geodesics. Also, we restrict the motion to be at the equatorial plane, thus $`\dot{z}=0`$. In this way, we obtain the following equation for the radial geodesic motion:
$$\dot{\rho }^2e^{2(\psi \gamma )}[E\dot{t}L\dot{\phi }1]=0.$$
(7)
In order to have stable circular motion, which is the motion we are interested in, we have to satisfy three conditions:
i) $`\dot{\rho }=0`$, and
ii)$`\frac{V(\rho )}{\rho }=0`$, where $`V(\rho )=e^{2(\psi \gamma )}[E\dot{t}L\dot{\phi }1]`$.
iii)$`\frac{^2V(\rho )}{\rho ^2}|_{extr}>0`$, in order to have a minimum.
With these conditions, from Eq.(7), we obtain a set of two equations constraining the motion to be circular extrema in the equatorial plane:
$`E\dot{t}L\dot{\phi }1`$ $`=`$ $`0,`$ (8)
$`{\displaystyle \frac{}{\rho }}\left(e^{2(\psi \gamma )}[E\dot{t}L\dot{\phi }1]\right)`$ $`=`$ $`0.`$ (9)
¿From Eq.(6), we can express $`\dot{t}`$, and $`\dot{\phi }`$ in terms of $`E,L`$, and the metric coefficients as
$`\dot{t}`$ $`=`$ $`e^{2\psi }E,`$ (10)
$`\dot{\phi }`$ $`=`$ $`{\displaystyle \frac{e^{2\psi }}{\mu ^2}}L.`$ (11)
Using these equations in the constraints ones and recalling that $`E`$ and $`L`$ are constants for each circular orbit, after some rearranging, we arrive to the following equations:
$`\mu ^2e^{2\psi }(1e^{2\psi }E^2)+L^2`$ $`=`$ $`0,`$ (12)
$`(e^{2\psi })_\rho E^2+\left({\displaystyle \frac{e^{2\psi }}{\mu ^2}}\right)_\rho L^2`$ $`=`$ $`0,`$ (13)
where the subindex stands for derivative with respect to $`\rho `$. Solving for $`E`$ and $`L`$, we obtain:
$`E`$ $`=`$ $`e^\psi \sqrt{{\displaystyle \frac{\frac{\mu _\rho }{\mu }\psi _\rho }{\frac{\mu _\rho }{\mu }2\psi _\rho }}},`$ (14)
$`L`$ $`=`$ $`\mu e^\psi \sqrt{{\displaystyle \frac{\psi _\rho }{\frac{\mu _\rho }{\mu }2\psi _\rho }}}.`$ (15)
The second derivative of the potential $`V(\rho )`$ evaluated at the extreme, in this case means evaluate at the values of $`E`$ and $`L`$ which constraint the motion to be circular and extrema, is given by:
$$V_{\rho \rho }|_{extr}=\frac{2e^{2(\psi \gamma )}}{\frac{\mu _\rho }{\mu }2\psi _\rho }\left(\frac{\mu _\rho }{\mu }\psi _{\rho \rho }\frac{\mu _{\rho \rho }}{\mu }\psi _\rho +4\psi _{\rho }^{}{}_{}{}^{3}6\frac{\mu _\rho }{\mu }\psi _{\rho }^{}{}_{}{}^{2}+3\left(\frac{\mu _\rho }{\mu }\right)^2\psi _\rho \right).$$
(16)
We can now obtain an expression for the angular velocity of a test particle, $`\mathrm{\Omega }`$, moving in a circular motion in the orbital plane, in terms of the metric coefficients, recalling that
$$\mathrm{\Omega }=\frac{d\phi }{dt}=\frac{\dot{\phi }}{\dot{t}},$$
(17)
thus, using Eqs.(11), and (15), in this last equation for the angular velocity, we obtain that:
$$\mathrm{\Omega }=\frac{e^{2\psi }}{\mu }\sqrt{\frac{\psi _\rho }{\frac{\mu _\rho }{\mu }\psi _\rho }}.$$
(18)
Finally, in order to express the tangential velocity of the test particles in circular motion in the equatorial plane, in terms of the metric coefficients, following Chandrasekhar , we rewrite the line element given in Eq.(1) as:
$$ds^2=e^{2\psi }dt^2+e^{2\psi }\mu ^2d\phi ^2+e^{2(\psi \gamma )}(d\rho ^2+dz^2),$$
(19)
thus, in terms of the proper time, $`d\tau ^2=ds^2`$, we have that
$`d\tau ^2`$ $`=`$ $`e^{2\psi }dt^2[1e^{4\psi }\mu ^2\left({\displaystyle \frac{d\phi }{dt}}\right)^2+`$ (21)
$`e^{2\gamma }e^{4\psi }(\left({\displaystyle \frac{d\rho }{dt}}\right)^2+\left({\displaystyle \frac{dz}{dt}}\right)^2)],`$
from which we can write that
$$1=e^{2\psi }u_{}^{0}{}_{}{}^{2}[1v^2],$$
(22)
where $`u^0=\frac{dt}{d\tau }`$ is the usual time component of the four velocity, and a definition of the spatial velocity, $`v^2`$, comes out naturally in this way.
$$v^2=e^{4\psi }\mu ^2\left(\frac{d\phi }{dt}\right)^2+e^{2\gamma }e^{4\psi }\left(\left(\frac{d\rho }{dt}\right)^2+\left(\frac{dz}{dt}\right)^2\right).$$
(23)
This spatial velocity is the 3-velocity of a particle measured with respect to an orthonormal reference system (see section 52 of ), thus has components:
$$v^2=v_{}^{(\phi )}{}_{}{}^{2}+v_{}^{(\rho )}{}_{}{}^{2}+v_{}^{(z)}{}_{}{}^{2}.$$
(24)
¿From these last two expressions we obtain for the $`\phi `$component the spatial velocity:
$$v^{(\phi )}=e^{2\psi }\mu \mathrm{\Omega },$$
(25)
and substituting $`\mathrm{\Omega }`$ from Eq.(18), we finally obtain an expression for the tangential velocity of a test particle in stable circular motion, in terms of the metric coefficients of the general line element given by Eq.(1), such tangential velocity has the form:
$$v^{(\phi )}=\sqrt{\frac{\psi _\rho }{\frac{\mu _\rho }{\mu }\psi _\rho }}.$$
(26)
It was our goal to obtain this expression for the tangential velocity for a general axisymmetric static space time, and to be able to describe it in terms of the metric coefficients alone, because now we can impose conditions on this tangential velocity, and deduce a constraint equation among the metric coefficients, which has to be satisfied in order to fulfill the condition imposed on the velocity. In particular, the tangential velocity for circular trajectories in each orbit is constant, that is $`v^{(\phi )}{}_{\rho }{}^{}=0`$, thus $`v^{(\phi )}=v_c^{(\phi )}`$, with $`v_c^{(\phi )}`$ a constant, representing the value of the velocity, from Eq. (26), we have that:
$$\frac{\mu _\rho }{\mu }=\frac{1+v_{c}^{(\phi )}{}_{}{}^{2}}{v_{c}^{(\phi )}{}_{}{}^{2}}\psi _\rho .$$
(27)
Finally, with respect to the $`z`$-motion, considering that at the equatorial plane not only $`\dot{z}=0`$, but also that $`\ddot{z}=0`$, that is that the forces above and below the plane cancel out, from the geodesic $`z`$-equation, using Eq.(27), we obtain that this relation among the metric coefficients must hold for the derivatives with respect to $`z`$ as well:
$$\frac{\mu _z}{\mu }=\frac{1+v_{c}^{(\phi )}{}_{}{}^{2}}{v_{c}^{(\phi )}{}_{}{}^{2}}\psi _z.$$
(28)
In this way, we arrive to an if and only if condition: If Eqs.(27,28) is satisfied, then the tangential velocity of circular stable equatorial orbits is constant. Furthermore, if the tangential velocity of circular stable equatorial orbits is constant, then the metric coefficients have to satisfy Eqs.(27,28). Notice then that if the function $`\psi `$ and $`\mu `$ are related by
$$e^\psi =(\frac{\mu }{\mu _0})^l.$$
(29)
with $`l=const,`$ we obtain that this a necessary and sufficient condition for the velocity $`v_{c}^{}{}_{}{}^{(\phi )}`$ to be the same for two orbits at different radii at the equatorial plane, provided that $`l=(v_{c}^{}{}_{}{}^{(\phi )})^2/\left(1+(v_{c}^{}{}_{}{}^{(\phi )})^2\right)`$.
We call your attention to the remarkable fact that the metric coefficient $`\gamma `$ does not play any role in this analysis, the motion analyzed is determined only by the other two metric coefficients, which now are related by this last equation, Eq.(28), thus leaving the problems in terms of only one metric coefficient. Actually this absence of $`\gamma `$ will be clear in the next section, where with the field equations we will see that it is determined in terms of the other metric coefficients and some components of the matter presented in the space-time, implying that it is not and independent function.
Thus, in order to have tangential velocities of equatorial objects circling the galaxy, and whose magnitude is radii independent, the form of the line element in the equatorial plane has to be
$$ds^2=\left(\frac{\mu }{\mu _0}\right)^{2l}dt^2+\left(\frac{\mu }{\mu _0}\right)^{2l}[e^{2\gamma }d\rho ^2+\mu ^2d\phi ^2].$$
(30)
Notice that this type of space time definitely can not be asymptotically flat. Neither it has the form of a space time related with a central black hole. What can be said is that this line element describes the region where the tangential velocity of the test particles is constant all over that region, and that it has to be joined in the interior and in the exterior regions with other types of space times if one wishes to have a central black hole, and that the influence of the middle region ends at some distance and thus has an asymptotically flat external region.
Taking into account the constraint between $`\psi `$ and $`\mu `$ given by Eq.(29), the energy, angular momentum, the rotational velocity , and the second derivative of the potential have the final expressions:
$`E`$ $`=`$ $`{\displaystyle \frac{\left(\frac{\mu }{\mu _0}\right)^l}{\sqrt{l_{}}}},`$ (31)
$`L`$ $`=`$ $`{\displaystyle \frac{\mu _0v_{c}^{}{}_{}{}^{(\phi )}\left(\frac{\mu }{\mu _0}\right)^{1/l_+}}{\sqrt{l_{}}}},`$ (32)
$`\mathrm{\Omega }`$ $`=`$ $`{\displaystyle \frac{v_{c}^{}{}_{}{}^{(\phi )}}{\mu _0}}\left({\displaystyle \frac{\mu }{\mu _0}}\right)^{l_{}/l_+},`$ (33)
$`{\displaystyle \frac{^2V(\rho )}{\rho ^2}}|_{extr}`$ $`=`$ $`2e^{2(\psi \gamma )}{\displaystyle \frac{l_+l_{}}{l_+^2}}\left({\displaystyle \frac{\mu _\rho }{\mu }}\right)^2.`$ (34)
being $`l_+=1+(v_{c}^{}{}_{}{}^{(\phi )})^2`$ and $`l_{}=1(v_{c}^{}{}_{}{}^{(\phi )})^2`$.
Notice that the second derivative of the potential at the extreme is always positive, thus the circular equatorial curves with constant tangential velocity are stable.
Before going to the field equations, we think it useful to present our derivations applied to the Schwarzschild case, thus testing the expressions while recovering the well known results. Starting with the line elements in spherical coordinates:
$$ds^2=\left(1\frac{2M}{r}\right)dt^2+\left(1\frac{2M}{r}\right)^1dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),$$
(35)
we perform the coordinate transformation
$$r=\sqrt{\rho ^2+z^2}+\frac{M^2}{4\sqrt{\rho ^2+z^2}}+M,\theta =\mathrm{tan}^1(\frac{\rho }{z}),$$
(36)
(the inverse transformation is $`\rho =R\mathrm{sin}\theta ,z=R\mathrm{cos}\theta `$, with $`R=\frac{1}{2}(rM+\sqrt{r^22Mr})`$), to obtain the line element in the Papapetrou form, Eq.(1), with $`e^{2\psi }=\left(\frac{\sqrt{\rho ^2+z^2}\frac{M}{2}}{\sqrt{\rho ^2+z^2}+\frac{M}{2}}\right)^2,e^{2\gamma }=\left(1\frac{M^2}{4(\rho ^2+z^2)}\right)^2,\mu =\rho \left(1\frac{M^2}{4(\rho ^2+z^2)}\right),\omega =0`$. The horizon in this coordinates is located at $`\sqrt{\rho ^2+z^2}=\frac{M}{2}`$.
Restricting the expressions to the equatorial plane, $`z=0`$, we have that $`\psi |_{z=0}=\mathrm{ln}\left(\left(\rho \frac{M}{2}\right)/\left(\rho +\frac{M}{2}\right)\right),`$ $`\mu |_{z=0}=\rho \left(1M^2/4\rho ^2\right)`$, thus from Eqs.(15), we obtain, for the energy and angular momentum:
$$E=\frac{(\rho \frac{M}{2})^2}{\left(\rho +\frac{M}{2}\right)\sqrt{\rho ^22M\rho +\frac{M^2}{4}}},L=\left(\rho +\frac{M}{2}\right)^2\sqrt{\frac{M}{\rho (\rho ^22M\rho +\frac{M^2}{4})}},$$
(37)
for the angular velocity test particles, from Eq.(18):
$$\mathrm{\Omega }=\sqrt{M}\rho ^{\frac{3}{2}}\left(1+\frac{M}{2\rho }\right)^3,$$
(38)
which gives us the Kepler law.. For the tangential velocity, from Eq.(26) we obtain
$$v^{(\phi )}=\sqrt{\frac{M}{\rho }}\left(1\frac{M}{2\rho }\right)^1,$$
(39)
with the known dependence as the inverse of the square root of the distance. Transforming back to spherical coordinates, it can be seen that our expressions agree with the usual ones, see for instance . For the second derivative of the potential, it is useful to write $`\rho `$ in terms of the horizon radius, as $`\rho =n\frac{M}{2}`$, with $`n`$ a number, and from Eq.(16) we obtain:
$$V_{\rho \rho }|_{extr}=\frac{4e^{2(\psi \gamma )}}{M^2n(n+1)^3}\frac{(n^210n+1)}{(n^24n+1)}.$$
(40)
This second derivative is positive down to $`\rho =(5+2\sqrt{6})\frac{M}{2}`$, which marks the last stable orbit, and corresponds to the known result of $`r=6M`$ in spherical coordinates.
In this way, we are confident on our expressions and can proceed to study the field equations.
A last remark about the geometric analysis. Recall that the observations are based on measurements of the red shift, not on the tangential velocity directly. If the space time is flat, the two quantities are proportional. But we are now working in curved space times, so we have to see if that proportionality is still valid. Following , we use the fact that the frequency of a photon is given by $`\nu =u^\alpha p_\alpha `$, with $`u^\alpha `$ the four velocity of the object and $`p^\alpha `$ the photon momentum, we have that the red shift, $`z`$, is given by
$$z=1\frac{\nu _{em}}{\nu _{rec}},$$
(41)
thus, for an object orbiting the galactic center in the equatorial plane at a distance $`\rho `$ from the center, with tangential velocity $`v^{(\phi )}`$ and emitting a photon with frequency $`\nu _0`$, and for an observer located at rest at infinity, that is far away from the emission, detecting the photon with a frequency $`\nu _{\mathrm{}}`$, it can be shown that the red shift is given by:
$$z=1\frac{(1+v^{(\phi )})}{\sqrt{1v_{}^{(\phi )}{}_{}{}^{2}}}\sqrt{\frac{g_{tt}(\rho )}{g_{tt}(\mathrm{})}}.$$
(42)
Now, for the observed velocities, we have that $`v1`$, i. e. they are much less than the speed of light, and we have to suppose that far away from the observed galaxy its gravitational influence ends, otherwise we could not detect the tangential velocity, we would then be moving along with the observed object! Thus we can take $`g_{tt}(\mathrm{})=1`$ and we have that at first order in the velocity, with $`g_{tt}(\rho )=e^\psi =1\psi \mathrm{}`$, we obtain:
$$z=(v^{(\phi )}+\psi +\mathrm{}).$$
(43)
But we have computed that in the case analyzed $`\psi =v_{}^{(\phi )}{}_{}{}^{2}/\left(1+v_{}^{(\phi )}{}_{}{}^{2}\right)\mathrm{ln}(\mu /\mu _0)2v_{}^{(\phi )}{}_{}{}^{2}(\mu 1)/(\mu +1)`$, thus $`\psi v_{}^{(\phi )}{}_{}{}^{2}`$, and we conclude that $`zv^{(\phi )}`$. In this way, we see that the radii independence of the value of the measured red shift can be related with the radii independence of the value of the tangential velocity, which is the fact that has been studied in this work.
## III Field equations
Now we are in a position where we can test any type of matter-energy to determine whether or not it produces a curvature in the space time such that the motion of the test particles can be circular stable and be such that the tangential velocity of those particles is constant for a large radial region in the equatorial plane.
We obtain the general form of the Einstein tensor for the axisymmetric static space time described by Eq. (1), with $`\omega =0`$ , and equate it to an arbitrary stress energy tensor. After some manipulations we conclude that the field equations are a set of two equations involving the metric coefficients $`\psi `$ , and $`\mu `$:
$`\mu (\psi _{\rho \rho }+\psi _{zz})+\mu _\rho \psi _\rho +\mu _z\psi _z=4\pi \mu [e^{2(\psi \gamma )}(e^{2\psi }T_{tt}+{\displaystyle \frac{e^{2\psi }}{\mu ^2}}T_{\phi \phi })+T_{\rho \rho }+T_{zz}],`$ (44)
$`\mu _{\rho \rho }+\mu _{zz}=8\pi \mu [T_{\rho \rho }+T_{zz}].`$ (45)
There are also two first order equations for the other metric coefficient $`\gamma `$:
$`\gamma _\rho \mu _\rho \gamma _z\mu _z\mu (\psi _{\rho }^{}{}_{}{}^{2}\psi _{z}^{}{}_{}{}^{2})+\mu _{zz}=8\pi \mu T_{\rho \rho },`$ (46)
$`\gamma _\rho \mu _z+\gamma _z\mu _\rho 2\mu \psi _\rho \psi _z\mu _{\rho z}=8\pi \mu T_{\rho z},`$ (47)
and finally, the field equations give us another equation for the second derivatives of $`\gamma `$ which thus is redundant, this equation is:
$$\gamma _{\rho \rho }+\gamma _{zz}+(\psi _\rho )^2+(\psi _z)^2=8\pi \frac{e^{2\gamma }}{\mu ^2}T_{\phi \phi }.$$
(48)
The analysis presented in the last section is exact and the relations between the metric coefficients and their first derivatives must be satisfied at the equatorial plane in order to describe the observed motion. Using the Einstein’s equations, we need the second derivatives of those metric coefficients. Thus, we have to make the approximation that the relations obtained among them, holds as well in a region close to the equatorial plane. Within this approximation, from Eqs.(27, 28,29), it can be obtained the following expression:
$$\mu (\psi _{\rho \rho }+\psi _{zz})+\mu _\rho \psi _\rho +\mu _z\psi _z=(v_{c}^{}{}_{}{}^{(\phi )})^2/\left(1+(v_{c}^{}{}_{}{}^{(\phi )})^2\right)(\mu _{\rho \rho }+\mu _{zz}).$$
(49)
Thus, with this last relation, from the Einstein’s equations Eqs. (44,45) we obtain a constraint among the stress energy tensor components which, within the approximation made on the validity of Eqs.(27, 28,29) out of the equatorial plane, has to be satisfied by any type of matter in order to have constant tangential velocities:
$$\left(\frac{1(v_{c}^{}{}_{}{}^{(\phi )})^2}{1+(v_{c}^{}{}_{}{}^{(\phi )})^2}\right)(T_{\rho \rho }+T_{zz})=e^{2(\psi \gamma )}\left(e^{2\psi }T_{tt}+\frac{e^{2\psi }}{\mu ^2}T_{\phi \phi }\right),$$
(50)
We finish our analysis by testing several types of matter, described by their respective stress energy tensor, to see whether or not they are able to deform the geometry of the space time in such a way that the tangential velocity of the equatorial rotational objects be constant, that is, that they satisfy Eq.(50).
### A Vacuum Fields
We start with the vacuum solutions, with $`T_{\mu \nu }=0`$. In this case Eq.(50) is trivially satisfied, thus we proceed to analyze the Einstein equations directly. From Eq.(45), the easiest solution implies $`\mu =\rho `$. Eq.(44), is a Laplace equation for $`\psi `$ and ‘imposing the condition that at the equatorial plane the flat curve condition, Eq.(28), be satisfied, as well as the one of symmetry with respect to the galactic plane, we obtain that $`\psi =l\mathrm{ln}\rho `$ and the other Einstein equation then imply $`\gamma =l^2\mathrm{ln}\rho `$, with $`l=(v_{c}^{}{}_{}{}^{(\phi )})^2/\left(1+(v_{c}^{}{}_{}{}^{(\phi )})^2\right)`$. In this way, we obtained an exact vacuum solution for the Einstein equations, which produces that the test particles circling at the equatorial plane behave in agreement with the observations:
$$ds^2=\rho ^{2l}dt^2+\rho ^{2l}(\rho ^{2l^2}(d\rho ^2+dz^2)+\rho ^2d\phi ^2).$$
(51)
The central object is string-like. Observations show that cosmic strings object are very unlikely to exist, nevertheless, this is an example of objects which could produce the observed motion of test particles, and in which the density does not go as $`r^2`$, because it is vacuum. Thus, such a behavior on the density is a sufficient but not a necessary condition for the flatness of the rotational curves.
### B Perfect Fluid
For the perfect fluid, $`T_{\mu \nu }=(d+p)u_\mu u_\nu +g_{\mu \nu }p`$, with $`d`$ the density of the fluid and $`p`$ its pressure. In this case we are thinking on a “dark fluid”, which could be composed of planetoids or WIMPS or MACHOS, which are not seen but it is thought that they could be there affecting the geometry in the way needed in order to have the observed behavior in the tangential velocities of the luminous matter. Taking this dark fluid as static, the four velocity is given by $`u^\alpha =(u^0,0,0,0)`$ with, for the line element given by Eq.(1) with $`\omega =0`$, $`u^0=Ee^{2\psi }`$, and $`L=0`$. Thus $`u_0=E`$, and from $`u_\alpha u^\alpha =1`$, we conclude that $`E=e^\psi `$. The stress energy tensor for the dark fluid then has the form:
$`T_{tt}=dE^2=de^{2\psi }`$ (52)
$`T_{\rho \rho }=T_{zz}=e^{2(\psi \gamma )}p,`$ (53)
$`T_{\phi \phi }=\mu ^2e^{2\psi }p.`$ (54)
Substituting in Eq.(50), we obtain that in the equatorial plane, in order to satisfy the observed behavior on the tangential velocities, the “dark fluid” has to satisfy:
$$2\left(\frac{1(v_{c}^{}{}_{}{}^{(\phi )})^2}{1+(v_{c}^{}{}_{}{}^{(\phi )})^2}\right)p=(d+p),$$
(55)
Thus, we obtain an equation of state for the “dark fluid” particles at the equatorial plane:
$$p=\frac{1+(v_{c}^{}{}_{}{}^{(\phi )})^2}{3(v_{c}^{}{}_{}{}^{(\phi )})^2}d,$$
(56)
which, compared to the equation of state for a perfect fluid, $`p=\omega d`$, implies that $`1<\omega <\frac{1}{3}`$, for $`v_{c}^{}{}_{}{}^{(\phi )}`$ between the speed of light and zero. This result is quite remarkable. It coincides with the type of equation of state derived within the Quintessence model at the cosmological level, and now we obtain similar results at the galactic level. This sort of matter has been called exotic matter and studied in several contexts . Our result points to the fact that the Dark Matter actually could be exotic. We want to stress that due to the approximation taken for the behavior of the metric coefficients off the galactic plane, we are not excluding the possibility that the dark fluid be composed of baryonic usual matter, actually from the Newtonian approach, we know that regular matter can produce the observed motion. This fact is not reproduced in the present analysis, due to our approximation, what we certainly can conclude is that exotic type of matter also can produce such ‘observed motion.
In order to recover the Newtonian case, where we know that the dust type fluid does work as the dark matter, we have to analyze the spherical case, which is introduced in Appendix V. For the stress energy tensor we again take the static perfect fluid, thus the four velocity of the test particle reads $`u^\mu =(u^0,0,0,0)`$, with $`u^0=\dot{t}`$. For the spherically symmetric metric, Eq.(73), we know that $`\dot{t}`$ is associated to a conserved quantity, the energy, and is given by: $`\dot{t}=\frac{E}{B(r)}`$, and from the normalization of the four velocity, $`u_\alpha u^\alpha =1`$, we get that $`E=\sqrt{B(r)}`$. Thus, considering those spacetimes for which the tangential velocity ($`v_c`$) of test particles in circular orbits is radii independent, we use Eq.(81), obtaining that $`u^0=\frac{1}{\sqrt{B_0}}r^{(v_c)^2}`$, thus, $`u_0=\sqrt{B_0}r^{(v_c)^2}`$. In this way, we get that the non zero components of the static spherically symmetric perfect fluid are:
$`T_{tt}`$ $`=`$ $`dB_0r^{2(v_c)^2},`$ (57)
$`T_{rr}`$ $`=`$ $`pA(r),`$ (58)
$`T_{\theta \theta }`$ $`=`$ $`pr^2.`$ (59)
Substituting these expressions in the Einstein equations Eq.(84), it turns out that the equations can be completely solved, yielding:
$`A(r)`$ $`=`$ $`{\displaystyle \frac{b}{2(1A_0r^{\frac{b}{a}})}},`$ (60)
$`d`$ $`=`$ $`{\displaystyle \frac{(2+a)A(r)ab}{8\pi ar^2A(r)}},`$ (61)
$`p`$ $`=`$ $`{\displaystyle \frac{1+2(v_c)^2A(r)}{8\pi r^2A(r)}},`$ (62)
with $`b=2(1+2(v_c)^2v_{c}^{}{}_{}{}^{4}),a=1+(v_c)^2`$, and $`A_0`$ an integration constant.
Taking the particular case for the integration constant $`A_0=0`$, we get
$`A(r)`$ $`=`$ $`{\displaystyle \frac{b}{2}},`$ (63)
$`d`$ $`=`$ $`{\displaystyle \frac{(v_c)^2(1\frac{(v_c)^2}{2})}{2\pi br^2}},`$ (64)
$`p`$ $`=`$ $`{\displaystyle \frac{(v_c)^4}{4\pi br^2}}.`$ (65)
In this way, we obtain the particular solution where the particles move in the observed way, in a space-time with a deficit angle, $`g_{rr}=constant`$, and in which the density goes as $`d(v_c)^2/r^2`$, and the pressure goes as $`p(v_c)^4/r^2`$, thus it is a dust like solution. Furthermore, for an equation of state $`p=\omega d`$, this dust like solution implies $`\omega =(v_c)^2/\left(2(1(v_c)^2/2)\right)`$, which is between $`0`$ and $`1`$, for the tangential velocity between $`0`$ and $`1`$, and thus is a perfectly well known fluid type. In this way we recover the dust hypotheses within our approach, and we clearly see that the‘ fact that we did not recover this case within the axi-symmetric analysis, was due to our approximation outside the equatorial plane.
For the general case when the integration constant $`A_0`$ is non zero, taking an equation of state as before $`p=\omega d`$, we obtain that the $`\omega `$ is a function of $`r`$ and again, as in the axial case, it is negative, so we are dealing with the exotic type of perfect fluid obtained in the axi-symmetric case.
This is a good moment to discuss the question on the need to use GR, even though the gravitational field is weak. In the Newtonian description it is well known that the space-time can be described as
$`ds^2=(1+2\mathrm{\Phi })dt^2+(12\mathrm{\Phi })dr^2+r^2d\mathrm{\Omega }^2`$
with $`\mathrm{\Phi }=\mathrm{\Phi }(r)`$ the Newtonian gravitational potential. For the spherically symmetric case we obtain that
$`g_{tt}=r^{2\left(v_c\right)^2}=e^{2\left(v_c\right)^2\mathrm{ln}(r)}=1+2\left(v_c\right)^2\mathrm{ln}(r)+\mathrm{}`$
from here we determine that $`\mathrm{\Phi }=\left(v_c\right)^2\mathrm{ln}(r)`$. In this Newtonian approximation the complete set of 10 Einstein equations reduces to one equation, the usual Poisson equation $`^2\mathrm{\Phi }=4\pi Gd`$. Collecting this last result, it is obtain that the density, $`d,`$ goes as $`d\left(v_c\right)^2/r^2`$, which is the expression for the dark matter density known from the astronomers’ work.
However, notice that in this approximation it can not be said anything else about the matter producing the observed motion. The Newtonian approximation fixes the matter to be dust-perfect-fluid-like type. This is the usual way of reasoning: it is suppose a priori that the dark matter is a completely Newtonian dust and at the end of the day one arrives to a consistent description of the dark matter determining only the shape of the Newtonian gravitational potential.
In this work we are proceeding in a different way. We are using Einstein equations backwards; we do not make any assumptions in the type of matter nor do approximations. From the observations on the motion of the test particles, we determine the geometry and then, by means of the Einstein’s equations we obtain constrains of the type of matter. We have shown that we do recover the Newtonian result, but also it is clear that this is a very particular case for a very specific type of matter. In the general reasoning, we do not fix neither the type of matter nor the equation of state, we let the equations themselves to do that obtaining more general results.
To end this argumentation we recall the reader that there is three conditions that have to be fulfill in order to reach the Newtonian limit; 1) the speeds of the study particles most be most less than the speed of light, 2) the gravitational field must be weak, and 3) the pressures associated with the matter study must to be must smaller than the corresponding density. It is this last condition which is taking a priori in the usual analysis and it is not satisfy in the general case as we have shown, thus justify the need of GR in order to be able to consider any type of matter.
As a last remark about these results on the type of perfect fluid is that the Big-Bang nucleosynthesis imposes very strong constraints to the percentage of the baryonic matter to the total content of the Universe. If the dark matter would be a dark fluid of baryonic matter, such percentage would be quite above the value settled by those constraints. Thus, even if the baryonic dark fluid can not be discarded by dynamical methods, the cosmological constraints make it unlikely, a fact which might strength the case for exotic type of dark fluids or for other type of dark matter.
### C Cosmological Constant
For a cosmological constant, $`\mathrm{\Lambda }`$, $`T_{\mu \nu }=\mathrm{\Lambda }g_{\mu \nu }`$. We thus have $`T_{tt}=e^{2\psi }\mathrm{\Lambda },T_{\rho \rho }=T_{zz}=e^{2(\psi \gamma )}\mathrm{\Lambda },T_{\phi \phi }=e^{2\psi }\mu ^2\mathrm{\Lambda }`$, thus from Eq.(50), we obtain
$$\left(\frac{1(v_{c}^{}{}_{}{}^{(\phi )})^2}{1+(v_{c}^{}{}_{}{}^{(\phi )})^2}\right)\mathrm{\Lambda }=0,$$
(66)
In this way we see that, within our approximation, a non zero cosmological constant can not explain the observed behavior because implies that the observed tangential velocity had to be equal to 1, i. e., they should be moving at the speed of light. Something similar occurs with the scalar field.
### D Scalar Field
For scalar field $`\varphi `$ with potential, $`T_{\mu \nu }=\varphi _{,\mu }\varphi _{,\nu }\frac{1}{2}g_{\mu \nu }\varphi ^\alpha \varphi _\alpha +g_{\mu \nu }V(\varphi )`$. We have that, due to the symmetry of our space time $`\varphi =\varphi (\rho ,z)`$ the $`T_{\mu \nu }`$ components are:
$`T_{tt}`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{2(2\psi \gamma )}(\varphi _{\rho }^{}{}_{}{}^{2}+\varphi _{z}^{}{}_{}{}^{2})e^{2\psi }V(\varphi ),`$ (67)
$`T_{\rho \rho }`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\varphi _{\rho }^{}{}_{}{}^{2}\varphi _{z}^{}{}_{}{}^{2})+e^{2(\psi \gamma )}V(\varphi ),`$ (68)
$`T_{zz}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\varphi _{\rho }^{}{}_{}{}^{2}\varphi _{z}^{}{}_{}{}^{2})+e^{2(\psi \gamma )}V(\varphi ),`$ (69)
$`T_{\rho z}`$ $`=`$ $`\varphi _\rho \varphi _z,`$ (70)
$`T_{\phi \phi }`$ $`=`$ $`{\displaystyle \frac{1}{2}}e^{2\gamma }\mu ^2(\varphi _{\rho }^{}{}_{}{}^{2}+\varphi _{z}^{}{}_{}{}^{2})+e^{2\psi }\mu ^2V(\varphi ).`$ (71)
Inserting these components in Eq.(50), we obtain that, as in the cosmological constant case
$$\left(\frac{1(v_{c}^{}{}_{}{}^{(\phi )})^2}{1+(v_{c}^{}{}_{}{}^{(\phi )})^2}\right)V(\varphi )=0.$$
(72)
Again, either the particles move at the speed of light, or the scalar field potential is zero at the equatorial plane. When the scalar field potential is zero, that is, we have a massless scalar field, Eq.(50) is satisfied, and we have to go back to the Einstein’s equations as in the vacuum case. Again from Eq.(45), we take the simplest solution $`\mu =\rho `$, and, as in the vacuum case, we obtain $`\psi =l\mathrm{ln}\rho `$. The last metric coefficient $`\gamma `$, can be solved in terms of the scalar field. Thus, to have a complete solution, it is only left to solve the Klein Gordon equation for the scalar field: $`D^2\varphi +\frac{1}{\mu }D\varphi D\mu =0`$, which turns out to be the same equation for the metric coefficient $`\psi `$. However, in this case we do not have boundary conditions well defined: The space time is not asymptotically flat; it is not known the form of the space time near and at the origin, we are only analyzing the region where the curves are flat; there are no conditions for the scalar field at the equatorial plane. What can be concluded at this stage, is that the scalar field does remain being a candidate for the dark matter, and thus to contribute with about 25% of the matter of the Universe.
## IV Conclusions
We have found that in a static, axisymmetric space-time, a sufficient and necessary condition in order to have a flat profile for the rotational curves in a plane of that space-time is that its metric tensor must have the form given by Eq. (30). This form of the metric must be the one required for galaxies in the region where the rotational curves profile of stars is flat. It is important to stress the fact that in the derivation of this expression, only the geometry of the space-time was involved, thus, it is independent of the type of matter which generates such a geometry; that is, whatever the matter might be, Eq. (30) must be the form of the line element at the galactic plane for a static axisymmetric space time which can be expressed in the Papapetrou form, Eq. (1) with $`\omega =0`$, thus, this result is not only the general relativistic analog of the Newtonian result for the gravitational Newtonian potential, $`\varphi (r)1/r`$, but it can be used for any type of matter, including those which do not have a clear Newtonian expression, such as the scalar field.
With this idea in mind, we proceeded further using the Einstein’s equations, which essentially describe the inter dependence of matter-energy and geometry. We had to accept some loose of the generality of our results, in making the assumption that the definite relation which we obtained for the metric coefficients at the galactic plane, Eqs.(29), are also valid in a close by region off the plane, and thus determine a relation for the second derivatives of the metric coefficients, Eq.(49). Within this approximation, we were able to obtain a constraint equation among the components of a general stress energy tensor, $`T^{\alpha \beta },`$ Eq.(50). We tested this expression in four types of stress energy tensors which included the traditional types of matter which have been used as candidates for the dark matter in the galactic halos, such as the perfect fluid or the cosmological constant. We obtained that for the vacuum case, a cosmic string type of matter does generates the observed motion of test particles. Even though the cosmic strings are unlikely objects to be the Universe, it was a clear example for the fact that the Newtonian behavior of the density, $`\rho (r)1/r^2`$, is not a necessary condition for describing the observed motion. We analyzed also the static perfect fluid, and it is interesting that in difference with the Newtonian description, within the general relativistic formulation, we are able to obtain conditions for the equation of state of the dark fluid. Even though we were not able to reproduce the well known Newtonian result for dust-like fluid, due to our approximation for the second derivatives of the metric coefficients, we did showed that a dark fluid with an exotic type of matter, is a candidate for being the dark matter. Furthermore, this results represents the Quintessence type of matter at a galactic level. We performed the study in the spherically symmetric static space time in appendix V, and when applied to the perfect fluid static case we did recover the Newtonian result, namely that the “dark fluid” could be a well behaved dust like fluid. We further analyzed the cosmological constant case, which within our approximation implied that it can not be, and the massive scalar field case, which again, within our approximation, turned out that it has to be massless, and the massless scalar field also remains being a candidate for the dark matter.
As we have mentioned above, our results are useful for describing the region where is observed the flat behavior of the test particles rotating around the galactic center, it is clearly needed to proceed further in order to be able to describe the motion in the complete region, from the center to the exterior. Some preliminary results indicate that a combination of perfect fluid with baryonic matter and some of the matter analyzed here, could be in the right direction, . Also, the approximation we made for the second derivatives of the metric coefficients, has to be further analyzed. It would be of great help, in order to be able to apply the present description in objets out from the glactic plane, to have a sample of the profile of velocities of such objects. In the simpler spherically symmetric case, there is no need to make such approximation, and more definite results can be obtained, though with more restrictions on the geometry. Besides, there are reasons to belive that the dark matter halo is spherical , thus it might be a good approximation the analysis made within this symmetry. We studied the scalar field in this symmetry in .
Finally, as the best way to study dark matter is through its effects on the dynamics of the visible objetcs, further studies along the lines presetned in this work can be performed using gravitational lensing or jets.
Overall, we consider that the analysis presented in this work, is on the right track in order to determine which is the type of matter which constitutes the 90% of the matter in the galaxies.
## V Acknowledgments
This work was partially supported by a grant CONACyT-DFG, by CONACyT, México, under grants 94890 (F.S.G.), and DGAPA-UNAM IN121298 (D.N.) We want to thank the relativity group in Jena for its kind hospitality and partial support.
In this appendix we present an analogous derivation for the conditions which the constancy of the tangential velocity of circular orbits impose on the metric coefficients for the spherical static case. Furthermore, we also present the Einstein equations in this case for a general stress energy tensor.
It is interesting that, due to the symmetries, in this case we do not have to restrict the analysis to equatorial orbits, and that the above mentioned conditions give a closed form for the metric coefficient $`g_{tt}`$.
We begin with the line element
$$ds^2=B(r)dt^2+A(r)dr^2+r^2(d\theta ^2+\mathrm{sin}^2\theta d\phi ^2)$$
(73)
The Lagrangian for a test particle reads
$$2=B(r)\dot{t}^2+A(r)\dot{r}^2+r^2(\dot{\theta }^2+\mathrm{sin}^2\theta \dot{\phi }^2).$$
(74)
We infer the conserved quantities, the energy $`E=B(r)\dot{t}`$, the $`\phi `$-momentum $`L_\phi =r^2\mathrm{sin}^2\theta \dot{\phi }`$, and the total angular momentum, $`L^2=L_{\theta }^{}{}_{}{}^{2}+(\frac{L_\phi }{\mathrm{sin}\theta })^2`$, with $`L_\theta =r^2\dot{\theta }`$. The radial motion equation can thus be written as:
$$\dot{r}^2+V(r)=0,$$
(75)
with the potential $`V(r)`$ given by
$$V(r)=\frac{1}{A(r)}(\frac{E^2}{B(r)}\frac{L^2}{r^2}1).$$
(76)
Notice that, due to the spherical symmetry, we do not need to restrict the study to equatorial orbits, this last radial motion is valid for any angle $`\theta `$. For circular stable orbits, we again have the conditions, $`\dot{r}=0,V_r=0`$, and $`V_{rr}>0`$, which imply the following expressions for the energy and total momentum of the particles in such orbits:
$`E^2`$ $`=`$ $`{\displaystyle \frac{2B(r)^2}{2B(r)rB(r)_r}},`$ (77)
$`L^2`$ $`=`$ $`{\displaystyle \frac{r^3B(r)_r}{2B(r)rB(r)_r}},`$ (78)
and for the second derivative of the potential evaluated at the extrema
$$V(r)_{rr}|_{extr}=2\frac{\frac{rB(r)_{rr}}{B}+\frac{B(r)_r}{B}(3\frac{2rB(r)_r}{B})}{rA(r)(2\frac{rB(r)_r}{B})}.$$
(79)
On the other hand, in a similar manner as it was presented in the text, we obtain that the tangential velocity, $`(v_c)^2=\frac{r^2}{B(r)}((\frac{d\theta }{dt})^2+\mathrm{sin}^2\theta (\frac{d\phi }{dt})^2)`$, for particles in stable circular orbits is given by:
$$(v_c)^2=\frac{rB(r)_r}{2B(r)}.$$
(80)
Thus, imposing the observed condition that this tangential velocity velocity is constant for several radii, this last equation can be integrated for the metric coefficient $`g_{tt}`$:
$$B(r)=B_0r^{2(v_c)^2},$$
(81)
with $`B_0`$ an integration constant.
In this way, we again arrive to a theorem, stating that: For a static spherically symmetric spacetime, $`v_c`$, the tangential velocity of particles moving in circular stable orbits is radii independent if and only if the $`g_{tt}`$ metric coefficient has the form $`g_{tt}=B_0r^{2(v_c)^2}`$.
Notice that in this case, one of the metric coefficients was completely integrated and the other one, $`A(r)`$ remains arbitrary. Also, as mentioned, the analysis made no suppositions on the plane of motion, so the result is valid for any circular stable trajectory.
Finally, we can construct the Einstein tensor and arrive to the following Einstein equations which give us information about the type of matter curving the spacetime in such a way that the motion corresponds to the observed one:
$`{\displaystyle \frac{B_0r^{2(v_c)^2}}{A(r)^2}}(rA(r)^{}+A(r)(A(r)1))`$ $`=`$ $`8\pi T_{tt},`$ (82)
$`{\displaystyle \frac{2(v_c)^2+1A(r)}{r^2}}`$ $`=`$ $`8\pi T_{rr},`$ (83)
$`{\displaystyle \frac{r((v_c)^2+1)A(r)^{}2v_{c}^{}{}_{}{}^{4}A(r)}{2A(r)^2}}`$ $`=`$ $`8\pi T_{\theta \theta },`$ (84)
where stands for derivative with respect to $`r`$. This study was applied for the scalar field in .
In this appendix we present the generalization of the derivation of the constraint equation among the metric coefficients, Eq.(28), for the stationary case, where $`\omega 0`$, described by the line element 1.
¿From Eq.(6), we express $`\dot{t}`$, and $`\dot{\phi }`$. in terms of $`E,L`$, and the metric coefficients as
$`\dot{t}`$ $`=`$ $`{\displaystyle \frac{e^{2\psi }}{\mu ^2}}[(\mu ^2e^{4\psi }\omega ^2)E\omega L],`$ (85)
$`\dot{\phi }`$ $`=`$ $`{\displaystyle \frac{e^{2\psi }}{\mu ^2}}(\omega E+L).`$ (86)
Using these equation in the constraints ones, Eqs.(9) we arrive at:
$`\mu ^2e^{2\psi }(1e^{2\psi }E^2)+(\omega E+L)^2`$ $`=`$ $`0,`$ (87)
$`(e^{2\psi })_\rho E^2+\left({\displaystyle \frac{e^{2\psi }}{\mu ^2}}\right)_\rho (\omega E+L)^2+{\displaystyle \frac{2e^{2\psi }}{\mu ^2}}(\omega E+L)\omega _\rho E`$ $`=`$ $`0,`$ (88)
Solving for $`E`$ and $`L`$, we obtain:
$`E`$ $`=`$ $`e^\psi \sqrt{{\displaystyle \frac{𝒜}{}}},`$ (89)
$`L`$ $`=`$ $`{\displaystyle \frac{\mu e^\psi }{\sqrt{}}}(\sqrt{𝒜}{\displaystyle \frac{\omega e^{2\psi }}{\mu }}\sqrt{𝒜}),`$ (90)
where
$`𝒜`$ $`=`$ $`2e^{4\psi }(\mu _\rho \mu \psi _\rho )(\mu _\rho 2\mu \psi _\rho )(\omega _\rho )^2+`$ (93)
$`\pm \omega _\rho \sqrt{(\omega _\rho )^24\mu \psi _\rho e^{4\psi }(\mu _\rho \mu \psi _\rho )},`$
$``$ $`=`$ $`2e^{4\psi }(\mu _\rho 2\mu \psi _\rho )^22(\omega _\rho )^2.`$ (94)
For the second derivative of the potential $`V(\rho )`$ evaluated at the extreme, we obtain:
$`V_{\rho \rho }|_{extr}`$ $`=`$ $`{\displaystyle \frac{2e^{2(\psi \gamma )}}{}}[(2\psi _{\rho \rho }{\displaystyle \frac{\mu _{\rho \rho }}{\mu }}+3\left({\displaystyle \frac{\mu _\rho }{\mu }}\right)^24{\displaystyle \frac{\mu _\rho }{\mu }}\psi _\rho +{\displaystyle \frac{e^{4\psi }}{\mu ^2}}(\omega _\rho )^2)𝒜`$ (97)
$`(\psi _{\rho \rho }{\displaystyle \frac{\mu _{\rho \rho }}{\mu }}+2(\psi _\rho )^2+3\left({\displaystyle \frac{\mu _\rho }{\mu }}\right)^24{\displaystyle \frac{\mu _\rho }{\mu }}\psi _\rho )+`$
$`(4(\psi _\rho {\displaystyle \frac{\mu _\rho }{\mu }})\omega _\rho +\omega _{\rho \rho }){\displaystyle \frac{\sqrt{𝒜(𝒜)}}{\mu }}].`$
On the other hand, using Eqs.(86), and (90), in the expression for the angular velocity, Eq.(17) we obtain that:
$$\mathrm{\Omega }=\frac{e^{2\psi }}{\mu }\frac{\sqrt{𝒜}}{\sqrt{𝒜}\frac{\omega e^{2\psi }}{\mu ^2}\sqrt{𝒜}},$$
(98)
where $`𝒜`$ and $``$ are given by Eqs.(94).
As in the static case, Following Chandrasekhar , we rewrite the line element given in Eq.(1) as:
$`ds^2`$ $`=`$ $`{\displaystyle \frac{\mu ^2e^{2\psi }}{\mu ^2e^{4\psi }\omega ^2}}dt^2+e^{2\psi }(\mu ^2e^{4\psi }\omega ^2)\left(d\phi {\displaystyle \frac{\omega }{\mu ^2e^{4\psi }\omega ^2}}dt\right)^2+`$ (100)
$`e^{2(\psi \gamma )}(d\rho ^2+dz^2),`$
thus, in terms of the proper time, $`d\tau ^2=ds^2`$, we have that
$`d\tau ^2`$ $`=`$ $`{\displaystyle \frac{\mu ^2e^{2\psi }}{\mu ^2e^{4\psi }\omega ^2}}dt^2[1{\displaystyle \frac{e^{4\psi }(\mu ^2e^{4\psi }\omega ^2)^2}{\mu ^2}}({\displaystyle \frac{d\phi }{dt}}{\displaystyle \frac{\omega }{\mu ^2e^{4\psi }\omega ^2}})^2+`$ (102)
$`{\displaystyle \frac{e^{2\gamma }(\mu ^2e^{4\psi }\omega ^2)}{\mu ^2}}(\left({\displaystyle \frac{d\rho }{dt}}\right)^2+\left({\displaystyle \frac{dz}{dt}}\right)^2)],`$
from which we can write that
$$1=\frac{\mu ^2e^{2\psi }}{\mu ^2e^{4\psi }\omega ^2}u_{}^{0}{}_{}{}^{2}[1v^2],$$
(103)
where $`u^0=\frac{dt}{d\tau }`$ is the usual time component of the four velocity, and a definition of the spatial velocity, $`v^2`$, again comes out naturally in this way.
$`v^2`$ $`=`$ $`{\displaystyle \frac{e^{4\psi }(\mu ^2e^{4\psi }\omega ^2)^2}{\mu ^2}}\left({\displaystyle \frac{d\phi }{dt}}{\displaystyle \frac{\omega }{\mu ^2e^{4\psi }\omega ^2}}\right)^2+`$ (105)
$`+{\displaystyle \frac{e^{2\gamma }(\mu ^2e^{4\psi }\omega ^2)}{\mu ^2}}\left(\left({\displaystyle \frac{d\rho }{dt}}\right)^2+\left({\displaystyle \frac{dz}{dt}}\right)^2\right),`$
which is the 3-velocity of a particle measured with respect to an orthonormal reference system, it has components:
$$v^2=v_{}^{(\phi )}{}_{}{}^{2}+v_{}^{(\rho )}{}_{}{}^{2}+v_{}^{(z)}{}_{}{}^{2}.$$
(106)
For the $`\phi `$component of the spatial velocity we obtain:
$$v^{(\phi )}=\frac{e^{2\psi }}{\mu }[(\mu ^2e^{4\psi }\omega ^2)\mathrm{\Omega }\omega ],$$
(107)
and substituting $`\mathrm{\Omega }`$ from Eq.(98), we finally obtain an expression for the tangential velocity of a test particle in stable circular motion:
$$v^{(\phi )}=\frac{\mu e^{2\psi }\sqrt{𝒜}\omega \sqrt{𝒜}}{\mu e^{2\psi }\sqrt{𝒜}\omega \sqrt{𝒜}},$$
(108)
where $`𝒜`$ and $``$ are given by Eqs.(94).
Imposing the condition of constancy for all radii, that is $`v_\rho ^{(\phi )}=0`$, thus $`v^{(\phi )}=v_c^{(\phi )}`$, with $`v_c^{(\phi )}`$ a constant, representing the value of the velocity, from Eq. (108), we finally have that:
$$=(1v_{c}^{(\phi )}{}_{}{}^{2})F^2𝒜,$$
(109)
where $`F=\left(\mu ^2e^{4\psi }\omega ^2\right)/\left(\mu e^{2\psi }+v_c^{(\phi )}\omega \right)^2`$. This last expression, represents a constraint among three of the metric coefficients, and we can express one of them, say $`\omega `$, in terms of the other two: $`\psi `$, and $`\mu `$. In this way, we again arrive to an iff condition, namely: The tangential velocity of a test particle moving in a circular equatorial motion in an axisymmetric stationary background, has a radii independent magnitude iff the metric coefficients satisfy the constraint equation (109).
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# Tame coverings of arithmetic schemes
## 1 Tame coverings
The concept of tame ramification stems from number theory: A finite extension of number fields $`L|K`$ is called tamely ramified at a prime $`𝔓`$ of $`L`$ if the associated extension of completions $`L_𝔓|K_𝔓`$ is a tamely ramified extension of local fields. The latter means that the ramification index is prime to the characteristic of the residue field. It is a classical result that composites and towers of tamely ramified extensions are again tamely ramified. This concept generalizes to arbitrary discrete valuation rings by requiring that the associated residue field extensions are separable.
As is well known, the concept of unramified extensions has found its generalization to arbitrary schemes by the notion of étale coverings. Tame ramification along a normal crossing divisor of a regular scheme has been studied in \[SGA1\], \[G-M\]. Let us first recall the definition given there. We assume throughout this paper that all schemes under consideration are noetherian.
Let $`X`$ be a regular and connected scheme and let $`D=D_i`$ be a divisor on $`X`$. We say that $`D`$ has normal crossings, if étale locally around every point $`xsupp(D)X`$, we have
$$D_i=\underset{j}{}div(s_{ij}),$$
where $`(s_{ij})_{i,j}𝒪_{X,x}`$ is part of a regular system of parameters of the regular local ring $`𝒪_{X,x}`$.
Let $`X`$ be regular and connected, $`DX`$ a divisor with normal crossings and $`U=XD`$. Let $`P_i`$, $`i=1,\mathrm{},n`$, be the generic points of the irreducible components $`D_1,\mathrm{},D_n`$ of $`D`$. Then the local rings $`𝒪_{X,P_i}`$ are discrete valuation rings, inducing discrete valuations $`v_1,\mathrm{},v_n`$ on the function field of $`X`$. Let $`U^{}U`$ be a finite étale morphism and assume for simplicity that $`U^{}`$ is connected. Let $`X^{}`$ be the normalization of $`X`$ in the function field $`K^{}`$ of $`U^{}`$.
###### Definition 1.1
(\[G-M\], 2.2.2.) The finite étale covering $`U^{}U`$ is called tamely ramified (or tame, for short) along $`D`$ if the extension of function fields $`K^{}|K`$ is tamely ramified at the discrete valuations associated to $`D_1,\mathrm{},D_n`$.
Étale locally, tame coverings are of a very simple structure by the following theorem which is known under the name Generalized Abhyankar’s Lemma.
###### Theorem 1.2
(\[SGA1\], Exp. XIII, 5.3.0) Let $`X`$ be a strictly henselian local regular scheme of residue characteristic $`p>0`$, $`D=_{i=1}^rdiv(f_i)`$ a divisor with normal crossings on $`X`$ and $`U=XD`$. Then every connected finite étale covering of $`U`$ which is tamely ramified along $`D`$ is a quotient of a (tamely ramified) covering of the form
$$U^{}=U[T_1,\mathrm{},T_r]/(T_1^{n_1}f_1,\mathrm{},T_r^{n_r}f_r),$$
where the $`n_i`$ are natural numbers prime to $`p`$.
Tame coverings as defined above satisfy the axioms of a Galois theory, so (omitting the base point) there exists a profinite group $`\pi _1^t(X,D)`$, which is a quotient of $`\pi _1(U)`$ and which classifies étale coverings of $`U`$ which are tame along $`D`$.
If $`X`$ is strictly local and $`D=_{i=1}^rdiv(f_i)`$ a normal crossing divisor, then theorem 1.2 yields a natural isomorphism
$$\pi _1^t(X,D)\left(\underset{\mathrm{}p}{}_{\mathrm{}}(1)\right)^r\stackrel{notation}{=}\left(\widehat{}^{(p^{})}(1)\right)^r.$$
If, more generally, $`DX`$ is an arbitrary divisor on $`X`$, then (as already remarked in \[G-M\], 2.2.3.4, without further elaboration) the above definition of tame ramification is not the ‘correct’ one. For example, it is not stable under base change.
###### Example 1.3
Let $`X=Spec([T])`$ be the affine line over $`Spec()`$ and consider the divisor
$$D=div(T+4)+div(T4),$$
which is not a normal crossing divisor. Let $`K=(T)`$ be the function field of $`X`$ and $`U=XD`$. Put $`f=(T+4)(T4)=T^216`$, $`L=K(\sqrt{f})`$ and consider the normalization $`X_L`$ of $`X`$ in $`L`$. The ramification locus of $`X_LX`$ is either $`D`$ or $`DX_2`$, where $`X_2`$ is the unique vertical divisor on $`X`$ over characteristic $`2`$. Let us show that $`X_LX`$ is unramified at the generic point of $`X_2`$. This is equivalent to the statement that $`L|K`$ is unramified at the unique discrete valuation $`v_2`$ of $`K`$ which corresponds to the prime ideal $`2[T][T]`$. Therefore it suffices to show that $`f`$ is a square in the completion $`K_2`$ of $`K`$ with respect to $`v_2`$. Consider the polynomial $`F(X)=X^2f=X^2T^2+16`$. We have $`F(T)0mod16`$ and the derivative $`F^{}(T)=2T`$ has the exact $`2`$-valuation $`1`$. By the usual approximation process (cf. \[Se\] 2.2. th.1), we see that $`f`$ has a square root in $`K_2`$. Hence the ramification locus of $`X_LX`$ is exactly $`D`$, and since $`D`$ is the sum of horizontal prime divisors, the morphism $`U_LU`$ is tame along $`D`$ (in the naive sense).
Now consider the closed subscheme $`YX`$ given by the equation $`T=0`$, so $`YSpec()`$. Then $`D_Y=DY`$ is the point on $`Y`$ which corresponds to the prime number $`2`$. Let $`V=UY=YD_Y`$. The base change $`V^{}=U_L\times _UVV`$ is the normalization of $`VSpec([\frac{1}{2}])`$ in $`(\sqrt{1})`$. But $`2`$ is wildly ramified in $`(\sqrt{1})`$, and so $`V^{}V`$ is not tame along $`D_Y`$.
Let us now give a definition of tame ramification in the general situation which will be shown in proposition 1.14 to generalize definition 1.1. Let $`X`$ be a scheme, $`YX`$ a closed subscheme and $`U=XY`$ the open complement. For a point $`yY`$ we write $`X_y^{𝑠ℎ}`$ for $`Spec(𝒪_{X,y}^{𝑠ℎ})`$, where $`𝑠ℎ`$ means strict henselization. By abuse of notation, we write $`U_y^{𝑠ℎ}`$ for the base change $`U\times _XX_y^{𝑠ℎ}`$. The scheme $`U_y^{𝑠ℎ}`$ is empty if $`y\overline{U}`$.
###### Definition 1.4
We say that a finite étale covering $`U^{}U`$ is tamely ramified along $`Y`$ if for every point $`yY`$ such that $`U_y^{𝑠ℎ}`$ is nonempty, the base change
$$U^{}\times _UU_y^{𝑠ℎ}U_y^{𝑠ℎ}$$
can be dominated by an étale covering of the form
$$V_1\mathrm{}^\text{.}\mathrm{}\mathrm{}^\text{.}V_rU_y^{𝑠ℎ},$$
such that each $`V_i`$ is a connected étale Galois covering of its image in $`U_y^{𝑠ℎ}`$ and the degree of $`V_i`$ over its image is prime to the characteristic of $`k(y)`$.
The next lemma follows in a straightforward manner from the definition of tame ramification.
###### Lemma 1.5
Let $`X`$ be a scheme, $`UX`$ an open subscheme and $`Y=XU`$. Let $`f:X_1X`$ be a morphism of schemes, $`U_1=f^1(U)`$ and $`Y_1=f^1(Y)`$. If $`U^{}U`$ is a finite étale morphism which is tamely ramified along $`Y`$, then the base change
$$U_1^{}=U^{}\times _UU_1U_1$$
is tamely ramified along $`Y_1`$.
The notion of a finite étale morphism $`U^{}U`$ which is tamely ramified along $`Y=XU`$ is independent from $`X`$ in the following sense.
###### Lemma 1.6
Suppose that $`U`$ is contained as an open subscheme in schemes $`\stackrel{~}{X}`$ and $`X`$ with closed complements $`\stackrel{~}{Y}`$ and $`Y`$, respectively. Assume that there exists a finite morphism $`\pi :\stackrel{~}{X}X`$ making the diagram
$$\begin{array}{ccc}U& \stackrel{}{⸦-\to }& \stackrel{~}{X}\\ & & \pi \\ U& \stackrel{}{⸦-\to }& X\end{array}$$
commutative. Then a finite étale morphism $`U^{}U`$ is tamely ramified along $`\stackrel{~}{Y}`$ if and only if it is tamely ramified along $`Y`$.
Proof: Let $`yY`$ be a point and consider the cartesian diagram
$$\begin{array}{ccc}\stackrel{~}{X}\times _XX_y^{𝑠ℎ}& & X_y^{𝑠ℎ}\\ \text{}& & \text{}\\ U\times _XX_y^{𝑠ℎ}& ==& U_y\end{array}$$
Since $`\pi :\stackrel{~}{X}X`$ is finite, $`\stackrel{~}{X}\times _XX_y^{𝑠ℎ}`$ is strictly henselian, and
$$\stackrel{~}{X}\times _XX_y^{𝑠ℎ}\underset{i}{\mathrm{}^\text{.}}X_{\stackrel{~}{y}_i}^{𝑠ℎ},$$
where the $`\stackrel{~}{y}_i`$ are the finitely many points of $`\stackrel{~}{X}`$ lying above $`y`$. Therefore, $`\pi `$ induces a natural isomorphism
$$\underset{i}{\mathrm{}^\text{.}}U_{\stackrel{~}{y}_i}^{𝑠ℎ}U_y^{𝑠ℎ}.$$
Now the statement of the lemma follows easily from the definition of tame ramification. $`\mathrm{}`$
###### Remark 1.7
Assume that $`X`$ is excellent or, slightly weaker, that all local rings of $`X`$ are Nagata rings. Then there exists a unique maximal scheme $`\stackrel{~}{X}`$ satisfying the conditions of lemma 1.6: the normalization of $`X`$ in $`U`$. This scheme is constructed in the following way. Denoting the open immersion by $`j:UX`$, the sheaf $`j_{}(𝒪_U)`$ is quasicoherent and contains $`𝒪_X`$. The integral closure of $`𝒜`$ of $`𝒪_X`$ in $`j_{}(𝒪_U)`$ is then a coherent sheaf of $`𝒪_X`$-algebras and the normalization of $`X`$ in $`U`$ is defined as $`𝒮pec(𝒜)`$ (cf. \[Mi\], I, §1, proof of th. 1.8).
If $`U^{}U`$ is a tame covering, then, by Zariski’s main theorem, the morphism $`U^{}X`$ factors in the form $`U^{}\stackrel{j}{}X^{}\stackrel{\pi }{}X`$ with $`\pi `$ finite and $`j`$ an open embedding. In general, $`X^{}`$ is not unique but if $`X_1^{}`$ and $`X_2^{}`$ are such schemes, we find a third scheme $`X_3^{}`$ dominating $`X_1^{}`$ and $`X_2^{}`$. If $`X`$ is normal and connected, then the normalization of $`X`$ in the function field of $`U^{}`$ is maximal among the schemes which are finite over $`X`$ and contain $`U^{}`$ as an open dense subscheme. The following lemma follows easily by the same considerations as in the proof of lemma 1.6.
###### Lemma 1.8
Let
$$X_1\stackrel{\pi _1}{}X_2\stackrel{\pi _2}{}X_3$$
be finite surjective morphisms, let $`U_3X_3`$ be an open subscheme and denote by $`U_2`$ and $`U_1`$ the corresponding inverse images. Suppose that $`\pi _1|_{U_1}:U_1U_2`$ and $`\pi _2|_{U_2}:U_2U_3`$ are étale. Furthermore, let $`Y_i`$, $`1i3`$, be the closed complement of $`U_i`$ in $`X_i`$.
Then $`(\pi _2\pi _1)|_{U_1}:U_1U_3`$ is tamely ramified along $`Y_3`$ if and only if $`\pi _1|_{U_1}:U_1U_2`$ is tamely ramified along $`Y_2`$ and $`\pi _2|_{U_2}:U_2U_3`$ is tamely ramified along $`Y_3`$.
Consider the category $`\mathrm{𝐅𝐄𝐭𝐓}/(X,Y)`$ of $`U`$-schemes, finite and étale over $`U`$ with tame ramification along $`Y`$. It is a full subcategory of the category $`\mathrm{𝐅𝐄𝐭}/U`$ of $`U`$-schemes, finite and étale over $`U`$. If $`\overline{x}U`$ is a geometric point of $`U`$, we consider the fibre functor
$$𝐅:\mathrm{𝐅𝐄𝐭𝐓}/(X,Y)\stackrel{}{}(Sets);(VU)Mor_U(\overline{x},V),$$
which is just the restriction of the usual fibre functor induced by $`\overline{x}`$ on $`\mathrm{𝐅𝐄𝐭}/U`$ to $`\mathrm{𝐅𝐄𝐭𝐓}/(X,Y)`$. One observes that the category $`\mathrm{𝐅𝐄𝐭𝐓}/(X,Y)`$ together with this fibre functor satisfies the axiomatic conditions for a Galois theory (see \[SGA1\] Exp.V,4). Thus one obtains a profinite fundamental group
$$\pi _1^t(X,Y;\overline{x})$$
which we call the tame fundamental group of $`U`$ with base point $`\overline{x}`$ and with respect to the embedding $`UX`$. It is a quotient of the étale fundamental group $`\pi _1(U;\overline{x})`$. If $`\overline{x}^{}U`$ is another base point in the same connected component of $`U`$, then the fundamental group $`\pi _1^t(X,Y;\overline{x}^{})`$ is isomorphic to $`\pi _1^t(X,Y;\overline{x})`$, the isomorphism is determined up to an inner automorphism. In the following, $`U`$ will always be connected and we will omit the base point from the notation. If $`X`$ is a scheme of characteristic $`0`$ (i.e. all residue fields have characteristic $`0`$), then
$$\pi _1^t(X,Y)=\pi _1(U).$$
If $`X`$ is normal and connected, then the functor $`\mathrm{𝐅𝐄𝐭}/X\mathrm{𝐅𝐄𝐭}/U`$ is fully faithful, and we obtain surjections
$$\pi _1(U)\pi _1^t(X,Y)\pi _1(X).$$
The next proposition follows in a straightforward manner from the theorem on the purity of the branch locus of Zariski-Nagata (\[SGA2\], Exp.X, th. 3.4).
###### Proposition 1.9
Assume that $`X`$ is regular and connected, and that $`Y`$ is of codimension $`2`$ in $`X`$. Then we have natural isomorphisms
$$\pi _1(U)\pi _1^t(X,Y)\pi _1(X).$$
Let $`U^{}U`$ be a finite étale morphism. The wild locus $`W_{(U^{}U)}`$, i.e. the set of points $`yY`$ such that $`U^{}U`$ is not tamely ramified at $`y`$, is a closed subscheme of $`Y`$. The theorem of Zariski-Nagata on the purity of the branch locus says that if $`X`$ is regular, the ramification locus of a quasifinite dominant morphism to $`X`$ is always pure of codimension $`1`$ in $`X`$. The same is not true of the wild locus, since in example 1.3 we constructed a cyclic cover of a regular scheme of dimension $`2`$ whose wild locus consists of a single closed point, i.e. is of codimension $`2`$. But if $`X`$ is equicharacteristic, we have the
###### Proposition 1.10
Assume that $`X`$ is regular and connected and that all points of $`X`$ have the same residue characteristic. Let $`U^{}U`$ be a finite nilpotent étale covering. Then $`W_{(U^{}U)}`$ is either empty or pure of codimension $`1`$ in $`X`$.
Proof: We may assume that all points on $`X`$ have the common residue characteristic $`p>0`$, because otherwise the wild locus is empty. Let $`K`$ be the function field of $`X`$ and $`L`$ the function field of $`U^{}`$. Then $`L|K`$ is a finite nilpotent extension and $`U^{}=U_L`$ is the normalization of $`U`$ in $`L`$. Writing $`L`$ as a composite $`L=L_1L_2`$, where $`[L_1:K]`$ is prime to $`p`$ and $`[L_2:K]`$ is a $`p`$-power, then $`U_{L_1}U`$ is tamely ramified along $`Y`$. The ramification locus of $`X_{L_2}X`$ is the same as the wild locus $`W_{U_{L_2}U}`$ since $`G(L_2|K)`$ is a $`p`$-group and all points of $`X`$ have residue characteristic $`p`$. Since $`X`$ is regular, the ramification locus of $`X_{L_2}X`$ is either empty or pure of codimension $`1`$ and coincides with $`W_{(U_{L_2}U)}=W_{(U^{}U)}`$. $`\mathrm{}`$
###### Remark 1.11
Under the assumptions of proposition 1.10, let $`D_1,\mathrm{},D_r`$ be the irreducible components of $`Y=XU`$ which are of codimension $`1`$ in $`X`$. Then a finite nilpotent covering $`U^{}U`$ is tamely ramified along $`Y`$ if and only if the discrete valuations $`v_1,\mathrm{},v_r`$ of $`K(X)=K(U)`$ which are associated to $`D_1,\mathrm{},D_r`$ are tamely ramified in $`K(U^{})|K(U)`$. This justifies the ‘naive’ definition of tame ramification used in \[S-S1\].
¿From now on assume that $`X`$ is a normal connected scheme and we denote the function field of $`X`$ by $`K`$. Every connected finite étale covering $`\stackrel{~}{X}X`$ coincides with the normalization of $`X`$ in some finite separable extension $`L`$ of $`K`$. Therefore (omitting suitable chosen base points from the notation), we obtain a natural surjection $`G(\overline{K}|K)=\pi _1(Spec(K))\stackrel{}{-}\pi _1(X)`$ and an isomorphism
$$G(\overline{K}|K)=\underset{}{lim}_{UX}\pi _1(U),$$
where $`U`$ runs through the open subschemes of $`X`$.
Let us collect some facts on decomposition and inertia of integrally closed domains (see \[B-Co\] Ch. V, §§2,3). Let $`A`$ be an integrally closed domain with quotient field $`K`$, $`L`$ a Galois extension of $`K`$ and $`B`$ the integral closure of $`A`$ in $`L`$. Let $`𝔓`$ be a prime ideal of $`B`$ and $`𝔭=𝔓A`$ the prime ideal in $`A`$ lying under $`𝔓`$. Let $`k(𝔓)=B_𝔓/𝔓B_𝔓`$ and $`k(𝔭)=A_𝔭/𝔭A_𝔭`$ be the residue fields of $`𝔓`$ and $`𝔭`$, respectively. Let $`G=G(L|K)`$ be the Galois group. Then $`G`$ acts transitively on the set of prime ideals of $`B`$ lying over $`𝔭`$. Then one has the following subgroups in the Galois group.
\- $`Z=Z_𝔓(L|K)=\{\sigma G|\sigma (𝔓)=𝔓\}`$ \- the decomposition group,
\- $`T=T_𝔓(L|K)=\{\sigma G_𝔓|\overline{\sigma }=𝑖𝑑:B/𝔓B/𝔓\}`$ \- the inertia group.
$`T`$ is a normal subgroup in $`Z`$. Let $`A^T`$ and $`A^Z`$ be the integral closures of $`A`$ in $`L^T`$ and $`L^Z`$ respectively. Put $`𝔓^T=𝔓A^T`$ and $`𝔓^Z=𝔓A^Z`$. For a proof of the following proposition see \[Ra\], X, th.2.
###### Proposition 1.12
If $`L=\overline{K}`$ is the separable closure of $`K`$, then
(i)$`𝔓`$ is the only prime ideal in $`B`$ extending the prime ideal $`𝔓^Z`$ in $`A^Z`$,
(ii)$`(A_{𝔓^Z}^Z,𝔓^Z)`$ is the henselization of $`(A_𝔭,𝔭)`$, in particular, $`k(𝔓^Z)=k(𝔭)`$,
(iii)$`(A_{𝔓^T}^T,𝔓^T)`$ is the strict henselization of $`(A_𝔭,𝔭)`$, in particular, $`k(𝔓^T)`$
is the separable closure of $`k(𝔭)`$.
Assume that $`X`$ is normal and connected. Then the normalization of $`X`$ in a finite separable extension of its function field is finite over $`X`$. As is well known, for every point $`x`$ on $`X`$, the scheme $`X_x^{𝑠ℎ}`$ is noetherian, normal and connected.
###### Corollary 1.13
Let $`X`$ be a normal connected scheme, $`UX`$ an open subscheme and $`Y=XU`$ the closed complement. Let $`L`$ be a finite Galois extension of the function field $`K`$ of $`X`$ and let $`U_L`$, $`X_L`$ be the normalizations of $`U`$ and $`X`$ in $`L`$. Then the following are equivalent
(i)$`U_LU`$ is étale and tamely ramified along $`Y`$,
(ii)For each point $`𝔓X_L`$ the following holds, where $`p`$ denotes the residue characteristic of $`𝔓`$:
a)$`T_𝔓(L|K)`$ is of order prime to $`p`$
b)$`T_𝔓(L|K)=0`$ if $`𝔓U_L`$.
Proof: The only nontrivial point is to see that $`U_LU`$ is flat, but this is well known (cf. \[Mi\], I,§3, th.3.21). $`\mathrm{}`$
###### Proposition 1.14
Let $`X`$ be a regular connected scheme and $`DX`$ a normal crossing divisor. Then both notions of tame ramification coincide.
Proof: Let (notation as above) $`U_LU`$ be tame in the sense of definition 1.4 and assume that $`v`$ is a discrete valuation on $`K`$ which corresponds to the generic point $`P`$ of an irreducible component of $`D`$. Let (for any embedding to the separable closure of $`K`$) $`K_P^{𝑠ℎ}`$ be the quotient field of the strict henselization of the discrete valuation ring $`𝒪_{X,P}`$. Then, by definition, $`LK_P^{𝑠ℎ}`$ is a product of fields each of which is contained in a finite Galois extension of degree prime to $`p=\text{char}k(P)`$ of $`K_P^{𝑠ℎ}`$. Therefore $`v`$ is tamely ramified in $`L|K`$. On the other hand, if all these discrete valuations are tamely ramified in $`L|K`$, then it follows from the Generalized Abhyankar’s lemma (theorem 1.2) that $`U_LU`$ is tamely ramified in the sense of definition 1.4. $`\mathrm{}`$
Remark: Corollary 1.13 asserts that for coverings of normal schemes our definition of tameness coincides with that given in \[Ab\]. Furthermore (in the situation of corollary 1.13), if $`U_LU`$ is tame, then $`X_LX`$ is numerically tame in the sense of \[C-E\]. It is also not difficult to show the inverse implication.
## 2 Valuations and tame ramification
The aim of this section is to relate tame covers of schemes of finite type over $`Spec()`$ to discrete valuations of higher rank. The quotient fields of the henselizations of such valuation rings are called higher dimensional henselian fields and are the basic constituents of Kato’s and Saito’s higher dimensional class field theory \[K-S\].
Let us recall some facts from valuation theory (see \[B-Co\], \[Z-S\], Ch. VI). Let $`K`$ be a field. A subring $`VK`$ is called valuation ring if $`xKVx^1V`$. A valuation ring is local with maximal ideal $`𝔪_V=\{xV|x^1V\}`$ and is integrally closed in $`K`$. The quotient $`\mathrm{\Gamma }_V=K^\times /V^\times `$ is a totally ordered abelian group ($`\overline{x}\overline{y}x^1yV`$) and is called the value group of $`V`$. Usually, one extends the natural map $`v:K^\times \mathrm{\Gamma }_V,`$ to $`K`$ by setting $`v(0)=\mathrm{}`$. As is well known, the valuation ring $`V`$ represents an equivalence class of abstractly defined valuations of $`K`$; the valuation $`v`$ is a natural representative of this class.
If $`L|K`$ is a finite separable field extension, then there are at least one and at most finitely many valuations $`w`$ of $`L`$ extending $`v`$, i.e. such that $`WK=V`$, where $`WL`$ is the valuation ring of $`w`$. We have an induced injective homomorphism $`\mathrm{\Gamma }_V\mathrm{\Gamma }_W`$. The index of $`\mathrm{\Gamma }_V`$ in $`\mathrm{\Gamma }_W`$ is called the ramification index and is usually denoted by $`e=e_{w|v}`$. We also have an associated field extension $`k(v)=V/𝔪_Vk(w)=W/𝔪_W`$ whose degree is denoted by $`f=f_{w|v}`$. The inequality
(1)
$$\underset{w|v}{}e_{w|v}f_{w|v}[L:K],$$
(see \[B-Co\] Ch.6, §8.3. th.1) shows, in particular, that all these numbers are finite. $`v`$ is called defectless in $`L`$ if equality holds in (1). By \[B-Co\], VI, §8.5, cor.2, a discrete valuation is defectless in a finite separable extension.
Let $`L`$ be a Galois extension of $`K`$, $`W`$ a valuation ring of $`L`$ and $`V=WK`$. Let $`𝔓`$ be the unique maximal ideal of $`W`$. Then, besides the inertia and decomposition group, we have the ramification group
$$R_𝔓=R_𝔓(L|K)=\{\sigma T_𝔓(L|K)|\frac{\sigma x}{x}1mod𝔓\text{ for all }0xL\}.$$
A proof of the following proposition can be found in \[End\], §20.
###### Proposition 2.1
If the residue characteristic of $`V`$ is zero, then $`R_𝔓=0`$ and $`T_𝔓`$ is abelian. If the residue characteristic of $`V`$ is $`p>0`$, then
(i)$`R_𝔓`$ is a normal subgroup in $`Z_𝔓`$,
(ii)$`R_𝔓`$ is the unique $`p`$-Sylow subgroup in $`T_𝔓`$,
(iii)$`T_𝔓/R_𝔓`$ is an abelian group of order prime to $`p`$.
The decomposition, inertia, ramification groups of different valuation rings $`W`$ of $`L`$ extending $`V`$ are conjugate in $`G(L|K)`$. One says that $`W|V`$ (resp. $`w|v`$) is tamely ramified if the ramification group is trivial. $`V`$ is tamely ramified in $`L|K`$ if $`W|V`$ is tamely ramified for one (every) valuation ring $`W`$ of $`L`$ extending $`V`$. One says that $`V`$ is tamely ramified in a separable extension if it is tamely ramified in the Galois closure of this extension.
If $`V`$ is tamely ramified in $`L|K`$, then it is defectless in $`L|K`$ (see \[End\], cor. 20.22). A valuation $`v`$ on $`K`$ is called discrete valuation of rank $`n`$ if its value group is (as an ordered group) isomorphic to $`\times \mathrm{}\times `$ ($`n`$ factors) with the lexicographic order. By the unspecified term discrete valuation we always mean discrete valuation of rank $`1`$, i.e. a discrete valuation in the usual sense.
Let $`V`$ be a discrete valuation ring of rank $`n`$ of $`K`$. Then $`V`$ has $`n`$ distinct prime ideals $`𝔭_0𝔭_1\mathrm{}𝔭_n`$. For $`0in`$, let $`V_i`$ be the localization of $`V`$ at $`𝔭_i`$ and let $`k_i=k(\pi _i)`$ be the residue field of $`V_i`$. The $`V_i`$ are also discrete valuation rings (of rank $`ni`$) of $`K`$ and
$$V=V_0V_1\mathrm{}V_n=K.$$
For $`0i<jn`$ the image of $`V_i`$ in $`k_j`$ is a discrete valuation ring of rank $`ji`$. For $`0i<n`$, we denote the image of $`V_i`$ in $`k_{i+1}`$ by $`\overline{V}_i`$. Then $`\overline{V}_i`$ is a usual (i.e. rank $`1`$) discrete valuation ring with quotient field $`k_{i+1}`$ and residue field $`k_i`$. By \[B-Co\], Ch. VI, §7, ex.6d, the ring $`V`$ is henselian if and only if all $`\overline{V}_i`$, $`0i<n`$ are henselian (see also \[Ri\], ch. F, prop.9).
Let $`VK`$ be a discrete valuation ring of rank $`n`$ and let $`L|K`$ be a finite separable extension. Let $`w|v`$ be a valuation of $`L`$ extending $`v`$. Then also $`w`$ is a discrete valuation of rank $`n`$. Assume that the residue characteristic of $`v`$ is $`p>0`$.
###### Lemma 2.2
Let $`v`$ be a discrete valuation of rank $`n`$ of a field $`K`$ and let $`L|K`$ be a separable extension of $`K`$. Then $`v`$ is tamely ramified in $`L|K`$ if and only if for every extension $`w`$ of $`v`$ to $`L`$ the ramification index $`e_{w|v}`$ is prime to $`p`$ and the successive residue extensions of $`\overline{W}_i|\overline{V}_i`$ are separable for $`0i<n`$.
Proof: We may assume that $`v`$ is strictly henselian and that $`L|K`$ is finite. In particular, $`v`$ has a unique extension $`w`$ to $`L`$. Assume that $`v`$ is tamely ramified in $`L|K`$. By proposition 2.1, (iii), $`L|K`$ is an abelian extension of degree prime to $`p`$. Since $`v`$ is defectless in $`L|K`$, the ramification index is prime to $`p`$. Furthermore, the degrees of the successive residue extensions are prime to $`p`$, and so these extensions are separable.
It remains to show the other direction. Since the successive residue extensions are separable, and since discrete valuations of rank $`1`$ are defectless in separable extensions, one shows inductively (cf. \[B-Co\] Ch.VI §7 ex.5 or \[Ri\], G th.5) that $`v`$ is defectless in $`L|K`$. In particular, $`[L:K]=e_{w|v}`$ is prime to $`p`$. By proposition 2.1 (ii),(iii), we conclude that $`L|K`$ is contained in a Galois extension of degree prime to $`p`$ and is therefore tamely ramified. $`\mathrm{}`$
We call a field $`k`$ an $`𝒏`$-dimensional henselian field if it is the quotient ring of a henselian discrete valuation ring of rank $`n`$ with finite residue field. Equivalently, one can give an inductive definition: A $`0`$-dimensional henselian field is just a finite field. An $`(n+1)`$-dimensional henselian field is a field which is henselian under a discrete (rank $`1`$) valuation whose residue field is an $`n`$-dimensional henselian field. So an $`n`$-dimensional henselian field $`\kappa `$ comes along with a sequence $`\kappa _{n1},\mathrm{},\kappa _0`$ of residue fields, where each $`\kappa _i`$, $`i=n1,\mathrm{},0`$ is an $`i`$-dimensional henselian field. In particular, $`\kappa _0`$ is finite.
We call (by abuse of notation) an extension of $`n`$-dimensional henselian fields tamely ramified, if it is tamely ramified with respect to the discrete valuations of rank $`n`$. The next result follows easily by induction from the corresponding $`1`$-dimensional result.
###### Proposition 2.3
Let $`\kappa `$ be a $`n`$-dimensional henselian field with last residue field $`\kappa _0`$ of characteristic $`p>0`$ and let $`\kappa ^t`$ be the maximal tamely ramified extension of $`\kappa `$. Then there is a natural isomorphism
$$G(\kappa ^t|\kappa )\left(\underset{ntimes}{\underset{}{\widehat{}^{(p^{})}(1)\times \mathrm{}\times \widehat{}^{(p^{})}(1)}}\right)G(\overline{\kappa }_0|\kappa _0),$$
where $`(\widehat{})^{(p^{})}`$ denotes the prime-to-$`p`$ part of $`\widehat{}`$ and $`(1)`$ means the first Tate-twist with respect to the cyclotomic character.
¿From now on let $`X`$ be a reduced, separated, equidimensional scheme of finite type over $`Spec()`$, in particular (\[Ma\], Ch.13, no.34), $`X`$ is excellent. Let $`d=dim(X)=dim_{Krull}(X)`$. A Parshin-chain $`P`$ on $`X`$ is a sequence $`(P_0,P_1,\mathrm{},P_d)`$ of points on $`X`$ such that
$$\overline{\{P_0\}}\overline{\{P_1\}}\mathrm{}\overline{\{P_d\}}=X$$
and $`dim\overline{\{P_i\}}=i`$ for $`i=0,\mathrm{},d`$.
An inductive localization-henselization procedure, which was proposed by Parshin with completion instead of henselization and which is described in \[K-S\], §3, provides us with a functor $`(X,P)\kappa _P^h`$
$$\left(\begin{array}{c}\text{Reduced schemes }X\text{ of}\\ \text{finite type over }Spec()\\ \text{with a Parshin-chain }P\\ \text{of length }d=dimX\end{array}\right)\left(\begin{array}{c}\text{finite products}\\ \text{of }d\text{-dimensional}\\ \text{henselian fields}\end{array}\right).$$
Let us briefly recall this construction. First one takes the henselization $`𝒪_{X,P_0}^h`$ of the local ring of $`X`$ at $`P_0`$. Then one considers the finitely many prime ideals in this ring which lie over $`P_1`$, passes to the product of the henselizations of the localizations with respect to this prime ideals, and so on.
If in the Parshin-chain $`P`$ each $`P_i`$ is a regular point of $`\overline{\{P_{i+1}\}}`$, then $`\kappa _P^h`$ is a $`d`$-dimensional henselian field rather than a finite product of such fields.
Assume that $`X`$ is integral and let $`K`$ be the function field of $`X`$. We say that a discrete valuation ring $`V`$ of rank $`d`$ in $`K`$ dominates a Parshin-chain $`P`$ if $`V_i`$ dominates $`𝒪_{X,P_i}`$ for $`i=1,\mathrm{},d`$. For a proof of the following proposition see \[K-S\], prop.3.3.
###### Proposition 2.4
Under the above assumptions,
$$\kappa _P^h=\underset{V}{}Quot(V^h),$$
where $`V`$ ranges over all discrete valuations of rank $`d`$ which dominate $`P`$ and $`Quot(V^h)`$ denotes the quotient field of the henselization $`V^h`$ of $`V`$.
###### Proposition 2.5
Let $`X`$ be a $`d`$-dimensional, regular, connected scheme of finite type over $`Spec()`$. Let $`UX`$ be an open subscheme and $`Y=XU`$ the complement. Let $`L`$ be a finite nilpotent extension of the function field $`K`$ of $`X`$. Then the following are equivalent.
(i)$`U_LU`$ is étale and tamely ramified along $`Y`$.
(ii)Every discrete valuation $`v`$ of rank $`d`$ in $`K`$ which dominates a Parshin chain on $`X`$ is tamely ramified in $`L`$, and unramified if the dominated chain in contained in $`U`$.
(iii)For every Parshin-chain $`P`$ on $`X`$ the extension of (finite products of) $`d`$-dimensional henselian fields $`\kappa _P^hL|\kappa _P^h`$ is tamely ramified, and unramified if $`P`$ is contained in $`U`$.
Proof: The implication (i)$``$(ii) is obvious and the equivalence (ii)$``$(iii) follows from proposition 2.4. Suppose that (ii) holds. Let $`P_{d1}`$ be a point of codimension $`1`$ on $`U`$. We extend $`P_{d1}`$ to a Parshin-chain $`(P_0,\mathrm{},P_{d1},P_d)`$ on $`U`$. Let $`V`$ be a discrete valuation ring in $`K`$ dominating $`P`$. Then (notation as above), $`V_{d1}`$ dominates and hence is equal to $`𝒪_{U,P_{d1}}`$. Since $`V`$ is unramified in $`L|K`$, so is $`V_{d1}`$. By the theorem of Zariski-Nagata on the purity of the branch locus, we conclude that $`U_LU`$ is étale.
It remains to show that $`U_LU`$ is tamely ramified along $`Y`$. Since $`L|K`$ is nilpotent, we can easily reduce to the case that the Galois group is a finite $`p`$-group for some prime number $`p`$. Let $`x_1,\mathrm{},x_n`$ be the finitely many codimension $`1`$ points which ramify in $`X_LX`$. Since the set of points in $`X`$ where $`U_LU`$ is not tame is closed, it suffices to show that it is tamely ramified at every closed point $`yY`$.
$`U_LU`$ is tamely ramified at every point of residue characteristic different from $`p`$ and at every point which is not contained in $`_{i=1}^n\overline{\{x_i\}}`$.
We show that a closed point $`y`$ of residue characteristic $`p`$ and with $`y\overline{\{x_i\}}`$ for some $`1in`$ does not exist. Assume that $`y`$ is such a point. Choose a Parshin-chain $`P=(P_0,\mathrm{},P_d)`$ on $`X`$ with $`P_0=y`$ and $`P_{d1}=x_i`$ (we find such a chain because $`X`$ is catenary). Choose a discrete valuation ring $`V`$ of rank $`d`$ in $`K`$ dominating $`P`$. Then (for appropriately chosen embeddings to $`\overline{K}`$) the quotient field $`K_{x_i}^{𝑠ℎ}`$ of the strict henselization of $`𝒪_{X,x_i}`$ contains $`Quot(V^{𝑠ℎ})`$. By assumption $`LQuot(V^{𝑠ℎ})|Quot(V^{𝑠ℎ})`$ is a field extension of degree prime to $`p`$, hence trivial. But, since $`x_i`$ is ramified in the $`p`$-extension $`L|K`$, $`LK_{x_i}^{𝑠ℎ}|K_{x_i}^{𝑠ℎ}`$ is nontrivial. This yields the required contradiction. $`\mathrm{}`$
From the proof of the last proposition, we obtain the
###### Corollary 2.6
Assume $`X`$ is connected but only normal. Let $`L|K`$ be a finite Galois extension of $`p`$-power degree. If $`U_LU`$ is étale and tamely ramified along $`Y`$, then $`L|K`$ is unramified at every discrete valuation associated to a codimension $`1`$ point of $`X`$ which either lies on $`U`$ or whose closure in $`X`$ contains a point of residue characteristic $`p`$. If $`X`$ is regular, the inverse implication is also true.
###### Proposition 2.7
Suppose that $`U`$ is regular and let $`X_1`$ and $`X_2`$ be regular schemes which are proper over $`Spec()`$ and which contain $`U`$ as an open subscheme. Let $`Y_i=X_iU`$, $`i=1,2`$, be the closed complements. Let $`U^{}U`$ be a finite nilpotent étale covering. Then $`U^{}U`$ is tamely ramified along $`Y_1`$ if and only if it is tamely ramified along $`Y_2`$.
Proof: We may suppose that the degree of $`U^{}U`$ is a power of some prime number $`p`$. Let $`X_3`$ be a normal scheme which is proper over $`Spec()`$, contains $`U`$ as an open subscheme and has proper surjective morphisms $`\pi _i:X_3X_1`$, $`i=1,2`$ making the diagrams
$$\begin{array}{ccc}U& \stackrel{}{⸦-\to }& X_3\\ & & \pi _i\\ U& \stackrel{}{⸦-\to }& X_i\end{array}$$
commutative. (The existence of such a $`X_3`$ is well known, for instance, take $`X_3`$ as the normalization of the closure of the image of $`U`$ in $`X_1\times _{}X_2`$.) If $`U^{}U`$ is tamely ramified along $`Y_1`$, then it is tamely ramified along $`Y_3=X_3U`$ (see lemma 1.5). Let $`L`$ be the function field of $`U^{}`$. Then, by corollary 2.6, $`L|K`$ is tamely ramified at every discrete valuation on $`L`$ which is defined by a codimension $`1`$ point of $`X_3`$ which lies on $`Y_3`$ and whose closure contains a point of residue characteristic $`p`$. This set of valuation contains the set of discrete valuation on $`L`$ which are defined by a codimension $`1`$ point of $`X_2`$ which lies on $`Y_2`$ and whose closure contains a point of residue characteristic $`p`$. Applying corollary 2.6 again, we conclude that $`U^{}U`$ is tamely ramified along $`Y_2`$. $`\mathrm{}`$
Theorem 1 of the introduction follows immediately from proposition 2.7. In particular, the abelianized tame fundamental group of a regular scheme of finite type over $`Spec()`$ does not depend on the choice of a compactification. So it is justified to use the notation $`\pi _1^t(X)^{ab}`$ for $`\pi _1^t(\overline{X},\overline{X}X)^{ab}`$ for this group.
## 3 A finiteness result
The aim of this section is to prove the following theorem 3.1 (=th.2 from the introduction). In the proof we essentially use a finiteness result of N. Katz and S. Lang on relative étale fundamental groups (\[K-L\], th.1).
###### Theorem 3.1
Let $`𝒪`$ be the ring of integers in a finite extension $`k`$ of $``$. Let $`X`$ be a flat $`𝒪`$-scheme of finite type whose geometric generic fibre $`X_𝒪\overline{k}`$ is connected. Assume that $`X`$ is normal and that the morphism $`XSpec(𝒪)`$ is surjective. Let $`U`$ be an open subscheme of $`X`$. Then the abelianized tame fundamental group $`\pi _1^t(X,XU)^{ab}`$ is finite.
For the proof of theorem 3.1 we need the following
###### Lemma 3.2
Let $`A`$ be a strictly henselian discrete valuation ring with perfect (hence algebraically closed) residue field and with quotient field $`k`$. Let $`k_{\mathrm{}}|k`$ be a $`_p`$-extension. Let $`K|k`$ be a regular field extension and let $`BK`$ be a discrete valuation ring dominating $`A`$. Then $`B`$ is ramified in $`Kk_{\mathrm{}}|K`$.
Proof: For each $`n0`$, let $`k_n|k`$ be the unique subextension of degree $`p^n`$ in $`k_{\mathrm{}}|k`$. Let $`A_n`$ (resp. $`B_n`$) denote the normalization of $`A`$ (resp. $`B`$) in $`k_n`$ (resp. $`Kk_n`$). $`A_n`$ is again a strictly henselian discrete valuation ring. $`B_n`$ is a semi-local Dedekind domain.
Suppose that $`B`$ is unramified in $`Kk_{\mathrm{}}`$. Fix an $`n`$. Let $`𝔭`$ be the prime ideal of $`B`$ and let $`𝔭_1,\mathrm{},𝔭_g`$ be the prime ideals of $`B_n`$. Since $`B_n|B`$ is Galois, all $`𝔭_i`$ have a common inertia index $`f=f_n`$, so $`N_{B_n|B}(𝔭_i)=𝔭^f`$ for $`i=1,\mathrm{},g`$. Since $`B_n|B`$ is étale, we have $`p^n=[k_n:k]=gf`$.
Let $`\pi _n`$ be a uniformizer of $`A_n`$. Since $`A_n|A`$ is totally ramified, we have
$$v_A(N_{k_n|k}(\pi _n))=1.$$
Considered as an element of $`B_n`$, the $`𝔭_i`$-valuation of $`\pi _n`$ is positive and independent of $`i`$, $`1ig`$. Denoting this positive number by $`a`$, we have the following equality of ideals in $`B_n`$
$$(\pi _n)=(𝔭_1\mathrm{}𝔭_g)^a.$$
Hence $`N_{B_n|B}((\pi _n))=𝔭^{afg}`$ and therefore
$$v_B(N_{k_n|k}(\pi _n))=ap^np^n.$$
On the other hand, we have
$$v_B(N_{k_n|k}(\pi _n))=v_B(\pi _0)v_A(N_{k_n|k}(\pi _n))=v_B(\pi _0),$$
where $`\pi _0`$ is any uniformizer in $`A=A_0`$. Since $`n`$ was arbitrary, the inequality $`v_B(\pi _0)p^n`$ yields a contradiction. $`\mathrm{}`$
Proof of theorem 3.1: Since $`X`$ is normal, for any open subscheme $`V`$ of $`U`$ the natural homomorphism $`\pi _1(V)\pi _1(U)`$ is surjective. Therefore also the homomorphism
$$\pi _1^t(X,XV)^{ab}\pi _1^t(X,XU)^{ab}$$
is surjective and so we may replace $`U`$ by a suitable open subscheme and assume that $`U`$ is smooth over $`\overline{S}=Spec(𝒪)`$. Let $`S\overline{S}`$ be the image of $`U`$. Consider the commutative diagram
$$\begin{array}{ccccccc}0& & Ker(U/S)& & \pi _1(U)^{ab}& & \pi _1(S)^{ab}\\ & & & & & & \\ 0& & Ker^t(U/S)& & \pi _1^t(X,XU)^{ab}& & \pi _1^t(\overline{S},\overline{S}S)^{ab}\end{array}$$
where the groups $`Ker(U/S)`$ and $`Ker^t(U/S)`$ are defined by the exactness of the corresponding rows, and the two right vertical homomorphisms are surjective. By a theorem of N. Katz and S. Lang (\[K-L\], th.1), the group $`Ker(U/S)`$ is finite. By classical one-dimensional class field theory, the group $`\pi _1^t(\overline{S},\overline{S}S)^{ab}`$ is finite (it is the Galois group of the ray class field of $`k`$ with modulus $`_{𝔭S}𝔭`$). The kernel of $`\pi _1(S)^{ab}\pi _1^t(\overline{S},\overline{S}S)^{ab}`$ is generated by the ramification groups of the primes of $`\overline{S}`$ which are not in $`S`$. Denoting the product of the residue characteristics of these primes by $`N`$, we see that $`\pi _1(S)^{ab}`$ is the product of a finite group and a topologically finitely generated pro-$`N`$ group. Therefore the same is also true for $`\pi _1(U)^{ab}`$ and for $`\pi _1^t(X,XU)^{ab}`$. By the snake lemma, it will suffice to show that the cokernel $`C`$ of the induced map $`Ker(U/S)Ker^t(U/S)`$ is a torsion group.
Let $`K`$ be the function field of $`X`$ and let $`k_1`$ be the maximal abelian extension of $`k`$ such that the normalization $`U_{Kk_1}`$ of $`U`$ in the composite $`Kk_1`$ is étale over $`U`$ and tamely ramified along $`Y`$. By \[K-L\], lemma 2, (2), the normalization of $`S`$ in $`k_1`$ is ind-étale over $`S`$. Let $`k_2|k`$ be the maximal subextension of $`k_1|k`$ such that the normalization $`S_{k_2}`$ of $`S`$ in $`k_2`$ is étale over $`S`$ and tamely ramified along $`\overline{S}S`$. Then $`G(k_2|k)=\pi _1^t(\overline{S},\overline{S}S)^{ab}`$ and, by the snake lemma, $`CG(k_1|k_2)`$.
In order to show that $`C`$ is a torsion group, we therefore have to show that $`k_1|k_2`$ does not contain a $`_p`$-extension of $`k_2`$ for any prime number $`p`$. Since $`k_2|k`$ is a finite extension and $`k_1|k`$ is abelian, this is equivalent to the assertion that $`k_1|k`$ contains no $`_p`$-extension of $`k`$ for any prime number $`p`$. Let $`p`$ be a prime number and suppose that $`k_{\mathrm{}}|k`$ is a $`_p`$-extension such that the normalization $`U_{Kk_{\mathrm{}}}`$ is ind-étale over $`U`$ and ind-tamely ramified along $`Y`$. A $`_p`$-extension is unramified outside $`p`$ and at least ramified at one prime dividing $`p`$, see e.g. \[NSW\], (10.3.20)(ii). Since the normalization of $`S`$ in $`k_1`$ is ind-étale, we may suppose that $`p|N`$. Let $`k^{}`$ be the maximal unramified subextension of $`k_{\mathrm{}}|k`$ and let $`\overline{S}^{}`$ be the normalization of $`\overline{S}`$ in $`k^{}`$. Then the base change $`X^{}=X\times _{\overline{S}}\overline{S}^{}X`$ is étale. Hence $`X^{}`$ is normal and the pre-image $`U^{}`$ of $`U`$ is smooth and geometrically connected over $`k^{}`$. So, after replacing $`k`$ by $`k^{}`$, we may suppose that $`k_{\mathrm{}}|k`$ is totally ramified at a prime $`𝔭|p`$, $`𝔭\overline{S}S`$.
Let $`𝒪_𝔭`$ be the local ring of $`\overline{S}`$ at $`𝔭`$. After a base change to the strict henselization $`A`$ of $`𝒪_𝔭`$, we arrive at the following situation (several letters get a new meaning):
(1)a strictly henselian discrete valuation ring $`A`$ with perfect residue field of characteristic $`p`$ and with quotient field $`k`$ of characteristic $`0`$
(2)a flat connected and normal $`A`$-scheme of finite type $`X`$ which projects surjectively to $`Spec(A)`$
(3)the function field $`K`$ of $`X`$ is a regular extension of $`k`$ (i.e. $`k`$ is algebraically closed in $`K`$)
(4)an open subscheme $`UX`$ which is smooth over $`K`$
(5)a $`_p`$-extension $`k_{\mathrm{}}|k`$ (in which $`A`$ is totally ramified)
(6)The normalization of $`U`$ in $`Kk_{\mathrm{}}`$ is ind-étale over $`U`$ and ind-tamely ramified along $`Y=XU`$.
Let us show that such a situation cannot occur. Let $`P`$ be the generic point of an irreducible component of the special fibre of $`X`$ over $`A`$. Then $`B=𝒪_{X_P}`$ is a discrete valuation ring dominating $`A`$. Now lemma 3.2 shows that $`B`$ ramifies in $`Kk_{\mathrm{}}`$. Therefore the order of the inertia group of $`P`$ in $`Kk_{\mathrm{}}`$ is divisible by $`p=\text{char}k(P)`$. By corollary 1.13, this contradicts assumption (6). $`\mathrm{}`$
In the special case $`U=X`$ we obtain the following corollary.
###### Corollary 3.3
Let $`𝒪`$ be the ring of integers in a finite extension $`k`$ of $``$. Let $`X`$ be a flat $`𝒪`$-scheme of finite type whose geometric generic fibre $`X_𝒪\overline{k}`$ is connected. Assume that $`X`$ is normal and that the morphism $`XSpec(𝒪)`$ is surjective. Then the abelianized étale fundamental group $`\pi _1(X)^{ab}`$ is finite.
Alexander Schmidt, Mathematisches Institut, Universität Heidelberg, Im Neuenheimer Feld 288, 69120 Heidelberg, Deutschland
e-mail: schmidt@mathi.uni-heidelberg.de
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# Energy levels and far-infrared spectroscopy for two electrons in a semiconductor nanoring
## I Introduction
Rapid progress in nanostructure technology has made it possible to fabricate various nanometer quantum devices which have potential applications in microelectronics. Such ultrasmall devices contain only a few electrons and the electron-electron interaction is proposed to be of great importance to theirs energy level structures and optical properties. It leads to a number of new quantum phenomena. One of the most interesting phenomena is the spin oscillation of the ground state in an external magnetic field, which is due to the interplay between three energies: the confinement potential, the Zeeman energy, and the electron-electron interaction. The simplest case, a pillar few electrons quantum dot (QD) with a parabolic potential, has been extensively investigated, and the spin oscillation in a QD is well understood. The theoretical predictions are also confirmed in the experiment through the conductance measurements in the finite drain-source voltage regime.
The semiconductor quantum ring is another interesting example. In 1993, D. Mailly et al. measured the persistent currents in a mesoscopic single GaAs ring induced by a magnetic flux threading the interior of ring, and the experimental results have attracted a lot of theoretical interests. One of the basic questions addressed by many theoretical explanation is concerned with the role of the electron-electron interaction. For a narrow-width rings, an adiabatic approximation allows one to decouple the radial motion from the angular motions and to arrive at analytical solutions for the wave functions and the energy spectra. As shown by L. Wendler et al., the interplay of the Coulomb repulsion between the electrons and the confining potential forms a relatively rigid rotator with internal azimuthal excitations and confined radial motions, i.e., the picture of a rotating Wigner molecule. However, at the moment, the influence of the electron-electron interaction in a finite-width and nanoscopic quantum ring, i.e., nanoring, is still less well understood.
Very recently, using the self-assembly techniques, A. Lorke and collaborators demonstrate the realization of nanoscopic semiconductor quantum rings inside a completed field-effect transistor (FET) structure. Quite different from the conventional sub-micron mesoscopic structures, the nanorings are in the true quantum limit. By applying two complementary spectroscopic techniques, capacitance-voltage (CV) spectra and far-infrared (FIR) spectroscopy, they investigate both the ground state transition and excitation’s properties of these two-electron nanorings in a magnetic field perpendicular to the plane of rings. Although the main experimental results can be qualitatively explained by the single electron picture, some contradictions remain, i.e. the coulomb interaction energy estimated roughly $`20`$ meV, is too large to be ignored safely. Hence, more in-depth theoretical works are desirable, especially in view of the existences of the very strong electron-electron coulomb interaction.
In this paper, we would like to study the energy levels and FIR spectroscopy of a two-electron nanoring, and pay special attention to the effects of the Coulomb interaction. First of all, for a ring-like confinement potential, the total Hamiltonian cannot be separated into the center-of-mass and the relative-motion terms. We develope a new theoretical method to handle this non-separability. It consists of the well-known series solution method, which is effective to solve the single particle problem, and the exact diagonalization method. On the other hand, we show that the electron-electron interaction can change the energy levels significantly. An obvious feature induced by the interaction is the intersection between the lower levels. It presents the spin oscillation of the ground state in an external field. Moreover, the results obtained by A. Lorke et al. can be explained more realistically in our model. Further, the generalized Kohn theorem breaks down in the FIR spectroscopy due to the mixing between the center-of-mass and the relative-motion modes as mentioned above. We outline here two prominent features of the FIR spectroscopy caused by electron-electron interaction: splitting and dis-continuous drops of the resonance energies, which are suggested to be detectable with circularly polarized far-infrared light.
The organization of this paper is as follows. In the next section, we introduce the model and Hamiltonian briefly. In Sec III, the main procedures of the numerical calculation are outlined. Sec IV is devoted to the results and discussions, followed by a summary in Sec V.
## II The model and Hamiltonian
The nanoring is considered to have two electrons with an effective conduction-band-edge mass $`m_e^{}`$ moving in the $`xy`$ plane, and a ring-like confining potential can be introduced $`U(\stackrel{}{r})=\frac{1}{2}m_e^{}\omega _0^2\left(rR_0\right)^2`$, where $`\omega _0`$ is the characteristic frequency of the radial confinement and $`R_0`$ is the ring’s radius. This system is subjected to a perpendicular uniform magnetic field which is described by a vector potential $`\stackrel{}{A}\left(\stackrel{}{r}\right)=\frac{1}{2}\stackrel{}{B}\times \stackrel{}{r}`$ in the symmetric (or circular) gauge. The resulting Hamiltonian is given by
$$=\underset{i=1,2}{}\left\{\frac{1}{2m_e^{}}\left(\stackrel{}{p}_i+e\stackrel{}{A}\left(\stackrel{}{r}_i\right)\right)^2+U(\stackrel{}{r}_i)\right\}+\frac{e^2}{4\pi \epsilon _0\epsilon _r\left|\stackrel{}{r}_1\stackrel{}{r}_2\right|},$$
(1)
where $`\stackrel{}{r}_i=(x_i,y_i)`$ and $`\stackrel{}{p}_i=i\mathrm{}\stackrel{}{\mathrm{}}_i`$ are respectively the position vector and momentum operator of the $`i`$-th electron with charge $`e`$. $`\epsilon _0`$ is the vacuum permittivity and $`\epsilon _r`$ is the static dielectric constant of the host semiconductor. In addition, there also exist the spin interaction term with the magnetic field $`_{spin}=g^{}\mu _B\left(\stackrel{}{S}_1+\stackrel{}{S}_2\right)\stackrel{}{B},`$ where $`g^{}`$ is the Landé factor, and $`\mu _B`$ is the Bohr magneton. Due to the small $`g^{}`$ in semiconductor, in general, $`_{spin}`$ is very small and can be ingored safely (i.e., for $`g^{}=0.44`$ in GaAs materials and $`B=10`$ T, the typical value of $`_{spin}`$ is $`0.25`$ meV). It is worthy to note that in the limit case $`R_0=0`$, the nanoring simply reduces to a parabolic quantum dot, which can be solved trivially by separating the Hamiltonian into center-of-mass and relative-motion terms. The appearance of $`R_0`$ breaks this separability, and it becomes more complex.
We apply the exact diagonalization method by constructing the basis with the single particle wavefunctions of Hamiltonian
$$_s=\frac{1}{2m_e^{}}\left(\stackrel{}{p}+e\stackrel{}{A}\left(\stackrel{}{r}\right)\right)^2+U(\stackrel{}{r}).$$
(2)
These wavefunctions labelled by the radial quantum number $`n`$ and orbital angular-momentum quantum number $`m`$ have the form
$$\psi _{nm}\left(\stackrel{}{r}\right)=R_{nm}(r)\mathrm{exp}\left(im\phi \right)n=0,1,2\mathrm{},m=0,\pm 1,\mathrm{},$$
(3)
where the radial part $`R_{nm}`$ will be solved exactly by using the series expansion method.
## III Formula and calculation methods
### A Series solution
For the sake of convenience, we use the effective atomic units, in which the effective Rydberg $`R_y^{}=\frac{m_e^{}e^4}{2\mathrm{}^2\left(4\pi \epsilon _0\epsilon _r\right)^2}`$ and the effective Bohr radius $`a^{}=\frac{4\pi \epsilon _0\epsilon _r\mathrm{}^2}{m_e^{}e^2}`$ are taken to be the energy and length units, respectively. Then, the Hamiltonian (2) has the form
$$_s=\stackrel{}{\mathrm{}}^2+\frac{1}{4}\gamma _b^2r^2+\frac{1}{4}\gamma _d^2\left(rR_0\right)^2+\gamma _b\widehat{L}_z,$$
(4)
where the magnetic field $`\gamma _b`$ is measured in the unit $`\frac{\mathrm{}\omega _c}{2R_y^{}}`$ with cyclotron frequency $`\omega _c=\frac{eB}{m_e^{}}`$, $`\gamma _b\widehat{L}_z`$ is the Zeeman term, and $`\gamma _d^{1/2}=\left(\frac{R_y^{}}{\mathrm{}\omega _0}\right)^{1/2}`$ is related to confinement region of the electrons. It is interesting to note how large the units of semiconductor materials are. For GaAs materials, for example, $`R_y^{}=5.8`$ meV, $`a^{}=10`$ nm, and $`\gamma _b=1`$ corresponds to $`B=6.75`$ T.
Now, we have to solve the Schrődinger-like equation
$$_s\left[R_{nm}(r)\mathrm{exp}\left(im\phi \right)\right]=E_{nm}\left[R_{nm}(r)\mathrm{exp}\left(im\phi \right)\right]$$
(5)
to obtain the energy $`E_{nm}`$ and radial part of wavefunction $`R_{nm}(r)`$. It is easy to find the equation satisfied by the function $`R_{nm}\left(r\right)`$
$$\left\{\frac{d^2}{dr^2}+\frac{1}{r}\frac{d}{dr}+\left[\left(E_{nm}m\gamma _b\right)\frac{m^2}{r^2}\frac{1}{4}\gamma _b^2r^2\frac{1}{4}\gamma _d^2\left(rR_0\right)^2\right]\right\}R_{nm}\left(r\right)=0.$$
(6)
We are prevented from analytically exact solutions of the eigenvalue problem because Eq. (6) with suitable boundary conditions is beyond the analytical problem of confluent hypergeometric equations. However, we can use the method of series expansion to obtain the series forms in different regions of Eq. (6) and the exact values of $`E_{nm}`$.
It should be noted that $`r=0`$ and $`r=+\mathrm{}`$ are respectively the regular and irregular singularity points of Eq. (6). So it is natural to divide the whole region $`[0,+\mathrm{})`$ into three parts, $`[0,+\mathrm{})=[0,r_0)[r_0,r_{\mathrm{}})[r_{\mathrm{}},+\mathrm{})`$, where $`r_0`$ and $`r_{\mathrm{}}`$ are the two dividing points. With suitable adjustment of $`r_0`$ and $`r_{\mathrm{}}`$, very high numerical precision can be archived. In all the three kinds of regions, the function $`R_{nm}\left(r\right)`$ is found to be the form:
$$R_{nm}\left(r\right)=\{\begin{array}{cc}Ar^l\underset{n=0}{\overset{\mathrm{}}{}}a_nr^n& 0r<r_0,\\ C_i\underset{n=0}{\overset{\mathrm{}}{}}c_{in}\left(rR_i\right)^n+D_i\underset{n=0}{\overset{\mathrm{}}{}}d_{in}\left(rR_i\right)^n& R_ir<R_{i+1},\\ B\mathrm{exp}\left(\frac{1}{4}\gamma r^2+\frac{\gamma _d^2}{2\gamma }R_0r\right)r^s\underset{n=0}{\overset{\mathrm{}}{}}b_nr^n& r_{\mathrm{}}r<+\mathrm{},\end{array}$$
(7)
where $`\gamma =\sqrt{\gamma _b^2+\gamma _d^2}`$, $`l=\left|m\right|`$ and $`s=\frac{\left(E_{nm}m\gamma _b\right)}{\gamma }1\frac{\gamma _b^2\gamma _d^2}{4\gamma ^3}R_0^2`$. In order to improve precision, we further divide the region $`[r_0,r_{\mathrm{}})`$ into $`N`$ pieces, denoted by $`R_i`$ ($`i=1,\mathrm{},N`$), here $`R_1=r_0`$ and $`R_N=r_{\mathrm{}}`$. $`A`$, $`C_i`$, $`D_i`$ and $`B`$ are constants, $`a_n`$, $`c_{in}`$, $`d_{in}`$ and $`b_n`$ are expanding coefficients and can be determined by the recurrence relation coming from Eq. (6), the initial values of coefficients are chosen to be $`a_0=1`$, $`c_{i0}=d_{i1}=1`$, $`c_{i1}=d_{i0}=0`$ and $`b_0=1`$.
Using the matching conditions at $`r=R_i`$ ($`i=1,\mathrm{},N`$), and the $`2\times 2`$ transfer matrices, we can deduce the equation for eigenenergies $`E_{nm}`$ easily. Then the constant $`A`$, $`C_i`$, $`D_i`$ and $`B`$ can be evaluated by normalization condition, and $`R_{nm}\left(r\right)`$ is obtained numerically.
To close this subsection, it is interesting to point out that the method mentioned above is very suitable for numerical calculations, and can be modified to handle various differential equations similar to Eq. (6).
### B Exact diagonalization
Once the single particle wavefunctions are obtained, we go ahead to construct the basis of the two-electron wavefunctions. It is obvious that the total $`z`$ component of the angular momentum operator of two electrons, $`\widehat{L}_{z,total}=\widehat{L}_{z1}+\widehat{L}_{z2}`$, is a constant of motion, i.e., $`[,\widehat{L}_{z,total}]=0`$ is valid, from which the rotational invariance of the problem follows. Thus, one can work in one subspace labelled by a good quantum number $`L`$ instead of the whole huge Hilbert space. We denote the corresponding two-electron wavefunction by $`\psi _L(\stackrel{}{r}_1\sigma _1;\stackrel{}{r}_2\sigma _2)`$ which depends on the spatial coordinates $`\left\{\stackrel{}{r}_i\right\}`$ and the spin coordinates $`\left\{\sigma _i\right\}`$ ($`i=1,2`$). Because the Hamiltonian of Eq. (1) does not depend on the spin operator, in our two electrons case, the wavefunction can separate into the orbital part $`\psi _L(\stackrel{}{r}_1,\stackrel{}{r}_2)`$ and the spin part $`\chi _S(\sigma _1,\sigma _2)`$:
$$\psi _{LS}(\stackrel{}{r}_1\sigma _1;\stackrel{}{r}_2\sigma _2)=\psi _L^{\lambda _1\lambda _2}(\stackrel{}{r}_1,\stackrel{}{r}_2)\chi _S(\sigma _1,\sigma _2)$$
(8)
with
$$\psi _L^{\lambda _1\lambda _2}(\stackrel{}{r}_1,\stackrel{}{r}_2)=c\left[\psi _{n_1m_1}\left(\stackrel{}{r}_1\right)\psi _{n_2m_2}\left(\stackrel{}{r}_2\right)+()^S\psi _{n_2m_2}\left(\stackrel{}{r}_1\right)\psi _{n_1m_1}\left(\stackrel{}{r}_2\right)\right],$$
(9)
where the normalized wavefunction $`\psi _{LS}(\stackrel{}{r}_1\sigma _1;\stackrel{}{r}_2\sigma _2)`$ is further labelled by the total spin number $`S=0`$ or $`1`$, corresponding to the singlet and triplet state. $`\lambda _i`$ $`\left(i=1,2\right)`$ stands for the quantum number pair $`\left(n_im_i\right).`$ $`L=m_1+m_2`$. $`\psi _{nm}\left(\stackrel{}{r}\right)`$ is the single particle wavefunction. $`c=\sqrt{\frac{1}{2}\text{ }}`$ as $`\left(n_1m_1\right)\left(n_2m_2\right),`$ and $`c=\frac{1}{2}`$ as $`\left(n_1m_1\right)=\left(n_2m_2\right).`$ Obviously, the wavefunction constructed above satisfies the antisymmetric condition,
$$𝒫_{12}\psi _{LS}(\stackrel{}{r}_1\sigma _1;\stackrel{}{r}_2\sigma _2)=\psi _{LS}(\stackrel{}{r}_1\sigma _1;\stackrel{}{r}_2\sigma _2),$$
(10)
where $`𝒫_{12}`$ is the permutation operator.
In the next step, we diagonalize the Hamiltonian (1) numerically in a restricted configuration space. The space is constructed by choosing the wavefunctions having the form (8) in the lowest $`f`$ levels. The secular equation of finite degree $`f`$ is given by
$$det\left(E_i^{(0)}E\right)\delta _{ij}+\mathrm{\Delta }_{ij}=0,i,j=1,\mathrm{},f,$$
(11)
where $`E_i^{(0)}=E_{n_1^im_1^i}+E_{n_2^im_2^i}`$ is the single particle energy, $`i=(n_1^im_1^i;n_2^im_2^i)`$ and $`j=(n_1^jm_1^j;n_2^jm_2^j)`$ represent the quantum number levels. $`\mathrm{\Delta }_{ij}`$ is the matrix element of electron-electron interaction in unit of $`R_y^{}`$
$$\mathrm{\Delta }_{ij}=i\left|\frac{2}{\left|\stackrel{}{r}_1\stackrel{}{r}_2\right|}\right|j=\mathrm{\Delta }_{ij}^c+()^S\mathrm{\Delta }_{ij}^e,$$
(12)
where
$$\mathrm{\Delta }_{ij}^c=𝑑\stackrel{}{r}_1𝑑\stackrel{}{r}_2\psi _{n_1^im_1^i}^{}\left(\stackrel{}{r}_1\right)\psi _{n_2^im_2^i}^{}\left(\stackrel{}{r}_2\right)\frac{2}{\left|\stackrel{}{r}_1\stackrel{}{r}_2\right|}\psi _{n_1^jm_1^j}\left(\stackrel{}{r}_1\right)\psi _{n_2^jm_2^j}\left(\stackrel{}{r}_2\right),$$
(13)
and
$$\mathrm{\Delta }_{ij}^e=𝑑\stackrel{}{r}_1𝑑\stackrel{}{r}_2\psi _{n_1^im_1^i}^{}\left(\stackrel{}{r}_1\right)\psi _{n_2^im_2^i}^{}\left(\stackrel{}{r}_2\right)\frac{2}{\left|\stackrel{}{r}_1\stackrel{}{r}_2\right|}\psi _{n_2^jm_2^j}\left(\stackrel{}{r}_1\right)\psi _{n_1^jm_1^j}\left(\stackrel{}{r}_2\right).$$
(14)
$`\mathrm{\Delta }_{ij}^c`$, $`\mathrm{\Delta }_{ij}^e`$ can be computed numerically from the series solution of $`\psi _{nm}\left(\stackrel{}{r}\right)`$.
By diagonalizing the secular equation (11) in each subspace $`(S,L)`$, we obtain the $`m`$-th energy level $`E_m`$ and the corresponding two-electron wave-function:
$$\mathrm{\Phi }_m(\stackrel{}{r}_1\sigma _1;\stackrel{}{r}_2\sigma _2)=\underset{\lambda _1\lambda _2}{}A_{\lambda _1\lambda _2}^m\psi _L^{\lambda _1\lambda _2}(\stackrel{}{r}_1,\stackrel{}{r}_2)\chi _S(\sigma _1,\sigma _2).$$
(15)
### C Optical absorption
In the electronic dipole approximation, the absorption coefficient is given by
$$\alpha \left(\omega \right)=c\omega \underset{fi}{}\left|f\left|\stackrel{}{e}\stackrel{}{d}\right|i\right|^2\delta \left(\omega \omega _{fi}\right)\left(f_i^{(0)}f_f^{(0)}\right),$$
(16)
where $`\stackrel{}{e}`$ is the complex polarization vector of the spatially constant external electronic field, $`\stackrel{}{d}=\left(\stackrel{}{r}_1+\stackrel{}{r}_2\right)`$ is the electronic dipole operator of two electrons, and $`f^{(0)}`$ is the equilibrium Fermi distribution. The summation is over all the two-electron eigenstates, $`\omega _{fi}`$ is proportional to the energy difference between the initial state $`i`$ and final state $`f:`$ $`\omega _{fi}=\omega _f\omega _i=\left(\frac{E_fE_i}{\mathrm{}}\right).`$ $`c`$ is a constant factor. Restricting ourselves to zero temperature $`T=0`$K, Eq. (16) reduces to
$$\alpha \left(\omega \right)=c\omega \underset{f}{}\left|f\left|\stackrel{}{e}\stackrel{}{d}\right|0\right|^2\delta \left(\omega \omega _{f0}\right),$$
(17)
where $`0`$ and $`f`$ represent the ground state and excited state respectively.
For circularly polarized light we have $`\stackrel{}{e}=\sqrt{\frac{1}{2}}(1,\pm i),`$ from which it follows $`\stackrel{}{e}\stackrel{}{d}=\sqrt{\frac{1}{2}}\left(r_1e^{\pm i\phi _1}+\left(12\right)\right)`$. Using Eq. (15), we obtain
$$f\left|\stackrel{}{e}\stackrel{}{d}\right|0=\sqrt{\frac{1}{2}}\underset{\lambda _1^0\lambda _2^0}{}\underset{\lambda _1^f\lambda _2^f}{}A_{\lambda _1^f\lambda _2^f}^fA_{\lambda _1^0\lambda _2^0}^0\left(\left(d_{\lambda _1^0\lambda _1^f}^\pm \delta _{\lambda _2^0\lambda _2^f}+()^Sd_{\lambda _1^0\lambda _2^f}^\pm \delta _{\lambda _2^0\lambda _1^f}\right)+\left(12\right)\right),$$
(18)
where the single particle matrix elements $`d_{\lambda \lambda ^{}}^\pm `$ are defined by
$$d_{\lambda \lambda ^{}}^\pm =\lambda ^{}\left|r\mathrm{exp}\left(\pm i\phi \right)\right|\lambda =\delta _{m\pm 1,m^{}}\underset{0}{\overset{\mathrm{}}{}}r^2R_\lambda ^{}\left(r\right)R_\lambda \left(r\right)𝑑r$$
(19)
and $`|\lambda `$ represents the single particle wave-function $`\psi _{nm}\left(\stackrel{}{r}\right)`$. It is obvious that the absorption coefficient satisfies the dipole selection rule $`\mathrm{\Delta }L=\pm 1.`$ Substituting Eq. (18) into Eq. (17) and taking
$$\delta \left(\omega \omega _{f0}\right)=\frac{\mathrm{\Gamma }/\pi }{\left(\omega \omega _{f0}\right)^2+\mathrm{\Gamma }^2},$$
(20)
where $`\mathrm{\Gamma }`$ is a phenomenological broadening parameter, we arrive at
$`\alpha ^\pm \left(\omega \right)`$ $`=`$ $`c{\displaystyle \frac{\omega }{2}}{\displaystyle \underset{f}{}}|{\displaystyle \underset{\lambda _1^0\lambda _2^0}{}}{\displaystyle \underset{\lambda _1^f\lambda _2^f}{}}A_{\lambda _1^f\lambda _2^f}^fA_{\lambda _1^0\lambda _2^0}^0((d_{\lambda _1^0\lambda _1^f}^\pm \delta _{\lambda _2^0\lambda _2^f}+()^Sd_{\lambda _1^0\lambda _2^f}^\pm \delta _{\lambda _2^0\lambda _1^f})+(12))|^2\times `$ (22)
$`{\displaystyle \frac{\mathrm{\Gamma }/\pi }{\left(\omega \omega _{f0}\right)^2+\mathrm{\Gamma }^2}},`$
where $`\pm `$ corresponds to right and left circularly polarized lights, respectively.
To check the numerical accuracy of the calculations we have used the $`f`$-sum rules for the dipole operators, which can be expressd in terms of ground-state quantities:
$$\underset{0}{\overset{\mathrm{}}{}}\left(\alpha ^+\left(\omega \right)+\alpha ^{}\left(\omega \right)\right)𝑑\omega =0\left|[\left(\stackrel{}{e}\stackrel{}{d}\right)^+,[,\left(\stackrel{}{e}\stackrel{}{d}\right)]]\right|0=N$$
(23)
in the effective atomic units, where $`c`$ is taken to be 1, and $`N=2`$ is the number of electrons.
It is important to point out that in our case, due to non-separability of the center-of-mass and the relative-motion modes, the generalized Kohn theorem, which means that FIR can only be used to excite the center-of-mass modes of electrons parabolically confined in circular quantum dot, will not be held further. Thus we expect our FIR absorption result may reflect an excitation of the relative motion of two electrons.
## IV Results and discussions
To explain the experimental measurements, we have taken the material parameters $`\epsilon _r=12.4`$ and $`m_e^{}=0.067m_e`$ for GaAs in our calculations. The radial confinement strength $`\mathrm{}\omega _0`$ and the ring’s radius $`R_0`$ are chosen to be $`12`$ meV and $`14`$ nm respectively. The corresponding width of ring is about $`15`$ nm, which means that the electrons are confined in a wide ring. Thus in contrast to the rotating Wigner molecule picture in a narrow-width quantum ring, the more pronounced energy spectra and optical properites are expected. For the calculation in each $`(S,L)`$ subspace, we first solve the single particle problem and save several hundreds single particle states, then pick up the suitable single particle states to construct thousands of two-electron states, among which only the lowest $`f`$ energy levels are selected. Here we simply point out that our numerical diagonalization scheme is very efficient and essentially exact in the sense that the accuracy can be improved as desired by increasing $`f`$. For instance, for the ground state in $`(0,0)`$ subspace, the use of 64 basis states allows the precision to be within the relative convergence of $`10^4.`$ On the other hand, by checking the $`f`$-sum rule, we find that the relative error of our FIR calculation is less than 1%.
### A Spin oscillation
The energy levels of two electrons in a nanoring as a function of the magnetic field have been plotted in Fig.1. As mentioned above, only the total spin and total angular momentum are conserved in our model, and then for the sake of clearness, we only plot the lowest level in each $`(S,L)`$ subspace. In order to understand the role of the electron-electron interaction in the two-electron spectra better, we first describe the characteristics of the energy levels of two electrons in the nanoring without interaction. As shown in Fig. 1(a), for small magnetic field, the spectra have the characteristics of disk-like quantum dot, since the electrons are confined in a wide nanoring whose radius is comparable to its effective width. As the magnetic field increases from zero, there are minima for the states with negative angular momentum $`L`$. These minima are caused by the interplay between the Zeeman term and the ring-like confinement potential. Moreover, the level with $`L=0`$ increases monotonically, and changes more dramatically than the others. This leads to a ground state transition from $`L=0`$ to $`L=2`$ around $`B=8`$ T, which also reflects the fact that the nanoring becomes more and more narrow with increasing the magnetic field. In the sufficient high field regime, it is quite safe to speculate that the levels will be the same as the levels in one-dimensional ring under a uniform magnetic field.
The electron-electron interaction can significantly change the characteristics of the spectra described above. As shown in Fig. 1(b), clearly visible is the ground state transitions around $`B=3`$, $`7`$ and $`10`$ T. These transitions are quite different from that mentioned in the previous paragraph: not only the total angular momentum but also the total spin are changed. They present the spin-singlet-spin-triplet oscillation of the ground state in the magnetic field, i.e., $`(0,0)(1,1)(0,2)(1,3)`$ states and so on. This phenomenon is indeed qualitatively similar to that seen in a quantum dot.
For a better understanding of the singlet-triplet oscillation, it is interesting to study the electron-electron interaction energies $`E_r`$, defined by the difference between the energy $`E`$ with interaction and $`E^{(0)}`$ without interaction. In Fig. 2, the $`E_r`$ for different states are plotted as a function of the magnetic field. It is readily seen that the $`E_r`$ increase with the magnetic field and the ordering is as follows: $`E_r(0,1)>E_r(0,0)>E_r(0,3)>E_r(1,2)>E_r(1,0)>E_r(1,1)>E_r(0,2)>E_r(1,3)\mathrm{}`$. Compared with the corresponding results in quantum dot, we find that the ordering depends significantly on the form of the confinement potential, e.g., the ordering $`E_r(0,0)>E_r(0,1)`$ in quantum dot is reversed in our case. However, the ordering $`E_r(0,0)>E_r(1,1)>E_r(0,2)>E_r(1,3)`$ is still preserved, thus the trivial crossover around $`B=8`$ T in Fig. 1(a) moves to the position $`B=3`$ T in Fig. 1(b).
We comment on some other effects caused by the interaction. ($`i`$) When the electron-electron interaction is excluded, the states $`(0,L)`$ and $`(1,L)`$ $`\left(L=1,3\right)`$ are degenerate in the whole regime of the magnetic field. This degeneracy can be understood by the fact that those states are constructed with the same single particle states, the only difference between them is spin part which have no influence to energy levels in the absence of interaction. When the electron-electron interaction turns on, the whole energy levels are shifted to high energies due to the repulsive coulomb interaction and the degenerated levels split. ($`ii`$) In the high magnetic field regime, the energy levels appear to be more separated in the presence of interaction, and it indicates that the interaction effect becomes larger in this regime.
### B FIR spectroscopy
In the previous subsection, we have demonstrated that the electron-electron interaction can induce some transitions of the ground state. It may be possible to observe these transitions in the lower energy optical absorption. Now, we present the FIR spectroscopy to elucidate such transitions. Here, the phenomenological broadening parameter $`\mathrm{\Gamma }`$ is assumed to be $`0.5`$ meV.
In Fig. 3, we plot the FIR absorption for the circularly polarized lights as a function of the magnetic field. Compared with the no interaction case, where the right and left circularly polarized (labeled by $`\sigma _\pm `$) resonance energies are roughly given by $`\omega _\pm =`$ $`\frac{1}{2}\left(\sqrt{\omega _c^2+4\omega _0^2}\pm \omega _c\right)`$in the low field (as shown in Fig. 4, $`\sigma _+`$ and $`\sigma _{}`$ are represented by the square and filled circle symbols, respectively), we find some unusual features in the presence of interaction. ($`i`$) As the magnetic field increases, the lowest right circularly polarized $`\sigma _+`$ resonance energy shows discontinuous drops around $`3`$, $`7`$ and $`10`$ T (see Fig. 3a). In contrast, this behavior is only found at $`B8`$ T in the case without the Coulomb interaction (see Fig. 4). Naturally, the discontinuous drops are originated from the spin oscillation of the ground state as mentioned above. Note that they are less reflected in the $`\sigma _{}`$ case. ($`ii`$) The absorption peak around $`20`$ meV splits into three or more subpeaks, especially in the $`\sigma _{}`$ polarization case. The striking behavior can be understood from the mixing of the center-of-mass and the relative-motion modes. Certainly, it has to keep in mind that due to the mixing there are many irregularly spaced and near-degenerated energy levels hybridized from the relative-motion mode in each $`(S,L)`$ subspace. The splitting can be identified as a transition from the ground state to those ”hybrid” states, showing its two-electron characteristic. ($`iii`$) For both circular polarizations, the lowest resonance energy shows slight blue shift due to the Coulomb interaction.
Motivated by the above prominent features, we suggest that the spin transitions of the ground state and the other effects induced by the electron-electron interaction can be observed with circularly polarized light. However, as pointed out by A. Lorke, in theirs experiment, the single particle states are quite accurate basis for the description of the measured FIR resonance. It seems in contrast to ours expectation. We argue here that the controversy comes from the resolution of the experiment. With low resolution, one is difficult to distinguish the details of the FIR spectroscopy discussed above, and only the profile is explored. As shown in Figs. (3a) and (3b), we can indeed observe that the profile of the FIR resonance can be described by the single particle picture. Therefore, we expect that our predictions about the electron-electron interaction may be confirmed by a high resolution experiment in the future.
## V Summary
In conclusion, we have investigated the energy levels and far-infrared spectroscopy of a two-electron nanoring in a magnetic field. Because of the electron-electron interaction as well as the interplay between the magnetic field and ring-like confinement potential, the nanoring exhibits rich electronic structures. An obvious feature induced by the interaction is the intersection between the lower levels. It presents the spin oscillation of the ground state on the magnetic field. i.e., $`(0,0)(1,1)(0,2)(1,3)`$ states and so on. This phenomenon is indeed qualitatively similar to that seen in a quantum dot, suggesting its intrinsic nature of zero-dimensional quantum structures.
The profile of the FIR spectroscopy is roughly captured by the single particle picture as indicated by a recent experiment. However, the ring-like confinement potential doesn’t allow the application of the generalized Kohn theorem. Thus the single particle picture is inadequate for seeing any effect due to electron-electron interaction. We have outlined here two prominent features of the FIR spectroscopy caused by electron-electron interaction: splitting and dis-continuous drops of the resonance energies. We suggest that those two features can be detectable by using the circularly polarized far-infrared light with high experimental resolution.
ACKNOWLEDGMENT
The financial support from NSF-China (Grant No.19974019) and China’s ”973” program is gratefully acknowledged.
Figures Captions
Fig. 1. The energy levels of two electrons in a nanoring are plotted as a function of the magnetic field with $`\mathrm{}\omega _0=12`$ meV and $`R_0=14`$ nm in the (a) absence and (b) presence of electron-electron interaction. The spin singlet and triplet states are labelled by solid and dashed lines, respectively. The quantum number of each state (total spin, total angular momentum) are also indicated. Note that only the lowest energy level of each $`(S,L)`$ subspace is selected.
Fig. 2. The net Coulomb energies are plotted as a function of the magnetic field. The other parameters are the same as in Fig. 1.
Fig. 3. A logarithmic 3D plot of the far-infrared absorption coefficient as a function of the magnetic field for (a) right and (b) left circularly polarized light.
Fig. 4. Far-infrared absorption resonance energies in the absence of electron-electron interaction as a function of the magnetic field. The square and filled circle symbols are corresponding to right and left circularly polarized lights, respectively.
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# Brownian motion in a magnetic field
## 1 <br>Introduction
We address an old-fashioned problem of the Brownian motion of a charged particle in a constant magnetic field. That issue has originated from studies of the diffusion of plasma across a magnetic field , and nowadays, together with a free Brownian motion example, stands for a textbook illustration of how transport and auto-correlation functions should be computed in generic situations governed by the Langevin equation (to a suitable degree of approximation of a kinetic theory, when collisions are stochastically modeled in terms of a random force), cf. but also , .
From a purely pragmatic point of view this white-noise strategy is quite satisfactory. After (formally) evaluating velocity auto-correlation functions, formulas for running and asymptotic diffusion coefficients easily follow. To that end an explicit form of the probability density or transition probability density of the involved stochastic diffusion processes (in velocity space, phase-space or configuration space) is not necessary, cf. , .
To our knowledge, except for the paper (mentioned in as a footnote reference for the purpose of evaluation of the mean square velocity and its mean square displacement at equilibrium), for a Brownian particle in a constant magnetic field no attempt was made in the literature to give a complete characterization of the stochastic process itself, nor pass to the associated macrosopic (hydrodynamical formalism) balance equations. (Cf. , , , , , for a number of reasons why to do that).
Surprisingly enough, in Ref. , the Brownian motion in a magnetic field is described in terms of operator-valued (matrix-valued functions) probability distributions that additionally involve fractional powers of matrices. In consequence, there is no clean path towards a (necessary) relationship with the associated Kramers-Smoluchowski equations (cf. Chap. 6.1 in Ref. ), nor ways to stay in conformity with the standard wisdom about probabilistic procedures valid in case of the free Brownian motion (Ornstein-Uhlenbeck process), cf. , , .
Therefore, we decided to address an issue of the Brownian motion in a magnetic field anew, to unravel its features of a fully-fledged stochastic diffusion process. In particular, we derive transition probability densities governing both the velocity, phase-space and the configuration space processes. Hydrodynamical balance equations and their behaviour in the Smoluchowski regime are discussed as well.
## 2 Velocity-space diffusion process
The standard analysis of the Brownian motion of a free particle employs the Langevin equation $`\frac{d\stackrel{}{u}}{dt}=\beta \stackrel{}{u}+\stackrel{}{A}\left(t\right)`$ where $`\stackrel{}{u}`$ denotes the velocity of the particle and the influence of the surrounding medium on the motion (random acceleration) of the particle is modeled by means of two independent contributions. A systematic part $`\beta \stackrel{}{u}`$ represents a dynamical friction. The remaining fluctuating part $`\stackrel{}{A}\left(t\right)`$ is supposed to display a statistics of the familiar white noise: (i)$`\stackrel{}{A}\left(t\right)`$ is independent of $`\stackrel{}{u}`$, (ii) $`A_i\left(s\right)=0`$ and $`A_i\left(s\right)A_j\left(s^{}\right)=2q\delta _{ij}\delta \left(ss^{}\right)`$ for $`i,j=1,2,3`$, where $`q=\frac{k_BT}{m}\beta `$ is a physical parameter. The well-known Ornstein-Uhlenbeck stochastic process comes out on that conceptual basis, ,, .
The linear friction model can be adopted to the case of diffusion of charged particles in the presence of a constant magnetic field which acts upon particles via the Lorentz force. The Langevin equation for that motion reads:
$$\frac{d\stackrel{}{u}}{dt}=\beta \stackrel{}{u}+\frac{q_e}{mc}\stackrel{}{u}\times \stackrel{}{B}+\stackrel{}{A}\left(t\right)$$
(1)
where $`q_e`$ denotes an electric charge of the particle of mass $`m`$.
Let us assume for simplicity that the constant magnetic field $`\stackrel{}{B}`$ is directed along the z-axis of a Cartesian reference frame: $`\stackrel{}{B}=(0,0,B)`$ and $`B=const`$. In this case Eq. (1) takes the form
$$\frac{d\stackrel{}{u}}{dt}=\mathrm{\Lambda }\stackrel{}{u}+\stackrel{}{A}\left(t\right)$$
(2)
where
$$\mathrm{\Lambda }=\left(\begin{array}{ccc}\beta & \omega _c& 0\\ \omega _c& \beta & 0\\ 0& 0& \beta \end{array}\right)$$
(3)
and $`\omega _c=\frac{q_eB}{mc}`$ denotes the Larmor frequency. Assuming the Langevin equation to be (at least formally) solvable, we can infer a probability density $`P(\stackrel{}{u},t|\stackrel{}{u}_0)`$, $`t>0`$ conditioned by the the initial velocity data choice $`\stackrel{}{u}=\stackrel{}{u}_0`$ at $`t=0`$. Physical circumstances of the problem enforce a demand: (i) $`P(\stackrel{}{u},t|\stackrel{}{u}_0)\delta ^3\left(\stackrel{}{u}\stackrel{}{u}_0\right)`$ as $`t0`$ and (ii) $`P(\stackrel{}{u},t|\stackrel{}{u}_0)\left(\frac{m}{2\pi k_BT}\right)^{\frac{3}{2}}\mathrm{exp}\left(\frac{m|\stackrel{}{u}_0|^2}{2k_BT}\right)`$ as $`t\mathrm{}`$.
A formal solution of Eq. (2) reads:
$$\stackrel{}{u}\left(t\right)e^{\mathrm{\Lambda }t}\stackrel{}{u}_0=_0^te^{\mathrm{\Lambda }\left(ts\right)}\stackrel{}{A}\left(s\right)𝑑s.$$
(4)
By taking into account that
$$e^{\mathrm{\Lambda }t}=e^{\beta t}\left(\begin{array}{ccc}\mathrm{cos}\omega _ct& \mathrm{sin}\omega _ct& 0\\ \mathrm{sin}\omega _ct& \mathrm{cos}\omega _ct& 0\\ 0& 0& 1\end{array}\right)=e^{\beta t}U\left(t\right)$$
(5)
we can rewrite (4) as follows
$$\stackrel{}{u}\left(t\right)e^{\beta t}U\left(t\right)\stackrel{}{u}_0=_0^te^{\beta \left(ts\right)}U\left(ts\right)\stackrel{}{A}\left(s\right)𝑑s.$$
(6)
Statistical properties of $`\stackrel{}{u}\left(t\right)e^{\mathrm{\Lambda }t}\stackrel{}{u}_0`$ are identical with those of $`\stackrel{}{A}\left(s\right)ds`$. In consequence, the problem of deducing a probability density $`P(\stackrel{}{u},t|\stackrel{}{u}_0)`$ is equivalent to deriving the probability distribution of the random vector
$$\stackrel{}{S}=_0^t\psi \left(s\right)\stackrel{}{A}\left(s\right)𝑑s$$
(7)
where $`\psi \left(s\right)=e^{\mathrm{\Lambda }\left(ts\right)}=e^{\beta \left(ts\right)}U\left(ts\right)`$.
The white noise term $`\stackrel{}{A}\left(s\right)`$ in view of the integration with respect to time is amenable to a more rigorous analysis that invokes the Wiener process increments and their statistics, . Let us divide the time integration interval into a large number of small subintervals $`\mathrm{\Delta }t`$. We adjust them suitably to assure that effectively $`\psi \left(s\right)`$ is constant on each subinterval $`(j\mathrm{\Delta }t,\left(j+1\right)\mathrm{\Delta }t)`$ and equal $`\psi \left(j\mathrm{\Delta }t\right)`$. As a result we obtain the expression
$$\stackrel{}{S}=\underset{j=0}{\overset{N1}{}}\psi \left(j\mathrm{\Delta }t\right)_{j\mathrm{\Delta }t}^{\left(j+1\right)\mathrm{\Delta }t}\stackrel{}{A}\left(s\right)𝑑s.$$
(8)
Here $`\stackrel{}{B}\left(\mathrm{\Delta }t\right)=_{j\mathrm{\Delta }t}^{\left(j+1\right)\mathrm{\Delta }t}\stackrel{}{A}\left(s\right)𝑑s`$ stands for the above-mentioned Wiener process increment. Physically, $`\stackrel{}{B}\left(\mathrm{\Delta }t\right)`$ represents the net acceleration which a Brownian particle may suffer (in fact accumulates) during an interval of time $`\mathrm{\Delta }t`$.
Equation (8) becomes
$$\stackrel{}{S}=\underset{j=0}{\overset{N1}{}}\psi \left(j\mathrm{\Delta }t\right)\stackrel{}{B}\left(\mathrm{\Delta }t\right)=\underset{j=0}{\overset{N1}{}}\stackrel{}{s}_j$$
(9)
where we introduce $`\stackrel{}{s}_j=\psi \left(j\mathrm{\Delta }t\right)\stackrel{}{B}\left(\mathrm{\Delta }t\right)=\psi _j\stackrel{}{B}\left(\mathrm{\Delta }t\right)`$.
The Wiener process argument , , allows us to infer the probability distribution of $`\stackrel{}{s}_j`$. It is enough to employ the fact that the distribution of $`\stackrel{}{B}\left(\mathrm{\Delta }t\right)`$ is Gaussian with mean zero and variance $`q=\frac{k_BT}{m}\beta `$. Then
$$w\left[\stackrel{}{B}\left(\mathrm{\Delta }t\right)\right]=\left(\frac{1}{4\pi q\mathrm{\Delta }t}\right)^{\frac{3}{2}}\mathrm{exp}\left(\frac{\left|\stackrel{}{B}\left(\mathrm{\Delta }t\right)\right|^2}{4q\mathrm{\Delta }t}\right)$$
(10)
and in view of $`\stackrel{}{s}_j=\psi _j\stackrel{}{B}\left(\mathrm{\Delta }t\right)`$ by performing the change of variables in (10) we get
$$\stackrel{~}{w}\left[\stackrel{}{s}_j\right]=det\left[\psi _j^1\right]w\left[\psi _j^1\stackrel{}{s}_j\right]=\frac{1}{det\psi _j}w\left[\psi _j^1\stackrel{}{s}_j\right].$$
(11)
Since $`det\psi \left(s\right)=e^{3\beta \left(ts\right)}`$ and $`\psi ^1\left(s\right)=U\left[\left(ts\right)\right]e^{\beta \left(ts\right)}`$ we obtain
$$\stackrel{~}{w}\left[\stackrel{}{s}_j\right]=\left(\frac{1}{4\pi q\mathrm{\Delta }t}\right)^{\frac{3}{2}}\frac{1}{e^{3\beta \left(tj\mathrm{\Delta }t\right)}}\mathrm{exp}\left(\frac{\left|e^{\beta \left(tj\mathrm{\Delta }t\right)}U\left[\left(tj\mathrm{\Delta }t\right)\right]\stackrel{}{s}_j\right|^2}{4q\mathrm{\Delta }t}\right)$$
(12)
and finally
$$\stackrel{~}{w}\left[\stackrel{}{s}_j\right]=\left(\frac{1}{4\pi q\mathrm{\Delta }t}\frac{1}{e^{2\beta \left(tj\mathrm{\Delta }t\right)}}\right)^{\frac{3}{2}}\mathrm{exp}\left(\frac{\left|\stackrel{}{s}_j\right|^2}{4q\mathrm{\Delta }te^{2\beta \left(tj\mathrm{\Delta }t\right)}}\right).$$
(13)
Clearly, $`\stackrel{}{s}_j`$ are mutually independent random variables whose distribution is Gaussian with mean zero and variance $`\sigma _j^2=2q\mathrm{\Delta }te^{2\beta \left(tj\mathrm{\Delta }t\right)}`$. Hence, the probability distribution of $`\stackrel{}{S}=_{j=0}^{N1}\stackrel{}{s}_j`$ is again Gaussian with mean zero. Its variance equals the sum of variances of $`\stackrel{}{s}_j`$ i.e. $`\sigma ^2=_j\sigma _j^2=2q_j\mathrm{\Delta }te^{2\beta \left(tj\mathrm{\Delta }t\right)}`$.
Taking the limit $`N\mathrm{}`$ $`(\mathrm{\Delta }t0)`$ we arrive at
$$\sigma ^2=2q_0^t𝑑se^{2\beta \left(ts\right)}=\frac{k_BT}{m}\left(1e^{2\beta t}\right).$$
(14)
Because of $`\stackrel{}{S}=\stackrel{}{u}\left(t\right)e^{\mathrm{\Lambda }t}\stackrel{}{u}_0`$ the transition probability density of the Brownian particle velocity, conditioned by the initial data $`\stackrel{}{u}_0`$ at $`t_0=0`$ reads
$$P(\stackrel{}{u},t|\stackrel{}{u}_0)=\left(\frac{1}{2\pi \frac{k_BT}{m}\left(1e^{2\beta t}\right)}\right)^{\frac{3}{2}}\mathrm{exp}\left(\frac{\left|\stackrel{}{u}e^{\mathrm{\Lambda }t}\stackrel{}{u}_0\right|^2}{2\frac{k_BT}{m}\left(1e^{2\beta t}\right)}\right).$$
(15)
The process is Markovian and time-homogeneous, hence the above formula can be trivially extended to encompass the case of arbitrary $`t_00`$ : $`P(\stackrel{}{u},t|\stackrel{}{u}_0,t_0)`$ arises by substituting everywhere $`tt_0`$ instead of $`t`$.
Physical arguments (cf. demand (ii) preceding Eq. (4)) refer to an asymptotic probability distribution (invariant measure density) $`P(u)`$ of the random variable $`\stackrel{}{u}`$ in the Maxwell-Boltzmann form
$$P\left(\stackrel{}{u}\right)=\left(\frac{m}{2\pi k_BT}\right)^{\frac{3}{2}}\mathrm{exp}\left(\frac{m\left|\stackrel{}{u}\right|^2}{2k_BT}\right).$$
(16)
This time-independent probability density together with the time-homogeneous transition density (15) uniquely determine a stationary Markovian stochastic process for which we can evaluate various mean values.
Expectation values of velocity components vanish: $`u_i\left(t\right)=_{\mathrm{}}^{\mathrm{}}u_iP\left(\stackrel{}{u}\right)𝑑\stackrel{}{u}=0`$ for $`i=1,2,3`$. The matrix of the second moments (velocity auto-correlation functions) reads
$$u_i\left(t\right)u_j\left(t_0\right)=_{\mathrm{}}^{\mathrm{}}u_iu_j^0P(\stackrel{}{u},t;\stackrel{}{u}_0,t_0)𝑑\stackrel{}{u}𝑑\stackrel{}{u}_0$$
(17)
where $`i,j=1,2,3`$ and in view of $`P(\stackrel{}{u},t;\stackrel{}{u}_0,t_0)=P(\stackrel{}{u},t|\stackrel{}{u}_0,t_0)P\left(\stackrel{}{u}_0\right)`$ we arrive at the compact expression
$$\frac{k_BT}{m}e^{\mathrm{\Lambda }\left|tt_0\right|}=\frac{k_BT}{m}e^{\beta \left|tt_0\right|}\left(\begin{array}{ccc}\mathrm{cos}\omega _c\left|tt_0\right|& \mathrm{sin}\omega _c\left|tt_0\right|& 0\\ \mathrm{sin}\omega _c\left|tt_0\right|& \mathrm{cos}\omega _c\left|tt_0\right|& 0\\ 0& 0& 1\end{array}\right).$$
(18)
In particular, the auto-correlation function (second moment) of the $`x`$-component of velocity equals
$$u_1\left(t\right)u_1\left(t_0\right)=\frac{k_BT}{m}e^{\beta \left|tt_0\right|}\mathrm{cos}\omega _c\left|tt_0\right|$$
(19)
in agreement with white noise calculations of Refs. and , cf. Chap.11, formula (11.25). In particular, the so-called running diffusion coefficient arises via straightforward integration of the function $`R_{11}(\tau )=<u_1(t)u_1(t_0)>`$ where $`\tau =tt_0>0`$:
$$D_1(t)=_0^t<u_1(0)u_1(\tau )>𝑑\tau =\frac{k_BT}{m}\frac{\beta +[\omega _csin(\omega _ct)\beta cos(\omega _ct)]exp(\beta t)}{\beta ^2+\omega _{c}^{}{}_{}{}^{2}}$$
(20)
with an obvious asymptotics (the same for $`D_2(t)`$): $`D_B=lim_t\mathrm{}D_1(t)=\frac{k_BT}{m}\frac{\beta }{\beta ^2+\omega _{c}^{}{}_{}{}^{2}}`$ and the large friction ($`\omega _c`$ fixed and bounded) version $`D=\frac{k_BT}{m\beta }`$.
## 3 Spatial process
The cylindrical symmetry of the problem allows us to consider separately processes running on the $`XY`$ plane and along the $`Z`$-axis (where the free Brownian motion takes place). We shall confine further attention to the two-dimensional $`XY`$-plane problem. Henceforth, each vector will carry two components which correspond to the $`x`$ and $`y`$ coordinates respectively. We will directly refer to the vector and matrix quantities introduced in the previous section, but while keeping the same notation, we shall simply disregard their $`z`$-coordinate contributions.
We define the spatial displacement $`\stackrel{}{r}`$ of the Brownian particle as folows
$$\stackrel{}{r}\stackrel{}{r}_0=_0^t\stackrel{}{u}\left(\eta \right)𝑑n$$
(21)
where $`\stackrel{}{u}\left(t\right)`$ is given by Eq. (2) (except for disregarding the third coordinate).
Our aim is to derive the probability distribution of $`\stackrel{}{r}`$ at time $`t`$ provided that the particle position and velocity were equal $`\stackrel{}{r}_0`$ and $`\stackrel{}{u}_0`$ respectively, at time $`t_0=0`$.
To that end we shall mimic procedures of the previous section. In view of:
$$\stackrel{}{r}\stackrel{}{r}_0_0^te^{\mathrm{\Lambda }\eta }\stackrel{}{u}_0=_0^t𝑑\eta _0^\eta 𝑑se^{\mathrm{\Lambda }\left(\eta s\right)}\stackrel{}{A}\left(s\right)$$
(22)
we have
$$\stackrel{}{r}\stackrel{}{r}_0\mathrm{\Lambda }^1\left(1e^{\mathrm{\Lambda }t}\right)\stackrel{}{u}_0=_0^t\mathrm{\Lambda }^1\left(1e^{\mathrm{\Lambda }\left(st\right)}\right)\stackrel{}{A}\left(s\right)𝑑s$$
(23)
where
$$\mathrm{\Lambda }^1=\frac{1}{\beta ^2+\omega _c^2}\left(\begin{array}{cc}\beta & \omega _c\\ \omega _c& \beta \end{array}\right)$$
(24)
is the inverse of the matrix $`\mathrm{\Lambda }`$ (regarded as a rank two sub-matrix of that originally introduced in Eq. (3)). Let us define two auxiliary matrices
$`\mathrm{\Omega }`$ $``$ $`\mathrm{\Lambda }^1\left(1e^{\mathrm{\Lambda }t}\right)`$ (25)
$`\varphi \left(s\right)`$ $``$ $`\mathrm{\Lambda }^1\left(1e^{\mathrm{\Lambda }\left(st\right)}\right).`$
Because of:
$$e^{\mathrm{\Lambda }t}=\mathrm{exp}\left\{t\left(\begin{array}{cc}\beta & \omega _c\\ \omega _c& \beta \end{array}\right)\right\}=e^{\beta t}\left(\begin{array}{cc}\mathrm{cos}\omega _ct& \mathrm{sin}\omega _ct\\ \mathrm{sin}\omega _ct& \mathrm{cos}\omega _ct\end{array}\right)=e^{\beta t}U\left(t\right)$$
(26)
we can represent matrices $`\mathrm{\Omega }`$, $`\varphi \left(s\right)`$ in more detailed form. We have:
$$\mathrm{\Omega }=\frac{1}{\beta ^2+\omega _c^2}\left\{\left(\begin{array}{cc}\beta & \omega _c\\ \omega _c& \beta \end{array}\right)e^{\beta t}\left(\begin{array}{cc}\beta & \omega _c\\ \omega _c& \beta \end{array}\right)\left(\begin{array}{cc}\mathrm{cos}\omega _ct& \mathrm{sin}\omega _ct\\ \mathrm{sin}\omega _ct& \mathrm{cos}\omega _ct\end{array}\right)\right\}$$
(27)
and
$$\varphi \left(s\right)=\mathrm{\Lambda }^1\left(1e^{\beta \left(ts\right)}U\left(ts\right)\right)=$$
(28)
$$\frac{1}{\beta ^2+\omega _c^2}\left(\begin{array}{cc}\beta & \omega _c\\ \omega _c& \beta \end{array}\right)\left(\begin{array}{cc}1e^{\beta \left(st\right)}\mathrm{cos}\omega _c\left(st\right)& e^{\beta \left(st\right)}\mathrm{sin}\omega _c\left(st\right)\\ e^{\beta \left(st\right)}\mathrm{sin}\omega _c\left(st\right)& 1e^{\beta \left(st\right)}\mathrm{cos}\omega _c\left(st\right)\end{array}\right).$$
Next steps imitate procedures of the previous section. Thus, we seek for the probability distribution of the random (planar) vector $`\stackrel{}{R}=_0^t\varphi \left(s\right)\stackrel{}{A}\left(s\right)𝑑s`$ where $`\stackrel{}{R}=\stackrel{}{r}\stackrel{}{r}_0\mathrm{\Omega }\stackrel{}{u}_0`$.
Dividing the time interval $`(0,t)`$ into small subintervals to assure that $`\varphi \left(s\right)`$ can be regarded constant over the time span $`(j\mathrm{\Delta }t,\left(j+1\right)\mathrm{\Delta }t)`$ and equal $`\varphi \left(j\mathrm{\Delta }t\right)`$, we obtain
$$\stackrel{}{R}=\underset{j=0}{\overset{N1}{}}\varphi \left(j\mathrm{\Delta }t\right)_{j\mathrm{\Delta }t}^{\left(j+1\right)\mathrm{\Delta }t}\stackrel{}{A}\left(s\right)𝑑s=\underset{j=0}{\overset{N1}{}}\varphi \left(j\mathrm{\Delta }t\right)\stackrel{}{B}\left(\mathrm{\Delta }t\right)=\underset{j=0}{\overset{N1}{}}\stackrel{}{r}_j$$
(29)
where $`\stackrel{}{r}_j=\varphi \left(j\mathrm{\Delta }t\right)\stackrel{}{B}\left(\mathrm{\Delta }t\right)=\varphi _j\stackrel{}{B}\left(\mathrm{\Delta }t\right)`$.
By invoking the probability distribution (10) we perform an appropriate change of variables: $`\stackrel{}{r}_j=\varphi _j\stackrel{}{B}\left(\mathrm{\Delta }t\right)`$ to yield a probability distribution of $`\stackrel{}{r}_j`$
$$\stackrel{~}{w}\left[\stackrel{}{r}_j\right]=det\left[\varphi _j^1\right]w\left[\varphi _j^1\stackrel{}{r}_j\right]=\frac{1}{det\varphi _j}w\left[\varphi _j^1\stackrel{}{r}_j\right].$$
(30)
Presently (not to be confused with previous steps (11)-(15)) we have
$$det\varphi \left(s\right)=\frac{1}{\beta ^2+\omega _c^2}\left(1+e^{2\beta \left(st\right)}2e^{\beta \left(st\right)}\mathrm{cos}\omega _c\left(st\right)\right)$$
(31)
and
$$\varphi ^1\left(s\right)=\frac{1}{1+e^{2\beta \left(st\right)}2e^{\beta \left(st\right)}\mathrm{cos}\omega _c\left(st\right)}\left[1e^{\beta \left(st\right)}U\left(\left(st\right)\right)\right]\mathrm{\Lambda }.$$
(32)
So, the inverse of the matrix $`\varphi _j`$ has the form:
$$\varphi _j^1=\frac{\stackrel{~}{A}_j}{\gamma _j}$$
(33)
where
$$\stackrel{~}{A}_j=\left(\begin{array}{cc}1e^{\beta \left(j\mathrm{\Delta }tt\right)}\mathrm{cos}\omega _c\left(j\mathrm{\Delta }tt\right)& e^{\beta \left(j\mathrm{\Delta }tt\right)}\mathrm{sin}\omega _c\left(j\mathrm{\Delta }tt\right)\\ e^{\beta \left(j\mathrm{\Delta }tt\right)}\mathrm{sin}\omega _c\left(j\mathrm{\Delta }tt\right)& 1e^{\beta \left(j\mathrm{\Delta }tt\right)}\mathrm{cos}\omega _c\left(j\mathrm{\Delta }tt\right)\end{array}\right)\left(\begin{array}{cc}\beta & \omega _c\\ \omega _c& \beta \end{array}\right)$$
(34)
and
$$\gamma _j=1+e^{2\beta \left(j\mathrm{\Delta }tt\right)}2e^{\beta \left(j\mathrm{\Delta }tt\right)}\mathrm{cos}\omega _c\left(j\mathrm{\Delta }tt\right).$$
(35)
There holds:
$$det\varphi _j^1=\left(det\varphi _j\right)^1=\left(\beta ^2+\omega _c^2\right)\frac{1}{\gamma _j}$$
(36)
and as a consequence the probability distribution of $`\stackrel{}{r}_j`$ becomes
$$\stackrel{~}{w}\left[\stackrel{}{r}_j\right]=\frac{1}{\frac{1}{\beta ^2+\omega _c^2}\gamma _j}\left(\frac{1}{4\pi q\mathrm{\Delta }t}\right)\mathrm{exp}\left\{\frac{\left|\stackrel{~}{A}_j\left(\begin{array}{c}r_j^x\\ r_j^y\end{array}\right)\right|^2}{\gamma _j^24q\mathrm{\Delta }t}\right\}.$$
(37)
In view of
$$\left|\stackrel{~}{A}_j\left(\begin{array}{c}r_j^x\\ r_j^x\end{array}\right)\right|^2=\left(\beta ^2+\omega _c^2\right)\gamma _j\left[\left(r_j^x\right)^2+\left(r_j^y\right)^2\right]$$
(38)
that finally leads to
$$\stackrel{~}{w}\left[\stackrel{}{r}_j\right]=\left(\frac{\beta ^2+\omega _c^2}{4\pi q\mathrm{\Delta }t\gamma _j}\right)\mathrm{exp}\left\{\frac{(\beta ^2+\omega _c^2)\left|\stackrel{}{r}_j\right|^2}{4q\mathrm{\Delta }t\gamma _j}\right\}.$$
(39)
Since this probability distribution is Gaussian with mean zero and variance $`\sigma _j^2=`$ $`2q\mathrm{\Delta }t\frac{1}{\beta ^2+\omega _c^2}\gamma _j`$, the random vector$`\stackrel{}{R}`$ as a sum of independent random variables $`\stackrel{}{r}_j`$ has the distribution
$$w\left(\stackrel{}{R}\right)=\frac{1}{2\pi _j\sigma _j^2}\mathrm{exp}\left(\frac{R_x^2+R_y^2}{2_j\sigma _j^2}\right).$$
(40)
$$\sigma ^2=\underset{j}{}\sigma _j^2=2q\underset{j}{}\mathrm{\Delta }t\frac{1}{\beta ^2+\omega _c^2}\gamma _j.$$
(41)
In the limit of $`\mathrm{\Delta }t0`$ we arrive at the integral
$$\sigma ^2=2q\frac{1}{\beta ^2+\omega _c^2}_0^t\gamma \left(s\right)𝑑s$$
(42)
with $`_0^t\gamma \left(s\right)𝑑s=t+\mathrm{\Theta }`$, where
$$\mathrm{\Theta }=\mathrm{\Theta }(t)=\frac{1}{2\beta }\left(1e^{2\beta t}\right)2\frac{1}{\beta ^2+\omega _c^2}\left[\beta +\left(\omega _c\mathrm{sin}\omega _ct\beta \mathrm{cos}\omega _ct\right)e^{\beta t}\right].$$
(43)
That gives rise to an ultimate form of the transition probability density of the spatial displacement process:
$$P(\stackrel{}{r},t|\stackrel{}{r}_0,t_0=0,\stackrel{}{u}_0)=\frac{1}{4\pi \frac{k_BT}{m}\frac{\beta }{\beta ^2+\omega _c^2}\left(t+\mathrm{\Theta }\right)}\mathrm{exp}(\frac{\left|\stackrel{}{r}\stackrel{}{r}_0\mathrm{\Omega }\stackrel{}{u}_0\right|^2}{4\frac{k_BT}{m}\frac{\beta }{\beta ^2+\omega _c^2}\left(t+\mathrm{\Theta }\right)})$$
(44)
with $`\mathrm{\Omega }=\mathrm{\Omega }(t)`$ defined in Eqs. (25), (27). Notice that an asymptotic diffusion coefficient $`D_B=D\frac{\beta ^2}{\beta ^2+\omega _c^2}`$ of Section 3 (cf. Eq. (20)) appears here as a spatial dispersion - attenuation signature (when $`\omega _c`$ grows up at $`\beta `$ fixed).
The spatial displacement process governed by the above transition probability density surely is not Markovian. That can be checked by inspection: the Chapman-Kolmogorov identity is not valid, like in the standard free Brownian motion example where the Ornstein-Uhlenbeck process induced (sole) spatial dynamics is non-Markovian as well.
## 4 Phase-space process
### 4.1 Free Brownian motion, Kramers equation and local conservation laws
We take advantage of the cylindrical symmetry of our problem, and consider separately the (free) Brownian dynamics in the direction parallel to the magnetic field vector, e.g. along the $`Z`$-axis.
That amounts to invoking a familiar Ornstein-Uhlenbeck process (in velocity/momentum) in its extended phase-space form. In the absence of external forces, the kinetic (Kramers-Fokker-Planck equation) reads:
$$_tW+u_zW=\beta _u(Wu)+q\mathrm{}_uW$$
(45)
where $`q=D\beta ^2`$. Here $`\beta `$ is the friction coefficient, $`D`$ will be identified later with the spatial diffusion constant, and (as before) we set $`D=k_BT/m\beta `$ in conformity with the Einstein fluctuation-dissipation identity.
The joint probability distribution (in fact, density) $`W(z,u,t)`$ for a freely moving Brownian particle which at $`t=0`$ initiates its motion at $`x_0`$ with an arbitrary inital velocity $`u_0`$ can be given in the form of the maximally symmetric displacement probability law:
$$W(z,u,t)=W(R,S)=[4\pi ^2(FGH^2)]^{1/2}exp\{\frac{GR^2HRS+FS^2}{2(FGH^2)}\}$$
(46)
where $`R=zu_0(1e^{\beta t})\beta ^1`$, $`S=uu_0e^{\beta t}`$ while $`F=\frac{D}{\beta }(2\beta t3+4e^{\beta t}e^{2\beta t})`$$`G=D\beta (1e^{2\beta t})`$ and $`H=D(1e^{\beta t})^2`$.
For future reference, let us notice that marginal probablity densities, in the Smoluchowski regime (take for granted that time scales $`\beta ^1`$ and space scales $`(D\beta ^1)^{1/2}`$ are irrelevant ) display familiar forms of the Maxwell-Boltzmann probability density $`w(u,t)=(\frac{m}{2\pi kT})^{1/2}exp(\frac{mu^2}{2k_BT})`$ and the diffusion kernel $`w(z,t)=(4\pi Dt)^{1/2}exp(\frac{z^2}{4Dt})`$ respectively.
A direct evaluation of the first and second local moment of the phase-space probability density gives
$$<u>=𝑑uuW(z,u,t)=w(R)[(H/F)R+u_0e^{\beta t}]$$
(47)
$$<u^2>=𝑑uu^2W(z,u,t)=(\frac{FGH^2}{F}+\frac{H^2}{F^2}R^2)(2\pi F)^{1/2}exp(\frac{R^2}{2F}).$$
(48)
Let us notice that after passing to the diffusion (Smoluchowski) regime, , one readily recovers the local (configuration space conditioned) moment $`<u>_z=\frac{1}{w}<u>`$ to be in the form
$$<u>_z=\frac{z}{2t}=D\frac{w(z,t)}{w(z,t)}$$
(49)
while for the second local moment $`<u^2>_z=\frac{1}{w}<u^2>`$ we would arrive at
$$<u^2>_z=(D\beta D/2t)+<u>_z^2.$$
(50)
By inspection one verifies that the transport (Kramers) equation for $`W(z,u,t)`$ implies local conservation laws:
$$_tw+(<u>_zw)=0$$
(51)
and
$$_t(<u>_zw)+_z(<u^2>_zw)=\beta <u>_zw.$$
(52)
At this point (we strictly follow the moment equations strategy of the traditional kinetic theory of gases and liquids, compare e.g. ) let us introduce the notion of the pressure function $`P_{kin}`$:
$$P_{kin}(z,t)=(<u^2>_z<u>_z^2)w(z,t)$$
(53)
in terms of which we can analyze the local momentum conservation law
$$(_t+<u>_z)<u>_z=\beta <u>_z\frac{P_{kin}}{w}.$$
(54)
One should realize that in the Smoluchowski regime the friction term is cancelled away by a counterterm coming from $`\frac{1}{w}P_{kin}`$ so that
$$(_t+<u>_z)<u>_z=\frac{D}{2t}\frac{w}{w}=\frac{P}{w}$$
(55)
where $`P=D^2w\mathrm{}lnw`$, called osmotic pressure in Ref. , is the net remnant of the kinetic pressure contribution.
Further exploiting the kinetic lore, we can tell few words about the temperature of Brownian particles as opposed to the (equilibrium) temperature of the thermal bath. Namely, in view of (we refer to the Smoluchowski regime with $`t\beta ^1`$) $`P_{kin}(D\beta \frac{D}{2t})w`$ where $`D=\frac{k_BT}{m\beta }`$, we can formally set:
$$k_BT_{kin}=\frac{P_{kin}}{w}(k_BT\frac{D}{2t})<k_BT.$$
(56)
That quantifies the degree of thermal agitation (temperature) of Brownian particles to be less than the thermostat temperature. Heat is continually pumped from the thermostat to the Brownian ”gas”, until asymptotically both temperatures equalize. This may be called a ”thermalization” of Brownian particles. In the process of that ”thermalization” the Brownian ”gas” temperature monotonically grows up until the mean kinetic energy of particles and that of mean flows asymptotically approach the familiar kinetic relationship: $`\frac{w}{2}(<u^2>_z<u>_z^2)𝑑x=k_BT`$, cf. Refs. , for more extended discussion of that medium $``$ particles heat transfer issue and its possible relevance while associating habitual thermal equilibrium conditions with essentially non-equilibrium phenomena.
Remark 1: Once local conservation laws were introduced, it seems instructive to comment on the essentially hydrodynamical features (compressible fluid/gas case) of the problem. Specifically, the ”pressure” term $`Q`$ is quite annoying from the traditional kinetic theory perspective. That is apart from the fact that our local conservation laws have a conspicuous Euler form appropriate for the standard hydrodynamics of gases and liquids. One should become alert that in the present (Brownian) context they convey an entirely different message. For example, in case of normal liquids the pressure is exerted upon any control volume (droplet) by the surrounding fluid. We may interpret that as a compression of a droplet. In case of Brownian motion, we deal with a definite decompression: particles are driven away from areas of higher concentration (probability of occurence). Hence, typically the Brownian ”pressure” is exerted by the droplet upon its surrounding.
Remark 2: The derivation of a hierarchy of local conservation laws (moment equations) for the Kramers equation can be patterned after the standard procedure for the Boltzmann equation. Those laws do not form a closed system and additional specifications (like the familiar thermodynamic equation of state) are needed to that end. In case of the isothermal Brownian motion, when considered in the large friction regime (e.g. Smoluchowski diffusion approximation), it suffices to supplement the Fokker-Planck equation by one more conservation law only to arrive at a closed system, and compare with the discussion of Ref. .
### 4.2 Planar process
Now we shall consider Brownian dynamics in the direction perpendicular to the magnetic field $`\stackrel{}{B}`$, hence (while in terms of configuration-space variables) we address an issue of the planar dynamics. We are interested in the complete phase-space process, hence we need to specify the transition probability density $`P\left(\stackrel{}{r},\stackrel{}{u},t|\stackrel{}{r}_0,\stackrel{}{u}_0,t_0=0\right)`$ of the Markov process conditioned by the initial data $`\stackrel{}{u}=\stackrel{}{u}_0`$ and $`\stackrel{}{r}=\stackrel{}{r}_0`$ at time $`t_0=0`$. That is equivalent to deducing the joint probability distribution $`W\left(\stackrel{}{S,}\stackrel{}{R}\right)`$ of random vectors $`\stackrel{}{S}`$ and $`\stackrel{}{R}`$, previously defined to appear in the form $`\stackrel{}{S}=\stackrel{}{u}\left(t\right)e^{\mathrm{\Lambda }t}\stackrel{}{u}_0`$ and $`\stackrel{}{R}=\stackrel{}{r}\stackrel{}{r}_0\mathrm{\Omega }\stackrel{}{u}_0`$ respectively, cf. Eqs. (15) and (44).
Let us stress that presently, all vectors are regarded as two-dimensional versions (the third component being simply disregarded) of the original random variables we have discussed so far in Sections 2 and 3.
Vectors $`\stackrel{}{S}`$ and $`\stackrel{}{R}`$ both share a Gaussian distribution with mean zero. Consequently, the joint distribution $`W\left(\stackrel{}{S,}\stackrel{}{R}\right)`$ is determined by the matrix of variances and covariances: $`C=\left(c_{ij}\right)=\left(x_ix_j\right)`$, where we abbreviate four phase-space variables in a single notion of $`x=(S_1,S_2,R_1,R_2)`$ and label components of $`x`$ by $`i,j=1,2,3,4`$. In terms of $`\stackrel{}{R}`$ and $`\stackrel{}{S}`$ the covariance matrix $`C`$ reads:
$$C=\left(\begin{array}{cccc}S_1S_1& S_1S_2& S_1R_1& S_1R_2\\ S_2S_1& S_2S_2& S_2R_1& S_2R_2\\ R_1S_1& R_1S_2& R_1R_1& R_1R_2\\ R_2S_1& R_2S_2& R_2R_1& R_2R_2\end{array}\right).$$
(57)
The joint probability distribution of $`\stackrel{}{S}`$ and $`\stackrel{}{R}`$ is given by
$$W\left(\stackrel{}{S,}\stackrel{}{R}\right)=W\left(\stackrel{}{x}\right)=\frac{1}{4\pi ^2}\left(\frac{1}{detC}\right)^{\frac{1}{2}}\mathrm{exp}\left(\frac{1}{2}\underset{i,j}{}c_{ij}^1x_ix_j\right)$$
(58)
where $`c_{ij}^1`$denotes the component of the inverse matrix $`C^1`$.
The probability distributions of $`\stackrel{}{S}`$ and $`\stackrel{}{R}`$, which were established in the previous sections, determine a number of expectation values:
$$gS_1S_1=S_2S_2=\frac{k_BT}{m}\left(1e^{2\beta t}\right)$$
(59)
cf. Eq. (14), while $`S_1S_2=S_2S_1=0`$. Furthermore:
$$fR_1R_1=R_2R_2=2\frac{k_BT}{m}\frac{\beta }{\beta ^2+\omega _c^2}\left(t+\mathrm{\Theta }\right)=2D_B(t+\mathrm{\Theta })$$
(60)
cf. Eqs. (20, (42), (43). In addition we have $`R_1R_2=R_2R_1=0`$.
As a consequence, we are left with only four non-vanishing components of the covariance matrix $`C`$: $`c_{13}=c_{31}=S_1R_1`$, $`c_{14}=c_{41}=S_1R_2`$, $`c_{23}=c_{32}=S_2R_1`$, $`c_{24}=c_{42}=S_2R_2`$ which need a closer examination.
We can obtain those covariances by exploiting a dependence of the random quantities $`\stackrel{}{S}`$ and $`\stackrel{}{R}`$ on the white-noise term $`\stackrel{}{A}\left(s\right)`$ whose statistical properties are known. There follows:
$$S_1=_0^t𝑑se^{\beta \left(ts\right)}\left[\mathrm{cos}\omega _c\left(ts\right)A_1\left(s\right)+\mathrm{sin}\omega _c\left(ts\right)A_2\left(s\right)\right]$$
(61)
$$S_2=_0^t𝑑se^{\beta \left(ts\right)}\left[\mathrm{sin}\omega _c\left(ts\right)A_1\left(s\right)+\mathrm{cos}\omega _c\left(ts\right)A_2\left(s\right)\right]$$
$`R_1`$ $`=`$ $`{\displaystyle _0^t}𝑑s{\displaystyle \frac{1}{\beta ^2+\omega _c^2}}\left[\beta \left(1e^{\beta \left(ts\right)}\mathrm{cos}\omega _c\left(ts\right)\right)+\omega _ce^{\beta \left(ts\right)}\mathrm{sin}\omega _c\left(ts\right)\right]A_1\left(s\right)+`$
$`{\displaystyle _0^t}𝑑s{\displaystyle \frac{1}{\beta ^2+\omega _c^2}}\left[\beta e^{\beta \left(ts\right)}\mathrm{sin}\omega _c\left(ts\right)+\omega _c\left(1e^{\beta \left(ts\right)}\mathrm{cos}\omega _c\left(ts\right)\right)\right]A_2\left(s\right)`$
$`R_2`$ $`=`$ $`{\displaystyle _0^t}𝑑s{\displaystyle \frac{1}{\beta ^2+\omega _c^2}}\left[\omega _c\left(1e^{\beta \left(ts\right)}\mathrm{cos}\omega _c\left(ts\right)\right)+\beta e^{\beta \left(ts\right)}\mathrm{sin}\omega _c\left(ts\right)\right]A_1\left(s\right)+`$
$`{\displaystyle _0^t}𝑑s{\displaystyle \frac{1}{\beta ^2+\omega _c^2}}\left[\omega _ce^{\beta \left(ts\right)}\mathrm{sin}\omega _c\left(ts\right)+\beta \left(1e^{\beta \left(ts\right)}\mathrm{cos}\omega _c\left(ts\right)\right)\right]A_2\left(s\right).`$
Multiplying together suitable components of vectors $`\stackrel{}{S}`$ and $`\stackrel{}{R}`$ and taking averages of those products in conformity with the rules $`A_i\left(s\right)=0`$ and $`A_i\left(s\right)A_j\left(s^{}\right)=2q\delta _{ij}\delta \left(ss^{}\right)`$, where $`q=\frac{k_BT}{m}\beta `$, $`i,j=1,2,3`$, we arrive at:
$$hR_1S_1=R_2S_2=2q\frac{1}{\beta ^2+\omega _c^2}_0^tdse^{\beta \left(ts\right)}[\beta \mathrm{cos}\omega _c(ts)+$$
(62)
$$\omega _c\mathrm{sin}\omega _c(ts)\beta e^{\beta \left(ts\right)}]=q\frac{1}{\beta ^2+\omega _c^2}(12e^{\beta t}\mathrm{cos}\omega _ct+e^{2\beta t})$$
and
$$kR_1S_2=R_2S_1=2q\frac{1}{\beta ^2+\omega _c^2}_0^tdse^{\beta \left(ts\right)}[\beta \mathrm{sin}\omega _c(ts)+$$
(63)
$$\omega _c\mathrm{cos}\omega _c(ts)\omega _ce^{\beta \left(ts\right)}]=q\frac{1}{\beta ^2+\omega _c^2}[2e^{\beta t}\mathrm{sin}\omega _ct\frac{\omega _c}{\beta }(1e^{2\beta t})].$$
The covariance matrix $`C=\left(c_{ij}\right)`$ has thus the form
$$C=\left(\begin{array}{cccc}g& 0& h& k\\ 0& g& k& h\\ h& k& f& 0\\ k& h& 0& f\end{array}\right)$$
(64)
while its inverse $`C^1`$ reads as follows:
$$C^1=\frac{1}{detC}\left(fgh^2k^2\right)\left(\begin{array}{cccc}f& 0& h& k\\ 0& f& k& h\\ h& k& g& 0\\ k& h& 0& g\end{array}\right)$$
(65)
where $`detC=\left(fgh^2k^2\right)^2`$.
The joint probability distribution of $`\stackrel{}{S}`$ and $`\stackrel{}{R}`$ can be ultimately written in the form:
$$W(\stackrel{}{S},\stackrel{}{R})=$$
(66)
$$\frac{1}{4\pi ^2\left(fgh^2k^2\right)}\mathrm{exp}\left(\frac{f\left|\stackrel{}{S}\right|^2+g\left|\stackrel{}{R}\right|^22h\stackrel{}{S}\stackrel{}{R}+2k\left(\stackrel{}{S}\times \stackrel{}{R}\right)_{i=3}}{2\left(fgh^2k^2\right)}\right).$$
(67)
In the above, all vector entries are two-dimensional. The specific $`i=3`$ vector product coordinate in the exponent is simply an abbreviation for the (ordinary $`R^3`$-vector product) procedure that involves merely first two components of three-dimensional vectors (the third is then arbitrary and irrelevant), hence effectively involves our two-dimensional $`\stackrel{}{R}`$ and $`\stackrel{}{S}`$.
### 4.3 Kramers equation and local conservation laws for the planar motion
For the purpose of evaluating local velocity averages of the Kramers equation, we need to extract the marginal configuration space distribution. Let us notice that
$$W(\stackrel{}{S},\stackrel{}{R})d\stackrel{}{S}=w(\stackrel{}{R})=P(\stackrel{}{r},t|\stackrel{}{r}_0,t_0=0,\stackrel{}{u}_0)$$
(68)
where the last transition probability density entry coincides with that of Eq. (44).
Let us introduce an auxiliary (weighted) distribution:
$$\stackrel{~}{W}\left(\stackrel{}{S}|\stackrel{}{R}\right)=\frac{W(\stackrel{}{S},\stackrel{}{R})}{W(\stackrel{}{S},\stackrel{}{R})𝑑\stackrel{}{S}}=$$
(69)
$$\frac{1}{2\pi \frac{1}{f}\left(fgh^2k^2\right)}\mathrm{exp}\left(\frac{\left|\stackrel{}{S}\stackrel{}{m}\right|^2}{2\frac{1}{f}\left(fgh^2k^2\right)}\right)$$
where
$$\stackrel{}{m}=\frac{1}{f}(hR_1kR_2,hR_2+kR_1)$$
(70)
and we recall that $`\stackrel{}{S}=\stackrel{}{u}\left(t\right)e^{\beta t}U\left(t\right)\stackrel{}{u}_0`$ and $`\stackrel{}{R}=\stackrel{}{r}\stackrel{}{r}_0\mathrm{\Omega }\stackrel{}{u}_0`$.
The local expectation values (compare e.g. calculations of the previous subsection) read: $`u_i_\stackrel{}{R}=u_i\stackrel{~}{W}𝑑\stackrel{}{u}`$ and $`u_i^2_\stackrel{}{R}=u_i^2\stackrel{~}{W}𝑑\stackrel{}{u}`$ where $`i=1,2`$.
By evaluating those averages we get:
$$\stackrel{}{u}_\stackrel{}{R}=(u_1_\stackrel{}{R},u_2_\stackrel{}{R})=e^{\beta t}U\left(t\right)\stackrel{}{u}_0+\stackrel{}{m}$$
(71)
and
$$u_1^2_\stackrel{}{R}u_1_\stackrel{}{R}^2=u_2^2_\stackrel{}{R}u_2_\stackrel{}{R}^2=\frac{1}{f}\left(fgh^2k^2\right)$$
(72)
The Fokker-Planck-Kramers equation, appropriate for the planar dynamics in its phase space version, reads
$$\frac{W}{t}+\stackrel{}{u}_\stackrel{}{r}W=\beta _\stackrel{}{u}\left(W\stackrel{}{u}\right)\omega _c\left[_\stackrel{}{u}\times W\stackrel{}{u}\right]_{i=3}+q_\stackrel{}{u}^2W$$
(73)
where the troublesome again (in the planar case all vectors are two-dimensional) vector product third component stands for
$$\left[_\stackrel{}{u}\times W\stackrel{}{u}\right]_{i=3}=\frac{}{u_1}\left(Wu_2\right)\frac{}{u_2}\left(Wu_1\right).$$
(74)
First two moment equations are easily derivable. Namely, the continuity (0-th moment) and the momentum conservation (first moment) equations come out in the form
$$_tw+\stackrel{}{}\left[\stackrel{}{u}_\stackrel{}{R}w\right]=0$$
(75)
and
$$_t\left[u_1_\stackrel{}{R}w\right]+\frac{}{r_1}\left[u_1^2_\stackrel{}{R}w\right]+\frac{}{r_2}\left[u_1_\stackrel{}{R}u_2_\stackrel{}{R}w\right]=\beta u_1_\stackrel{}{R}w+\omega _cu_2_\stackrel{}{R}w$$
(76)
$$_t\left[u_2_\stackrel{}{R}w\right]+\frac{}{r_2}\left[u_2^2_\stackrel{}{R}w\right]+\frac{}{r_1}\left[u_1_\stackrel{}{R}u_2_\stackrel{}{R}w\right]=\beta u_2_\stackrel{}{R}w\omega _cu_1_\stackrel{}{R}w.$$
That implies
$$\left[_t+u_1_\stackrel{}{R}\frac{}{r_1}+u_2_\stackrel{}{R}\frac{}{r_2}\right]u_1_\stackrel{}{R}=\beta u_1_\stackrel{}{R}+\omega _cu_2_\stackrel{}{R}\frac{1}{w}\frac{}{r_1}\left[\left(u_1^2_\stackrel{}{R}u_1_\stackrel{}{R}^2\right)w\right]$$
(77)
$$\left[_t+u_1_\stackrel{}{R}\frac{}{r_1}+u_2_\stackrel{}{R}\frac{}{r_2}\right]u_2_\stackrel{}{R}=\beta u_2_\stackrel{}{R}\omega _cu_1_\stackrel{}{R}\frac{1}{w}\frac{}{r_2}\left[\left(u_2^2_\stackrel{}{R}u_2_\stackrel{}{R}^2\right)w\right]$$
which finally sums up to a local momentum conservation law (here, the standard $`R^3`$ vector product on the right-hand-side contributes its first and second components only)
$$\left[_t+\stackrel{}{u}_\stackrel{}{R}\stackrel{}{}\right]\stackrel{}{u}_\stackrel{}{R}=\mathrm{\Lambda }\stackrel{}{u}_\stackrel{}{R}\frac{1}{w}\stackrel{}{}\stackrel{}{P}_{kin}=\beta \stackrel{}{u}_\stackrel{}{R}+\frac{q_e}{mc}\stackrel{}{u}_\stackrel{}{R}\times \stackrel{}{B}\frac{1}{w}\stackrel{}{}\stackrel{}{P}_{kin}$$
(78)
where $`\stackrel{}{P}_{kin}`$ has tensor components $`P_{ij}^{kin}`$, and $`\stackrel{}{}\stackrel{}{P}_{kin}`$ stands for a vector whose $`i`$-th component is equal $`_j\frac{P_{ij}^{kin}}{r_j}`$ and $`i,j=1,2`$. Here, obviously $`P_{ij}^{kin}=\left(u_iu_i_\stackrel{}{R}\right)\left(u_ju_j_\stackrel{}{R}\right)_\stackrel{}{R}w`$ and only diagonal entries do not vanish. Clearly $`P_{ii}^{kin}=\left(u_i^2_\stackrel{}{R}u_i_\stackrel{}{R}^2\right)w`$.
Because of
$$\sigma ^2=u_1^2_\stackrel{}{R}u_1_\stackrel{}{R}^2=u_2^2_\stackrel{}{R}u_2_\stackrel{}{R}^2=\frac{1}{f}\left(fgh^2k^2\right)=g\frac{h^2+k^2}{f}$$
(79)
we can introduce again $`\sigma ^2=\frac{P_{kin}}{w}=k_BT_{kin}`$ and pass to an asymptotic regime $`t>>t_c=\frac{1}{\beta }`$. Then, we obtain Eq. (56) to be valid in the present case as well thus quantifying an overall (magnetic field independent) heating process involved .
In that asymptotic regime we have $`\sigma ^2=D\beta \frac{D}{2t}`$ and by employing an asymptotic form of $`w(\stackrel{}{R})`$, Eq. (44) we recover:
$$\stackrel{}{}\stackrel{}{P}_{kin}=\left(D\beta \frac{D}{2t}\right)\frac{\stackrel{}{}w}{w}$$
(80)
together with
$$\mathrm{\Lambda }\stackrel{}{u}_\stackrel{}{R}=D\beta \frac{\stackrel{}{}w}{w}.$$
(81)
So, asymptotically ($`t>>\beta ^1`$) the momentum conservation law takes the form (to be compared with considerations of section $`4.1`$)
$$\left[_t+\stackrel{}{u}_\stackrel{}{R}\stackrel{}{}\right]\stackrel{}{u}_\stackrel{}{R}=\frac{D}{2t}\frac{\stackrel{}{}w}{w}$$
(82)
However an asymptotic regime does not yet imply that the right-hand-side of Eq. (83) represents an acceptable ”osmotic pressure” gradient contribution. We additionally need a large friction regime to deal with a consistent picture of a Markov diffusion process in the Smoluchowski form. Indeed, to reproduce a universal (Ref. ) pressure-type functional dependence on $`P`$ we must employ a suitable form of the diffusion coefficient. The usage of $`D`$ alone to define an ”osmotic pressure” implies an apparent failure. On the other hand, the usage of $`D_B=D\frac{\beta ^2}{\beta ^2+\omega _{}^{2}{}_{c}{}^{}}`$ as suggested by Eq. (44) leads to:
$$P=D_B^2\mathrm{\Delta }\mathrm{ln}w\frac{\stackrel{}{}P}{w}=\frac{D_B}{2t}\frac{\stackrel{}{}w}{w}.$$
(83)
Then however the momentum conservation law displays a supplementary scaling of the osmotic pressure contribution, which may trivialize (die out) only for large values of $`\beta `$ (an ultimate Smoluchowski regime). Namely, we have:
$$\left[_t+\stackrel{}{u}_\stackrel{}{R}\stackrel{}{}\right]\stackrel{}{u}_\stackrel{}{R}=\left(\frac{\beta ^2}{\beta ^2+\omega _c^2}\right)^1\frac{\stackrel{}{}P}{w}.$$
(84)
Such process is yet non-Markovian and its approximation by the Smoluchowski process becomes reliable when $`\beta `$ is large while $`\omega _c`$ is kept moderately small.
Basically, the large friction regime cancels all rotational features (arising due to the Lorentz force) on very short time scales. If we are satisfied with the non-Markovian regime of moderate friction but arbitrarily varying $`\omega _c`$ (i.e. $`B`$) then, in conformity with Eq. (78), mean flows would display signatures of rotation that is bound to die out asymptotically. This effect can be analyzed in $`R^3`$ by observing that a three-dimensional extension of the vector $`\stackrel{}{m}`$ of Eqs. (70), (71) asymptotically reads $`\stackrel{}{m}=\frac{1}{2t\beta }\mathrm{\Lambda }(\stackrel{}{r}\stackrel{}{r}_0)`$ where $`\mathrm{\Lambda }`$ comes from Eq. (3) and $`(\stackrel{}{r}\stackrel{}{r}_0)R^3`$. In view of that, we have (cf. Eq. (71)) $`curl\stackrel{}{u}_\stackrel{}{R}curl\stackrel{}{m}(0,0,\frac{\omega _c}{t\beta })`$ and the circulation asymptotically vanishes. The effect can be slowed down by a suitable adjustment of $`\omega _c`$ against $`\beta `$.
Acknowledgement: One of the authors (P. G.) receives financial support from the KBN research grant No. 2 PO3B 086 16.
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# On the existence and convergence of polyhomogeneous expansions of zero-rest-mass fields.
## 1 Introduction
The analysis of the asymptotic behaviour of asymptotically flat spacetimes and the diverse fields propagating on it has been carried out by assuming that both the gravitational field and the test fields (e.g. zero-rest-mass fields) are smooth at the boundary of the *unphysical* (i.e. conformally rescaled) spacetime manifold (null infinity, $`I`$) . This smoothness at null infinity implies that the fields in the physical spacetime manifold have the peeling behaviour. However, this assumption may be too strong for a number of physical applications (see for example and references within). In particular, the case for the so-called *polyhomogeneous spacetimes and zero-rest-mass fields* has been rised . A generic peeling (smooth) field possesses an analytic expansion around null infinity (given by $`\mathrm{\Omega }=0`$). By contrast, the expansions of polyhomogeneous fields around $`I`$ are Laurent-like (so that the field in the physical spacetime is non-peeling), and also contain powers of $`\mathrm{ln}`$.
So far, the study of polyhomogeneous fields has been carried out assuming that the diverse fields admit expansions of the required form. However, the expansions are formal in the sense that the convergence (and hence the existence and uniqueness) of the series has not been discussed. This point, so many times overlooked, is of fundamental importance. It could be well the case that none of the formal expansions converge!
The objective of this article is to provide a first step towards the understanding of the convergence of polyhomogeneous fields. Here, I analyse the possible convergence of a polyhomogeneous spin-1 zero-rest-mass field (Maxwell field) propagating on an asymptotically simple spacetime. The spin-1 zero-rest-mass field has been chosen in order to ease the discussion and the calculations, but the treatment here discussed can be directly extended to the case of say, a spin-2 field. It is also expected that similar techniques can be used to analyse the analogous question for the gravitational field. However, in that case one should expect a number of complications. Furthermore, again for concreteness, only the simplest case of spin-1 zero-rest-mass fields will be discussed: the so-called minimally polyhomogeneous. This case contains all the difficulties one could expect from the generic polyhomogeneous situation, but the algebra is much simpler!
The question of the convergence of the polyhomogeneous expansions will be tackled by posing an *asymptotic characteristic initial value problem* for the minimally polyhomogeneous spin-1 zero-rest-mass field, and proving that such a problem has a (unique) solution in a given neighbourhood of the spacetime which is polyhomogeneous (i.e. the solution has convergent polyhomogeneous expansion in the given neighbourhood).
## 2 Preliminaries, notation, conventions
Let $`(\stackrel{~}{M},\stackrel{~}{g}_{ab})`$ be an asymptotically simple spacetime. It will be assumed that null infinity is smooth (non-polyhomogeneous). Let $`(M,g_{ab})`$ be the unphysical spacetime obtained by means of the conformal rescaling $`g_{ab}=\mathrm{\Omega }^2\stackrel{~}{g}_{ab}`$. The spacetime will be taken to be completely known. In particular this means that the asymptotic expansions of the spin coefficients, and the tetrad functions are known. Let $`𝒩_0`$ be the initial (future oriented) null hypersurface, and let $`𝒵_0`$ be its intersection with future null infinity, $`𝒵_0=𝒩_0I^+`$. Following Penrose & Rindler’s notation, $``$ will denote equality at null infinity.
To describe the spacetime and the different fields involved, I will make use of the coordinates and null tetrad described by Friedrich & Stewart. In their original version, these coordinates and null tetrad were used to describe the spacetime between two intersecting null hypersurfaces . Here they are adapted to describe the region of spacetime in the causal past of the intersection of future null infinity ($`I^+`$), and a future oriented null hypersurface. The construction is of interest in itself, and its highlights are given for the sake of reference.
The physical spacetime $`(\stackrel{~}{M},\stackrel{~}{g}_{ab})`$ is asymptotically flat, thus $`𝒵_0`$ has the topology of the 2-sphere, $`𝒮^2`$. Let $`p`$ be an arbitrary point of $`𝒵_0`$, and let $`\gamma `$ and $`\gamma ^{}`$ be the null generators of $`𝒩_0`$ and $`I^+`$ through $`p`$ respectively. A null tetrad $`\{l^a,n^a,m^a,\overline{m}^a\}`$ at $`p`$ can be chosen such that $`m^a`$ and $`\overline{m}^a`$ span $`\text{T}_p(𝒵_0)`$, and $`l^a`$ and $`n^a`$ are tangent to $`\gamma `$ and $`\gamma ^{}`$ respectively. One can also choose an arbitrary coordinate system, $`(x^\alpha )`$, at $`𝒵_0`$. Let $`u`$ (retarded time) be a parameter along the null geodesic generators $`\gamma ^{}`$ of $`I^+`$, such that $`u=0`$ at $`p`$. Note that in principle $`u`$ is not an affine parameter. Repeating this construction for different points,$`p`$, one obtains a scalar field that vanishes precisely at $`𝒵_0`$. The coordinates $`(x^\alpha )`$ can now be extended to the remainder of $`I^+`$ by requiring the coordinates to be constant along the generators of $`I^+`$. In this way we have constructed a coordinate chart $`(u,x^\alpha )`$ for $`I^+`$ valid at least around $`𝒵_0`$. One can undertake an analogous construction for the generators $`\gamma `$ of $`𝒩_0`$, obtaining a parameter $`v`$ (advanced time) such that at $`p𝒵_0`$, $`v=0`$. Again, the parameter $`v`$ is not affine in principle. Thus, one ends up with a coordinate chart $`(v,x^\alpha )`$ for $`𝒩_0`$.
Now, on $`I^+`$ there is a 1-parameter family of spacelike hypersurfaces, $`𝒮_u`$ (cuts) on which $`u`$ is constant. For each $`𝒮_u`$ there is a unique null hypersurface $`𝒩_u`$ containing $`𝒮_0`$ ($`𝒮_u𝒩_u`$) extending into the interior of the unphysical spacetime. This yields a foliation of the spacetime by null hypersurfaces at least close to $`I^+`$ —note that the foliation could give rise to caustics and cusps. Similarly, one can obtain another 1-parameter family of null hypersurfaces $`𝒩_v^{}`$ for which $`v`$ is constant and that intersect $`𝒩_0`$. In particular one has that $`𝒩_0^{}=I^+`$. The coordinates $`(x^\alpha )`$ can now also be extended to the interior of the spacetime by requiring them to be constant along the null generators of $`𝒩_u`$. Note that $`(x^\alpha )`$ are not constant along the generators of the $`𝒩_v^{}`$ foliation (except for $`I^+`$) for the null vectors $`l^a`$ and $`n^a`$ which will be taken as tangent to the null generators of $`𝒩_u`$ and $`𝒩_v^{}`$ respectively do not in general commute.
To extend the null tetrad defined in $`𝒵_0`$, one begins by defining $`l_a`$ on $`M`$ by $`l_a=_au`$. Therefore $`l^a`$ is tangent to the null generators of $`𝒩_0`$, and is parallel to the $`l^a`$ already defined at $`𝒵_0`$. A boost on the null tetrad defined at $`𝒵_0`$ can be used to make the two vectors equal. Thus, there exists a scalar function $`Q`$ such that $`l^a_a=Q_v`$. It can be shown that $`Q=𝒪(1)`$. Next, the $`n_a`$ is defined on $`M`$ by $`n_a=Q^1_v`$, so that $`n^a`$ is tangent to the null generators of $`N_v^{}`$ and the normalisation condition $`l^an_a=1`$ is satisfied. Together with $`n^an_a=0`$ this implies, $`n^a_a=_u+c^\alpha _\alpha `$. Due to the constancy of the $`(x^\alpha )`$ coordinates along the generators of $`I^+`$, one has that $`C^\alpha 0`$.
Following the standard NP notation, let $`D=l^a_a`$ and $`\mathrm{\Delta }=n^a_a`$. The vector $`m^a`$ defined so far only on $`𝒵_0`$ is propagated onto the remaining of $`I^+`$ via the equation $`\mathrm{\Delta }m^a=\tau n^a`$. Finally $`m^a`$ is propagated from $`I^+`$ to the interior of $`M`$ via $`Dm^a=\overline{\pi }l^a`$. This together with the nullity and orthogonality requirements imply that $`m^a_a=\xi ^\alpha _\alpha `$.
Hence, the aforediscussed null tetrad has the form:
$`l^a=Q\delta _v^a,`$ (1)
$`n^a=\delta _u^a+C^\alpha \delta _\alpha ^a,`$ (2)
$`m^a=\xi ^\alpha \delta _\alpha ^a,`$ (3)
$`\overline{m}^a=\overline{\xi }^\alpha \delta _\alpha ^a,`$ (4)
where $`C^\alpha 0`$, and $`Q`$, $`\xi ^\alpha `$ are of order $`𝒪(1)`$. From this tetrad construction one also deduces that $`\kappa =\nu =0`$, $`\rho =\overline{\rho }`$, $`\mu =\overline{\mu }`$, $`ϵ=0`$, $`\tau =\overline{\alpha }+\beta `$ in $`M`$. In particular one also has that all the spin coefficients with the exception of $`\sigma `$ vanish at $`I^+`$: that is, they are of order $`𝒪(\mathrm{\Omega })`$. On the other hand, $`\sigma `$ is $`𝒪(1)`$.
Finally, notice that the aforediscussed construction fixes also the associated spin basis $`\{o^A,\iota ^A\}`$.
### 2.1 Polyhomogeneous zero-rest-mass fields
A spin-s zero-rest-mass field is a totally symmetric spinor field ($`\varphi _{AB\mathrm{}C}=\varphi _{(AB\mathrm{}C)}`$), satisfying:
$$^{AA^{}}\varphi _{AB\mathrm{}C}=0,$$
(5)
where the indices $`AB\mathrm{}C`$ number $`2s`$. In order to ease the discussion, in this article we will restrict ourselves to spin 1 zero-rest-mass fields (Maxwell field). The extension of the results to spin 2 fields (linear gravity) is direct. The components of the spin 1 zero-rest-mass field in the aforementioned spin basis are: $`\varphi _0=\varphi _{AB}o^Ao^B`$, $`\varphi _1=\varphi _{AB}o^A\iota ^B`$, $`\varphi _2=\varphi _{AB}\iota ^A\iota ^B`$.
Let us consider zero-rest-mass fields belonging to $`C_{loc}^{\mathrm{}}(M)`$, i.e. fields that are infinite differentiable (smooth) in the interior of the unphysical spacetime ($`\stackrel{}{M}=\stackrel{~}{M}`$) but need not to be so at the boundary ($`I`$) of the spacetime; this is the meaning of the $`loc`$ label. In particular our attention will be centered on the so-called *polyhomogeneous fields*, which are $`C_{loc}^{\mathrm{}}(M)`$ zero-rest-mass fields with asymptotic expansions in terms of powers of a particular parameter (e.g. the conformal factor, and affine parameter of outgoing null geodesics, a luminosity parameter, an advanced time), and powers of the logarithm.
The most general polyhomogeneous spin-1 zero-rest-mass field is of the form,
$`\stackrel{~}{\varphi }_0=\varphi _0^2\mathrm{\Omega }^2+\varphi _0^3\mathrm{\Omega }^3+\mathrm{},`$ (6)
$`\stackrel{~}{\varphi }_1=\varphi _1^2\mathrm{\Omega }^2+\varphi _1^3\mathrm{\Omega }^3+\mathrm{},`$ (7)
$`\stackrel{~}{\varphi }_2=\varphi _2^1\mathrm{\Omega }+\varphi _2^2\mathrm{\Omega }^2+\mathrm{},`$ (8)
where the $`\varphi _n^i`$’s are polynomials in $`z=\mathrm{ln}\mathrm{\Omega }`$ with coefficients in $`C^\omega (\text{}\times 𝒮^2)`$. It is customary to require absolute convergence of the polyhomogeneous series and its derivatives to all orders .
Note that the field defined by equations (6)-(8) is clearly non-peeling. A peeling field should behave like $`\varphi _n=𝒪(\mathrm{\Omega }^{(3n)})`$. Furthermore, looking at the spin-1 zero-rest-mass field in the unphysical (conformally rescaled) spacetime one has,
$`\varphi _0=\varphi _0^2\mathrm{\Omega }^1+\varphi _0^3+\mathrm{},`$ (9)
$`\varphi _1=\varphi _1^2+\varphi _1^3\mathrm{\Omega }+\mathrm{},`$ (10)
$`\varphi _2=\varphi _2^1+\varphi _2^2\mathrm{\Omega }+\mathrm{},`$ (11)
so that the field is not regular (diverges!) at $`I^+`$ (where $`\mathrm{\Omega }=0`$), first because of the term $`\varphi _0^2\mathrm{\Omega }^1`$, and also because of the presence of logarithmic terms in $`\varphi _1^2`$.
In order to ease the discussion, the discussion in this article will be reduced to the case of the so-called minimally polyhomogeneous fields. An extension to more general polyhomogeneous fields can be easily done. A minimally polyhomogeneous spin-1 zero-rest-mass field is of the form (in the physical spacetime):
$`\stackrel{~}{\varphi }_0=\varphi _0^{2,0}\mathrm{\Omega }^2+\left(\varphi _0^{3,1}\mathrm{ln}\mathrm{\Omega }+\varphi _0^{3,0}\right)\mathrm{\Omega }^3+\mathrm{},`$ (12)
$`\stackrel{~}{\varphi }_1=\left(\varphi _1^{2,1}\mathrm{ln}\mathrm{\Omega }+\varphi _1^{2,0}\right)\mathrm{\Omega }^2+\mathrm{},`$ (13)
$`\stackrel{~}{\varphi }_2=\varphi _0^{1,0}\mathrm{\Omega }+\mathrm{},`$ (14)
which in the unphysical spacetime looks like,
$`\varphi _0=\varphi _0^{2,0}\mathrm{\Omega }^1+\left(\varphi _0^{3,1}\mathrm{ln}\mathrm{\Omega }+\varphi _0^{3,0}\right)+\mathrm{},`$ (15)
$`\stackrel{~}{\varphi }_1=\left(\varphi _1^{2,1}\mathrm{ln}\mathrm{\Omega }+\varphi _1^{2,0}\right)+\mathrm{},`$ (16)
$`\stackrel{~}{\varphi }_2=\varphi _0^{1,0}+\left(\varphi _2^{2,1}\mathrm{ln}\mathrm{\Omega }+\varphi _2^{2,0}\right)\mathrm{\Omega }\mathrm{}.`$ (17)
Using the “constraint equations”, i.e. the $`D`$-Maxwell equations one can deduce several relations connecting the diverse coefficients of the different components of the spin-1 zero-rest-mass field. In particular one has that
$$\varphi _1^{2,1}=\overline{ð}\varphi _0^{2,0},$$
(18)
where $`\overline{ð}`$ is the ethbar differential operator (see for example ). This relation will be used later.
So far, the asymptotic expansions have been in $`\mathrm{\Omega }`$. However, it can be seen that,
$$\mathrm{\Omega }=a(u,x^\alpha )v+\mathrm{},$$
(19)
so that one could alternatively use the retarded time as the expansion parameter. *This will be done in the remainder of the article*. One can even redefine $`a(u,x^\alpha )`$ so that $`a(u,x^\alpha )=1`$.
A minimally polyhomogeneous field contains at the most first powers of $`\mathrm{ln}`$. If one provides $`\stackrel{~}{\varphi }_0`$ on an initial null hypersurface $`\stackrel{~}{𝒩}_0`$, and the coefficients $`\varphi _1^{2,0}`$ at $`x^0=0`$ and $`\varphi _2^{1,0}`$ for all times, one can use the equations (22) and (23) on to obtain the components $`\stackrel{~}{\varphi }_1`$ and $`\stackrel{~}{\varphi }_2`$ on $`\stackrel{~}{𝒩}_0`$, and then the evolution equations (20) and (21) to generate the field by “slices”. This procedure allows us to calculate formal polyhomogeneous expansions for the zero-rest-mass field. As mentioned previously in the introduction, the question to be investigated here regards the convergence of such formal polyhomogeneous expansions, i.e. we want to know if such fields exist, at least for a region of the spacetime very close to the initial null hypersurface $`𝒩_0`$.
## 3 Symmetric hyperbolicity and the asymptotic characteristic initial value problem
The question of the convergence of the polyhomogeneous expansions of the zero-rest-mass fields discussed in the previous section will be addressed by posing a so-called asymptotic characteristic initial value problem, i.e. an initial value problem with initial data given on a future oriented null hypersurface (a characteristic of the zero-rest-mass field equations), and on future null infinity (which, by the way is also a characteristic of the field equations). Once a suitable initial value problem has been posed, an existence result will be proven. This existence result will exhibit the convergence of the polyhomogeneous series. The aforesaid initial value problem is *non-regular* in a sense to be clarified later. Hence the standard techniques used to prove existence results for analytic fields cannot be directly used. Some manipulations will have to be carried. In order to understand the rationale behind the forthcoming discussion, the standard regular problem is briefly reviewed.
### 3.1 The initial value problem with regular initial data
Under the tetrad choice described in the previous sections, the spin-1 zero rest-mass field equations (aka source-free Maxwell equations) are given by:
$`\mathrm{\Delta }\varphi _0\delta \varphi _1=(2\gamma \mu )\varphi _02\tau \varphi _1+\sigma \varphi _2,`$ (20)
$`\mathrm{\Delta }\varphi _1\delta \varphi _2=2\mu \varphi _1+(2\beta \tau )\varphi _2,`$ (21)
$`D\varphi _2\overline{\delta }\varphi _1=\lambda \varphi _0+2\pi \varphi _1+\rho \varphi _2,`$ (22)
$`D\varphi _1\overline{\delta }\varphi _0=(\pi 2\alpha )\varphi _0+2\rho \varphi _1.`$ (23)
Recall that the zero-rest-mass fields are totally symmetric, and thus in our case $`\varphi _{AB}o^A\iota ^B=\varphi _{AB}\iota ^Ao^B`$. Therefore the system (20)-(23) is overdetermined: 4 equations for 3 complex components. This difficulty is overcome by discarding the last equation (23). Then, one has to verify that this equation is implied by the other ones if it held initially. This can be done by defining,
$$\mathrm{\Phi }=D\varphi _1\overline{\delta }\varphi _0(\pi 2\alpha )\varphi _02\rho \varphi _1.$$
(24)
Using the field equations, the commutators and the Ricci identities (NP field equations), it is not difficult to show that,
$$_u\mathrm{\Phi }=0.$$
(25)
Thus, if equation (23) held initially, it will also held for later retarded times.
The principal part (symbol) of the reduced system (20)-(22) is contained in the terms
$`\mathrm{\Delta }\varphi _0\delta _1,`$ (26)
$`\mathrm{\Delta }\varphi _1\delta _2,`$ (27)
$`D\varphi _2\overline{\delta }\varphi _1,`$ (28)
and can be written concisely in a matricial way as $`A^a\varphi _{b,a}`$ where $`\varphi =(\varphi _0,\varphi _1,\varphi _2)`$, and
$$A^u=\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right),$$
(29)
$$A^v=\left(\begin{array}{ccc}0& 0& 0\\ 0& 0& 0\\ 0& 0& Q\end{array}\right),$$
(30)
and
$$A^\alpha =\left(\begin{array}{ccc}C^\alpha & 0& 0\\ 0& C^\alpha & \xi ^\alpha \\ 0& \overline{\xi }^\alpha & 0\end{array}\right).$$
(31)
Whence the reduced system (20)-(22) is clearly symmetric hyperbolic .
Given suitable initial data on a future oriented null hypersurface $`𝒩_0`$ and on future null infinity $`I^+`$ one would like to obtain to obtain the spin-1 zero-rest-mass field to the past of the two initial hypersurfaces $`J^{}(I^+𝒩_0)`$. This is the *asymptotic characteristic initial value problem*. Alternatively one could give initial data on a past oriented null hypersurface $`𝒩_0^{}`$ and on $`I^{}`$ and reconstruct the field on $`J^+(I^{}𝒩_0^{})`$. For the sake of concreteness, here I will restrict to the first possibility.
From the structure of the principal part of the spin-1 zero-rest-mass field equations, one finds that the right way of setting such an *asymptotic characteristic initial value problem* for $`\varphi _{AB}`$ is set by prescribing:
1. $`\varphi _0`$ on $`𝒩_0`$,
2. $`\varphi _1`$ on $`𝒵_0=𝒩_0I^+`$,
3. and $`\varphi _2`$ on $`I^+`$.
If the initial data on $`𝒩_0`$ ($`\varphi _0`$) peels, i.e. $`\stackrel{~}{\varphi }_0=𝒪(\mathrm{\Omega }^3)`$ or alternatively $`\varphi _0=𝒪(1)`$), then $`\varphi _0`$ is bounded at $`I^+`$, and one has a *regular* asymptotic characteristic initial value problem. Furthermore, if the data on $`𝒩_0`$ is analytic with respect to v, and the data on $`I^+`$ is analytic with respect to u, one can make use of Duff and Friedrich’s techniques <sup>1</sup><sup>1</sup>1Duff showed that the ideas of the Cauchy-Kovalevskaya theorem for a Cauchy problem can be extended to the case ofa characteristic initial value problem. In particular he considered linear systems with data specified on an initial hypersurface which is part characteristic and part spacelike. The construction for the full chacteristic case has been done by Friedrich. theorem and prove that there is a unique analytic ($`C^\omega `$) solution to the posed initial value problem. The idea behind the proof is to construct another system of partial differential equations whose solutions can be shown to exist, be analytic and with an expansion that majorises the formal series solutions of the original system of equations. This implies the analyticity (and existence) of the solutions. In order to implement this construction the different terms in the equations are required to be analytic as well. The Duff-Friedrich’s construction works for quasilinear equations. Here, the situation is much simpler as the zero-rest-mass field equations are linear.
Further refinements can be implemented in order to show existence for the case when the initial data is $`C^{\mathrm{}}`$ rather than analytic . In such a case the field is found to be $`C^{\mathrm{}}`$. However, these further results may not concern us here.
### 3.2 The non-regular initial value problem
So much for the so-called regular initial value problem. What happens if the initial data over $`𝒩_0`$ is non-peeling, like for instance $`\varphi _0=\varphi _0^{2,0}v^1+\mathrm{}`$? As discussed elsewhere , such initial data gives rise to logarithms in the asymptotic expansions of the remaining components of the zero-rest-mass field spinor $`\varphi _{AB}`$, and in $`\varphi _0`$ itself for later times, even in the case where the logarithmic terms were not present in the initial data. An initial value problem with this kind of data that is not bounded at $`I^+`$, and henceforth it will be called *non-regular*.
In the case that concerns us here, that of a minimally polyhomogeneous spin-1 zero-rest-mass field, the non-regular nature of $`\varphi _0`$ and $`\varphi _1`$ at $`I^+`$ precludes us from setting an initial value problem as the one discussed above, and using directly Duff and Friedrich’s techniques.
Assume for a moment that we have a polyhomogeneous spin-1 zero-rest-mass field. Due to the absolute convergence of the formal series expansion, one can reshuffle the terms and write:
$`\varphi _0=\underset{0}{\overset{(1)}{\varphi }}\mathrm{ln}v+\underset{0}{\overset{(0)}{\varphi }},`$ (32)
$`\varphi _1=\underset{1}{\overset{(1)}{\varphi }}\mathrm{ln}v+\underset{1}{\overset{(0)}{\varphi }},`$ (33)
$`\varphi _2=\underset{2}{\overset{(1)}{\varphi }}\mathrm{ln}v+\underset{2}{\overset{(0)}{\varphi }},`$ (34)
where the coefficients $`\underset{n}{\overset{(1)}{\varphi }}`$ and $`\underset{n}{\overset{(0)}{\varphi }}`$ are completely logarithm free! Asymptotically one has that
$$\begin{array}{ccc}\{\begin{array}{c}\stackrel{~}{\stackrel{(1)}{\varphi }}_0=𝒪(\mathrm{\Omega }^3)\hfill \\ \stackrel{~}{\stackrel{(1)}{\varphi }}_1=𝒪(\mathrm{\Omega }^2)\hfill \\ \stackrel{~}{\stackrel{(1)}{\varphi }}_2=𝒪(\mathrm{\Omega }^2)\hfill \end{array}& \text{ or }& \{\begin{array}{c}\underset{0}{\overset{(0)}{\varphi }}=𝒪(1)\hfill \\ \underset{1}{\overset{(0)}{\varphi }}=𝒪(1)\hfill \\ \underset{2}{\overset{(0)}{\varphi }}=𝒪(\mathrm{\Omega })\hfill \end{array}\end{array}$$
(35)
and that
$$\begin{array}{ccc}\{\begin{array}{c}\stackrel{~}{\stackrel{(0)}{\varphi }}_0=𝒪(\mathrm{\Omega }^2)\hfill \\ \stackrel{~}{\stackrel{(0)}{\varphi }}_1=𝒪(\mathrm{\Omega }^2)\hfill \\ \stackrel{~}{\stackrel{(0)}{\varphi }}_2=𝒪(\mathrm{\Omega }^1)\hfill \end{array}& \text{ or }& \{\begin{array}{c}\underset{0}{\overset{(0)}{\varphi }}=𝒪(\mathrm{\Omega }^1)\hfill \\ \underset{1}{\overset{(0)}{\varphi }}=𝒪(1)\hfill \\ \underset{2}{\overset{(0)}{\varphi }}=𝒪(1)\hfill \end{array}\end{array}$$
(36)
for the other auxiliary field.
Thus, upon substitution on the field equations (20)-(22), one obtains the following equations for the *auxiliary field*,
$`\mathrm{\Delta }\underset{0}{\overset{(1)}{\varphi }}\delta \underset{1}{\overset{(1)}{\varphi }}=(2\gamma \mu )\underset{0}{\overset{(1)}{\varphi }}2\tau \underset{1}{\overset{(1)}{\varphi }}+\sigma \underset{2}{\overset{(1)}{\varphi }},`$ (37)
$`\mathrm{\Delta }\underset{1}{\overset{(1)}{\varphi }}\delta \underset{2}{\overset{(1)}{\varphi }}=2\mu \underset{1}{\overset{(1)}{\varphi }}+(2\beta \tau )\underset{2}{\overset{(1)}{\varphi }},`$ (38)
$`D\underset{2}{\overset{(1)}{\varphi }}\overline{\delta }\underset{1}{\overset{(1)}{\varphi }}=\lambda \underset{0}{\overset{(1)}{\varphi }}+2\pi \underset{1}{\overset{(1)}{\varphi }}+\rho \underset{2}{\overset{(1)}{\varphi }},`$ (39)
$`D\underset{1}{\overset{(1)}{\varphi }}\overline{\delta }\underset{0}{\overset{(1)}{\varphi }}=(\pi 2\alpha )\underset{0}{\overset{(1)}{\varphi }}+2\rho \underset{1}{\overset{(1)}{\varphi }}.`$ (40)
which are in fact identical to the spin-1 zero-rest-mass field equations (20)-(23). Again, the system is overdetermined, and thus one discards the last equation.
In an ordinary asymptotic characteristic initial value problem, the initial data over $`𝒩_0`$ is prescribed, that is
$$\varphi _0=\varphi _0^{4,0}v^1+\left(\varphi _0^{5,1}\mathrm{ln}v+\varphi _0^{5,0}\right)+\mathrm{},$$
(41)
is given. Thus, the value of
$$\underset{0}{\overset{(1)}{\varphi }}=\varphi _0^{5,1}+\mathrm{},$$
(42)
at $`𝒩_0`$ is also known. Furthermore, as seen before $`\varphi _1^{2,1}=\overline{ð}\varphi _0^{2,0}`$, and therefore we also know the value of $`\underset{1}{\overset{(1)}{\varphi }}`$ at $`𝒵_0`$. Finally, as discussed previously,$`\underset{2}{\overset{(1)}{\varphi }}=𝒪(v)`$, and therefore $`\underset{2}{\overset{(1)}{\varphi }}=0`$ on $`I^+`$, i.e. the field is non-radiative <sup>2</sup><sup>2</sup>2This means essentially that the radiation field component $`\varphi _0`$ has no type N term. Thus, it does not contribute to the Poynting vector and consequently no energy loss can be ascribed to it! . Consequently, we are are in the possession of all the ingredients needed to set an asymptotic characteristic value problem for the spin-1 zero-rest-mass field $`\underset{AB}{\overset{(1)}{\varphi }}`$. The initial data has been constructed out of the (non-regular) initial data for the original spin-1 zero-rest-mass field $`\varphi _{AB}`$. Duff and Friedrich’s techniques allows us to state that there is an (unique) analytic solution for this initial value problem in a neighbourhood of $`𝒵_0`$.
What happens with the other auxiliary field? In this case the equations are,
$`\mathrm{\Delta }\underset{0}{\overset{(0)}{\varphi }}\delta \underset{1}{\overset{(0)}{\varphi }}=(2\gamma \mu )\underset{0}{\overset{(0)}{\varphi }}2\tau \underset{1}{\overset{(0)}{\varphi }}+\sigma \underset{2}{\overset{(0)}{\varphi }},`$ (43)
$`\mathrm{\Delta }\underset{1}{\overset{(0)}{\varphi }}\delta \underset{2}{\overset{(0)}{\varphi }}=2\mu \underset{1}{\overset{(0)}{\varphi }}+(2\beta \tau )\underset{2}{\overset{(0)}{\varphi }},`$ (44)
$`D\underset{2}{\overset{(0)}{\varphi }}\overline{\delta }\underset{1}{\overset{(0)}{\varphi }}+v^1Q\underset{2}{\overset{(1)}{\varphi }}=\lambda \underset{0}{\overset{(0)}{\varphi }}+2\pi \underset{1}{\overset{(0)}{\varphi }}+\rho \underset{2}{\overset{(0)}{\varphi }},`$ (45)
$`D\underset{1}{\overset{(0)}{\varphi }}\overline{\delta }\underset{0}{\overset{(0)}{\varphi }}+v^1Q\underset{1}{\overset{(1)}{\varphi }}=(\pi 2\alpha )\underset{0}{\overset{(0)}{\varphi }}+2\rho \underset{1}{\overset{(0)}{\varphi }}.`$ (46)
Note the presence of the extra terms $`v^1Q\underset{2}{\overset{(1)}{\varphi }}`$ and $`v^1Q\underset{1}{\overset{(1)}{\varphi }}`$. Thus the auxiliary field $`\underset{AB}{\overset{(0)}{\varphi }}`$ satisfies a kind of sourced Maxwell equations. The source happens to be precisely the vacuum field $`\underset{AB}{\overset{(1)}{\varphi }}`$. As it has been done before, one discards the last equation (46).
From a previous analysis, one knows that $`\underset{2}{\overset{(1)}{\varphi }}=𝒪(v)`$, and thus the whole new term in equation (45) is such that $`v^1Q\underset{2}{\overset{(1)}{\varphi }}=𝒪(1)`$. Hence, it is regular at $`I^+`$. However, Duff and Friedrich’s techniques cannot be applied to the reduced system (43)-(45), as it is, for the initial data over the initial hypersurface $`𝒩_0`$ is not regular.
As discussed, $`\underset{0}{\overset{(0)}{\varphi }}=\varphi _0^{2,0}v^1+𝒪(1)`$. Substitution of this into equation (43) yields,
$$_u\varphi _0^{2,0}=0.$$
(47)
Therefore, the non-regular term is a constant of motion . Let $`\underset{0}{\overset{(0)}{\varphi }}=\varphi _0^{2,0}v^1+\underset{0}{\overset{(0)}{\varphi }}`$, with $`\underset{0}{\overset{(0)}{\varphi }}=𝒪(1)`$, i.e. from a physical spacetime point of view the component $`\underset{0}{\overset{(0)}{\varphi }}=𝒪(1)`$ satisfies the peeling behaviour. Substitution into equations (20)-(22) yields,
$`\mathrm{\Delta }\underset{0}{\overset{(0)}{\varphi }}\delta \underset{1}{\overset{(0)}{\varphi }}=(2\gamma \mu )\varphi _0^{2,0}v^1+(2\gamma \mu )\underset{0}{\overset{(0)}{\varphi }}2\tau \underset{1}{\overset{(0)}{\varphi }}+\sigma \underset{2}{\overset{(0)}{\varphi }},`$ (48)
$`\mathrm{\Delta }\underset{1}{\overset{(0)}{\varphi }}\delta \underset{2}{\overset{(0)}{\varphi }}=2\mu \underset{1}{\overset{(0)}{\varphi }}+(2\beta \tau )\underset{2}{\overset{(0)}{\varphi }},`$ (49)
$`D\underset{2}{\overset{(0)}{\varphi }}\overline{\delta }\underset{1}{\overset{(0)}{\varphi }}+v^1Q\underset{2}{\overset{(1)}{\varphi }}=\lambda \varphi _0^{2,0}v^1\lambda \underset{0}{\overset{(0)}{\varphi }}+2\pi \underset{1}{\overset{(0)}{\varphi }}+\rho \underset{2}{\overset{(0)}{\varphi }}.`$ (50)
Now, $`\gamma =𝒪(v)`$, $`\mu =𝒪(v)`$, and $`\lambda =𝒪(v)`$. Therefore, the new extra terms $`(2\gamma \mu )\varphi _0^{2,0}v^1`$ and $`\lambda \varphi _0^{2,0}v^1`$ are regular at $`I^+`$ and known. Moreover, they are also analytic at $`v=0`$. Therefore we have obtained a perfectly regular symmetric hyperbolic system for $`\underset{0}{\overset{(0)}{\varphi }}`$, $`\underset{1}{\overset{(0)}{\varphi }}`$, and $`\underset{2}{\overset{(0)}{\varphi }}`$.
Furthermore, we also know the value of $`\underset{0}{\overset{(0)}{\varphi }}`$ at the initial hypersurface $`𝒩_0`$ (obtainable from the non-regular initial value problem for $`\varphi _{AB}`$). The value of $`\underset{1}{\overset{(0)}{\varphi }}`$ at $`𝒵_0`$ is also known. Finally $`\underset{2}{\overset{(0)}{\varphi }}\varphi _2`$, and thus we also know the data on $`I^+`$. Consequently we have obtained a regular system of symmetric hyperbolic equations for the scalars $`\underset{0}{\overset{(0)}{\varphi }}`$, $`\underset{1}{\overset{(0)}{\varphi }}`$, $`\underset{2}{\overset{(0)}{\varphi }}`$ with suitable initial data on which Duff and Friedrich’s ideas can be used. The conclusion is that there exists a neighbourhood of $`𝒵_0`$ on which there exists an (unique) analytic solution ($`C^\omega `$) to the posed initial value problem. If $`\underset{0}{\overset{(0)}{\varphi }}`$ is analytic in a neighbourhood of $`I^+`$ then clearly $`\varphi _0=\varphi _0^{2,0}v^1`$ has a convergent Laurent expansion on the same neighbourhood.
So, we have seen that there exists a neighbourhood $`U_1`$ of $`𝒵_0`$ for which the field $`\underset{AB}{\overset{(1)}{\varphi }}`$ is analytic (in $`v`$), and another neighbourhood $`U_0`$ for which the field $`\underset{AB}{\overset{(0)}{\varphi }}`$ has a Laurent expansion (again in $`v`$). Let $`U=U_0U_1`$. Thus, one can say that the non-regular asymptotic characteristic initial value problem for the zero-rest-mass field $`\varphi _{AB}=\underset{AB}{\overset{(0)}{\varphi }}+\underset{AB}{\overset{(1)}{\varphi }}\mathrm{ln}v`$ has a (unique) solution in $`UI^+`$ which is in fact polyhomogeneous. Note that althought the field is not well defined at $`I^+`$, the radiation field is well defined ($`\varphi _2`$) as it is part of the regular initial data. Furthermore, it can be shown that the net electromagnetic charge of the field (as measured from null infinity) is also well defined.
These results can be sumarised in the form of the following
###### Theorem 1
Let $`(\stackrel{~}{M},\stackrel{~}{g}_{ab})`$ be a Ricci flat ($`R_{ab}=0`$) asymptotically simple spacetime, and let $`(M,g_{ab})`$ be the corresponding unphysical spacetime obtained by the conformal rescaling. Furthermore, let $`𝒩_0`$ be a future oriented null hypersurface in $`M`$ that intersects $`I^+`$ at $`𝒵_0`$. Then there exist a neighbourhood $`U`$ of $`𝒵_0`$ such that the non-regular asymptotic characteristic initial value problem for a spin-1 zero-rest-mass field with minimally polyhomogeneous data on $`𝒩_0`$ has a unique (polyhomogeneous) solution in $`UI^+`$.
## 4 Conclusions
An existence result for the asymptotic characteristic initial value problem for a spin-1 zero-rest-mass field with polyhomogeneous initial data has been proved. The polyhomogeneous initial data is not bounded at $`I^+`$, and thus, the initial value problem is non-regular. This means essentially that standard existence theorems for symmetric hyperbolic equations cannot be directly used. However, as seen it possible to reformulate this non-regular initial value problem as two coupled initial value problems for some auxiliary fields. These auxiliary fields satisfy partial differential equations with regular initial data. Moreover, these auxiliary systems have the same principal part (symbol) as the spin-1 zero-rest-mass field.
The existence result shows that the polyhomogeneous structure of the initial data is preserved at least in a neighbourhood of $`𝒵_0`$. And for the same price, one also settles the matter of the convergence of polyhomogeneous expansions for zero-rest-mass fields.
### 4.1 Further extensions
The proof given here is limited to minimally polyhomogeneous spin-1 zero-rest-mass field equations. A proof for an arbitrary polyhomogeneous spin-2 zero-rest-mass field (higher spins are not so interesting as the Buchdahl constraint has to be satisfied) follows using the same techniques. Arbitrary polyhomogeneity in the initial data means that further auxiliary fields and their respective systems should be included in the discussion. As it has been pointed out elsewhere, these auxiliary fields form a hierarchy. This hierarchy is such that the field equations for a given auxiliary field has the auxiliary field just above in the hierarchy as source term (cfr. with the equations for $`\underset{AB}{\overset{(1)}{\varphi }}`$ and $`\underset{AB}{\overset{(0)}{\varphi }}`$). The consideration of a spin-2 field rather than a spin-1 one would only increase the amount of algebra involved.
A much more interesting and natural extension of this work is to trying to apply similar techniques to polyhomogeneous gravitational fields. This would require the use of Friedrich’s regular conformal field equations, and a good understanding of the formal polyhomogeneous solutions to the aforementioned equations. These concerns will be and are the concern of future work.
## Acknowledgements
I want to thank my supervisor Prof. M.A.H. MacCallum for his advice and encouragement, and for suggesting me to work on the topics discussed here and in previous articles. I hold a scholarship (110441/110491), from the Consejo Nacional de Ciencia y Tecnología (CONACYT), Mexico.
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# The INT Search for Metal-Poor Stars. Spectroscopic Observations and Classification via Artificial Neural Networks1footnote 11footnote 1Based on observations made with INT operated on the island of La Palma by the Isaac Newton Group in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofisica de Canarias.
## 1 Introduction
High-resolution, high-signal-to-noise stellar spectroscopy provides a wealth of information on the physics and chemistry of stellar atmospheres, as well as on the dynamical and chemical evolution of our Galaxy. In the past, however, the small number of targets has limited the depth of information useful for constraining several important problems of the early history of the Galaxy. For example, measurements of chemical abundances and stellar motions of a sample of thousands of metal-deficient stars, ideally selected without kinematic bias, is required to address the most important problems relating to the kinematic structure of the stellar populations of the Milky Way (see, e.g., Beers et al. 2000; Chiba & Beers 2000). Similarly, the identification of trends, and in particular, the measurement of spreads in the relative abundances of individual elements provides required constraints on the nature of what was most likely a very inhomogeneous early evolution history (Ryan, Norris, & Beers 1996; Tsujimoto & Shigeyama 1998; Argast et al. 2000). To achieve these goals, large collaborative efforts are already in progress (see Beers 1999 for a summary).
Here we report on medium-resolution ($`2`$ Å ) spectroscopic observations for 731 stars obtained at the Isaac Newton Telescope (INT) on La Palma. The remaining $`400`$ stars with spectroscopy obtained during the INT follow-up observations turned out to be generally hotter than the main-sequence turnoff of an old stellar population, and are dominated by field horizontal-branch stars, blue metal-poor stars or blue stragglers, and other hot stars such as O- and B-type subdwarfs. These hot stars provide additional information on the kinematics of the Galaxy (e.g., Sommer-Larsen et al. 1997; Wilhelm, Beers, & Gray 1999), and will be the subject of a future paper.
The stars under discussion were previously identified as metal-poor candidates from the HK objective-prism/interference-filter survey of Beers and collaborators (as described by Beers, Preston, & Shectman 1985; Beers, Preston, & Shectman 1992, hereafter BPS). The HK objective-prism technique obtains wide-field (5 $`\times `$ 5 degree) photographic plates of widened spectra, with the interference filter selecting a $`150`$ Å region around the Ca ii H and K absorption lines, reaching limiting magnitudes of $`B1515.5`$. Initially, the HK survey made use of the Curtis Schmidt telescope at Cerro Tololo Inter-American Observatory for a survey of some 4000 square degrees of the southern-hemisphere sky. Later, the survey was extended to include some 3000 square degrees in the northern hemisphere, using the Burrell Schmidt telescope at Kitt Peak National Observatory. Visual inspection with a low-power binocular telescope was originally carried out by Preston, and later by Beers, to produce a list of $``$ 10,000 candidate metal-poor stars, a subset of which was observed with the INT. Other follow-up photometry and spectroscopy campaigns have been conducted in the northern and southern hemispheres; papers describing their results have either already been published (Schuster et al. 1996; Norris, Ryan, & Beers 1999; Anthony-Twarog et al. 2000) or are in preparation.
Beers et al. (1990) originally reported on a method for the estimation of stellar metal abundance based on spectra of similar resolution to our present program. The Beers et al. procedure made use of synthetic spectra generated from the Kurucz model atmospheres (Kurucz 1993) to derive a relationship between the metallicity \[Fe/H\], the pseudo-equivalent width of the Ca ii K line, and predicted broadband $`(BV)_0`$ color, as obtained from the Revised Yale Isochrones (Green 1988; King, Demarque, & Green 1988). This approach employed interpolated polynomial fits in the Ca ii K vs $`(BV)_0`$ plane, forcing agreement between the metallicities and colors provided by the models, and those derived from high-resolution spectroscopy and broadband photometry for a set of standard stars covering the range of metallicities and colors of interest. This method has recently been refined (Beers et al. 1999) by making use of locally weighted multi-dimensional regressions instead of polynomial fits, and by greatly expanding the number of standards used in the calibration. The 551 standard stars with accurate high-resolution metallicity determinations and broadband photometry discussed by Beers et al. (1999) are used in the present paper to train and test the performance of the ANNs. No overlap exists between the standard stars and the 731 program stars.
The outline of this paper is as follows. In §2 we summarize the spectroscopic observations and data reduction, as well as the measurement of radial velocities and line indices for the stars. In §3 we discuss the use of the ANN technique for the estimation of stellar metallicity and broadband colors for the INT sample, and compare with preliminary results derived from application of the Beers et al. (1999) approach. In §4 we discuss the results of our search for low metallicity stars, and compare with the previous searches of BPS. Section 5 presents our conclusions.
## 2 Observations, Reductions, and Measurement of Radial Velocities and Line Indices
Spectroscopic observations of our sample were carried out with the 2.5-m INT at the Observatorio del Roque de los Muchachos, on the island of La Palma (Spain), during three runs between 1995 March and 1996 June (see Table 1). The 235 mm camera and the 1200 groove mm<sup>-1</sup> grating of the Intermediate Dispersion Spectrograph (IDS) provided a resolving power $`\lambda /\mathrm{\Delta }\lambda `$ 2000, and a spectral range of 1100 Å centered at 4150 Å, as shown in the sample spectrum in Figure 1. The $`1024\times 1024`$ TEK3 CCD camera that was used is roughly 60% efficient over the observed spectral range. Only 256 pixels were retained in the spatial direction, reducing the readout time to about 40 seconds while still providing a sufficient range for optimal sky subtraction. The cosmetics on this chip were excellent, with a typical pixel-to-pixel sensitivity variation at a level of 1%. Our target signal-to-noise ratio was $`S/N>20`$, and required exposure times, under conditions of good seeing, varied between 1 and 30 minutes for stars in the apparent magnitude range 11 $`V`$ 15.
Data reduction was performed in the standard way within the IRAF<sup>2</sup><sup>2</sup>2IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation. environment, and consisted of bias correction, flatfielding, variance-weighted spectral extraction, and wavelength calibration.
### 2.1 Radial Velocity Determinations
The IDS is attached to the INT’s Cassegrain focus, resulting in displacements of the individual spectra due to flexure during the course of the night that might enlarge the errors in the measured radial velocities. To reduce this source of uncertainty, we obtained calibration-lamp spectra for groups of 5–10 stellar spectra of targets selected from the same HK survey plate, and which were therefore located in the same region of the sky. One calibration lamp spectrum was obtained at the end of the exposure of the last star in the group, before pointing to the next star, and another was taken before starting on the new group, with the telescope already pointing to the new target. In this way we collected two calibration spectra for every group, as well as spectra for two or three radial-velocity standards per night. The geocentric radial velocity for each star was measured by two methods: averaging the spectral shifts measured in prominent absorption lines, and by cross-correlating the spectra with a set of artificial templates and standard stars (see Beers et al. 1999, and references therein, for a complete description). Corrections for the Earth’s rotation and orbital motion were then applied to take the velocities to the heliocentric system. We estimate a final accuracy for the radial velocities of about 10 km s<sup>-1</sup>, as verified by comparison with the radial velocity standards.
### 2.2 Line Index Measurement
The measured geocentric radial velocities were used to place the wavelength scale of the spectra at zero rest velocity. We then determined pseudo-equivalent widths for a number of prominent spectral features, summing the total counts around the required line, and using a straight average of the flux in nearby side-bands to estimate the continuum level. In some cases, several wavelength intervals were used, as previous work has been shown that the optimal band size depends on line strength. Table 2 lists the bands employed. KP uses the indices K6, K12, and K18; HP2 uses the indices HD12 and HD24 (see Beers et al. 1999 for more details). KP is sensitive to the stellar calcium abundance, which is assumed to track the overall metallicity (as quantified by \[Fe/H\]<sup>3</sup><sup>3</sup>3$`[A/B]=\mathrm{log}\left(\frac{N_A}{N_B}\right)\mathrm{log}\left(\frac{N_A}{N_B}\right)_{}`$, where $`N_X`$ and is the number density of nuclei of an element X.), but it also depends on effective temperature, and more weakly, on surface gravity. Figure 2 illustrates the relationship between an independent estimate of \[Fe/H\] from high-resolution spectroscopic analyses in the literature, \[Fe/H\]<sub>lit</sub>, as a function of the KP index for the 551 standard stars of Beers et al. used in this study. The observed scatter about the simple linear fit arises primarily from real differences in the effective temperatures of the stars in the sample.
There are well-documented differences between the solar elemental abundance ratios and those found in the photospheres of most metal-poor stars. In particular, the ratio of calcium-to-iron abundance is known to vary with metallicity (see, e.g., McWilliam 1997). The procedure we have employed to assign metallicity estimates remains valid, however, since the overabundance of calcium with respect to iron increases monotonically and with little scatter towards lower iron abundances for the vast majority of stars, and the empirical metallicity calibration described below takes this variation into account. The critical pair of indices is completed with the HP2 or the H8 index, which provide estimates of the strengths of the Balmer lines H-$`\delta `$ or H-8, respectively. The well-known temperature sensitivity of the hydrogen lines (see, e.g., Fuhrmann, Axer, & Gehren 1993) is captured by these indices. Figure 3 illustrates the relationship between another temperature-sensitive quantity, the broadband $`(BV)_0`$ color, with HP2, using the same set of standard stars from Beers et al. (1999). In the present application, we explored supplementing these temperature-sensitive indices with two others, centered on the Ca i $`\lambda `$ 4226 Å (the CAP index) and the CH G-band at 4300 Å (the GP index). The additional spectral information allowed us to improve estimates of the metallicities and colors for the program stars, and might be used in the future to measure other quantities, such as carbon abundance.
A number of white dwarfs, hot O- and B-type subdwarfs, other stars bluer than $`BV0.3`$, including numerous field horizontal-branch and other A-type stars, as well as a number of emission-line objects, mostly M-dwarfs with active photospheres, were mistaken for metal-poor stars in the (visual) identification of candidates from the HK objective-prism survey. To exclude such objects, only the stars having metallic-line indices in the range of sensitivity to \[Fe/H\], and with estimated colors in the range $`0.3BV1.2`$ ($`0.1`$ HP2 $`6.0`$), were extracted from the full set of spectra, and are the subject of this paper.
Column (1) of Table 3 lists the star names of our sample. Columns (2) and (3) list the (1950.0 epoch) coordinates, generally accurate to a few arcseconds. The Galactic longitude and latitude of the stars are provided in columns (4) and (5), respectively. Column (6) lists the measured heliocentric radial velocities, in units of km s<sup>-1</sup>. Column (7) lists an estimate of the Johnson $`V`$-band apparent magnitude, when available in the Guide Star Catalogue (Lasker et al. 1989). Column (8) is the estimated de-reddened $`(BV)_0`$ color, listed as $`BV`$ (described below), in order to distinguish this estimate from a true photometric measurement. Columns (9)–(13) list the observed line indices for the program stars in units of Å. Column (14) lists the derived estimate of \[Fe/H\], as described below.
For a number of candidates, the presence of two or more objects in the field of view caused difficulties for the identification of the HK survey target. In many such situations, it was decided to put multiple candidates on the slit (by rotating the spectrograph) and obtain spectra of all possible sources. These multiple cases have been labeled in Table 3 according to their approximate orientation on the sky.
A subset of the plates from the HK survey overlapped with one another, hence there are occasional instances of multiple identifications. Table 4 lists the names of the multiply identified stars that were included among the INT program sample. A search of the SIMBAD catalog within 1 arcmin of the listed positions of our program stars revealed a small number of likely previous identifications; these are listed in Table 5.
Figure 4 shows the positions of the program stars in Galactic coordinates, the upper and lower panels corresponding to the northern and southern hemispheres, respectively.
## 3 Classification of the Stars
Previous application of ANNs to the classification of astronomical spectra has demonstrated the great utility of this new tool. For example, ANNs have been applied to low-dispersion IUE spectra by Vieira & Ponz (1995), allowing the determination of spectral class to an accuracy of $``$1.1 subclasses. More recently, Bailer-Jones, Irwin & von Hippel (1998) used ANNs to classify spectra of resolution $``$ 3 Å extracted from the Michigan Spectral Survey (Houk 1994) with an accuracy of 1.09 spectral subclasses, and obtained correct luminosity classes for over 95% of dwarfs and giants in their test sample.
An optimum procedure would make use of the entire measured spectrum of a given star as an input for an ANN, and ideally, provide a classification method to extract the effective temperature, surface gravity, and metallicity of each star (see, e.g., Snider et al. 2000). Unfortunately, the full implementation of this approach was not available at the beginning of our analysis, thus we decided to adopt a previous step, one making use of the measured spectral indices described above. This classification method, based on limited information input to the ANN, has nevertheless enabled reasonably accurate estimates to be made of the $`(BV)_0`$ colors and \[Fe/H\] for the sample of metal-poor stars observed at the INT, and is described below. It is worthwhile noting that Bailer-Jones (2000) has made use of synthetic spectra to show that, when spectrophotometry is available, i.e., the absolute energy distribution is measured, ANNs fed with the full spectrum behave quite well in the face of degradations in the signal-to-noise ratio and the resolving power.
The Beers et al. (1999) calibration has been applied to provide preliminary estimates of \[Fe/H\] for the sample of metal-poor stars observed at the INT, and serves as an independent reference to compare with the metallicities derived from the ANN approach. This calibration is based on the combination of the KP index, and an Auto-Correlation Function (ACF) of the stellar spectrum, which produces abundance estimates having external errors on the order of 0.15–0.20 dex over nearly the complete range of metallicities expected in the stellar populations of the Galaxy, $`4.0[\mathrm{Fe}/\mathrm{H}]0.0`$. The ACF employed by Beers et al. takes advantage of the fact that, when a stellar spectrum is correlated with itself, the peak value of the resulting (non-trivial overlapping) function depends on the strengths of various metallic features in the spectrum “beating” against one another. The greater the numbers of such lines, and the greater their strengths, the larger the signal. As Beers et al. discuss, the ACF approach provides complimentary information to the KP index methodology, and performs particularly well for cooler or more metal-rich stars where the KP index approaches saturation.
There are a number of iterative steps that are applied in order to arrive at a final estimate of stellar metallicity using the Beers et al. techniques. In the case of the present application, we do not yet have measured ACFs for all of the INT program stars, so we base our comparison spectroscopic abundance estimates solely on the KP index. The estimate we make use of, \[Fe/H\]<sub>K3</sub> (adopting the Beers et al. nomenclature), is expected to be close to, but not exactly the same as, the metallicity estimate that will be obtained once the full set of ACFs is available. The largest differences are expected for the stars with metallicities greater than \[Fe/H\] $`1.0`$, where the KP index begins to saturate, especially for the cooler stars. An approximate correction for this effect, even in the absence of ACF measurements, has been applied before making our comparisons. Beers et al. also describe an independent means of predicting broadband $`(BV)_0`$ colors, based on the Balmer line indices, which is used as a check on the ANN’s assignment of colors.
### 3.1 Artificial Neural Networks
Artificial neural networks are computational systems whose structure is based on that of the human brain and nervous system. In the present application we employ a feed-forward neural network, where “neurons” are grouped into layers (Murtagh 1991). There is one input layer where the information comes in, one or several hidden layers where the information is processed, and one layer that yields the output of the network computations. For instance, a 2-5-6-1 ANN architecture is comprised of a network with two input neurons, one output neuron, and two hidden layers with five and six neurons each. The neural network provides a mapping, $`f`$, from a set of inputs to a set of desired outputs. The set of “input”–“desired output” pairs forms the pattern sample. Depending on how the “desired output” is defined, learning will be supervised (Serra-Ricart et al. 1996) or unsupervised (Serra-Ricart et al. 1993). The mapping, $`f`$, is approximated by application of a specific learning algorithm (in this case, a back-propagation algorithm – see Rumelhart, Hinton, & Williams 1986) to adjust the weights of connections from examples contained in the pattern sample.
In the present case the pattern sample was the set of line indices listed in Table 2 for a randomly chosen subset of 50% of the standard stars described in Beers et al. (1999). The other half of the sample is set aside to carry out an independent test of the accuracy of the ANN. The list of standards includes 551 stars with high-resolution determinations of \[Fe/H\], and previously measured $`(BV)_0`$ colors, which had medium-resolution spectroscopy obtained with nine different telescope/spectrograph combinations at essentially the same resolution as the INT follow-up spectra, over the course of the HK survey. The heterogeneity of the sample does not necessarily introduce significant additional random error, as it was checked that variations of the KP and HP/H8 indices measured at different observatories were at the level of 5 % and 10 %, respectively, roughly the expected measurement error for a single spectrum with our minimum signal-to-noise ratio ($`S/N=20`$).
The two primary difficulties generally encountered with presently studied ANN classifiers are:
1. Poor mapping from input examples of the $`f`$ function. In Figure 5 the standard stars are plotted (with open circles) in the KP–HP2 plane together with the INT program stars (filled circles), demonstrating that the standards essentially map the input parameter space. It is possible to distinguish some stellar classes in the KP–HP2 diagram, as can be observed in Figure 6.
2. The ANN becomes “over-trained.” This means that the ANN has fully matched the input data. A large ANN can learn the entire pattern sample, essentially memorizing the complete set of input examples instead of providing a valid approximation of the mapping function $`f`$. To prevent over-training we randomly separate the pattern sample into two groups of almost the same size, then use one of them for the training task (the training sample), and the other to test for convergence of the training task (the testing sample).
Some of the spectral features quantified by the line indices were not measurable in all the stars, for example, the HP2 index for the cooler stars. Hence we considered a number of different ANNs, depending on the available indices. This scheme was also adopted in order to identify the ANN with the highest performance for a given set of available indices and output parameters. We trained and ran 49 different configurations with some or all of the input indices described in Table 2 (KP, HP2, H8, CAP, and GP) to extract some or all of the following parameters: the estimated $`(BV)_0`$ color index, \[Fe/H\], surface gravity, and the related parameter, luminosity. The experiments showed that the ANN was able to recover estimates of $`(BV)_0`$ and \[Fe/H\] with reasonable accuracy, but that surface gravity and luminosity remain a challenge. The ANN configurations that lacked at least the KP and one of the two hydrogen indices (HP2 and H8) were shown to be unsatisfactory and were discarded. Finally, two different configurations were used to provide estimates of the $`(BV)_0`$ color index, and three different configurations were used to provide estimates of metallicity. Most of the stars in the INT sample were classified by the optimal ANN configurations (#4 and #11), although the differences with the other configurations considered, listed in Table 6, were not large.
Figure 7 shows the results obtained for the testing sample of standard stars after training with ANN configuration #11. The errors in metallicity exhibit a Gaussian distribution, and do not depend significantly on \[Fe/H\]. Likewise, Figure 8 shows that the ANN estimate of $`(BV)_0`$, $`BV`$, correlates well with the measured broadband color. Note, however, that the errors on estimates of color are much smaller for the hotter (smaller $`BV`$) stars than for the cooler (larger $`BV`$) stars, for which the error distribution deviates from a Gaussian distribution. This is a natural consequence of the Balmer lines’ weakening towards lower effective temperatures, and the resulting increase in the errors of their line index measurements. In order to check for possible correlations between errors in \[Fe/H\] and $`BV`$ we made use of one of the ANN configurations with two output neurons (config. #20). Although the single-output configurations were preferred for processing of the INT stars, the uncertainty in the metallicity derived from this configuration is similar to that achieved with the preferred configuration, #11, as visual comparison of Figure 9a with Figure 7a confirms. No correlation between the errors in \[Fe/H\] and the color index were found, as demonstrated in Figure 9b.
The quoted $`\sigma _{\mathrm{rms}}`$ value for the $`[\mathrm{Fe}/\mathrm{H}]_{\mathrm{ANN}}`$ estimate from the analyses of the standard stars included in the testing set, which were not previously seen by the ANN, is 0.32 dex. This value corresponds to configuration #11, which was used to classify most of the program stars, but the results for other configurations are very similar (see Table 6). The $`\sigma _{\mathrm{rms}}`$ values for $`BV`$ fall in the range 0.07–0.15, depending on the color itself.
There were 22 stars in the original INT sample with indices that were likely to have been affected by an underlying emission component in the spectral lines under consideration, as simple visual inspection of their spectra suggests. The metallicities derived for such stars using the ANN could be severely in error, usually underestimated, and hence these stars were excluded.
### 3.2 Comparison Between the ANN and Beers et al. (1999) Methods
These two methods quote similar errors for the estimated $`(BV)_0`$ colors of the standard stars. The ANN provides, using the KP and HP2 indices, $`\sigma _{\mathrm{rms}}`$ values for $`BV_{\mathrm{ANN}}`$ on the order of 0.07 for $`(BV)_00.8`$, and on the order of 0.15 for $`(BV)_0>0.8`$. The Beers et al. (1999) calibration of the HP2 index yields errors for $`BV_{\mathrm{HP2}}`$ on the order of 0.04 and 0.15 for these same color intervals, respectively. At least some of this error in predicted color, perhaps as much as half, may arise from uncertainties in the reddening corrections applied to the standard stars.
Comparing the \[Fe/H\] estimates, the Beers et al. (1999) method exhibits a nominal scatter of iron abundance estimates for the standard stars of 0.29 dex, while the ANN reflects a scatter (for configuration #11) of 0.32 dex. These error estimates include the contribution due to the error in the abundances of the standards themselves, which are expected to be on the order of 0.15–0.2 dex, so the intrinsic error of the techniques employed here is $``$ 0.2 dex.
Clearly, the formal quoted errors are quite similar for both methods. They share the same standard stars, but the ANN uses only half of them for the training sample. Figure 10 shows the comparison of the results obtained using both methods. The standard deviation of the differences for the estimated $`(BV)_0`$ colors is 0.03 magnitudes, while the same quantity for the differences in the predicted metallicities is 0.24 dex, well within the requirements for the agreement with the quoted errors for each method. There is a low-order “wiggle” in the residuals of the comparison of abundances and estimated colors obtained by the two methods, as can be seen in Figure 10, that is presently of unknown origin.
The ANN method is able to achieve an accuracy similar to the Beers et al. (1999) approach for estimates of $`(BV)_0`$ and \[Fe/H\]. Its character is completely empirical; the ANN is trained to reproduce the classification scheme employed for the standard stars. We expect that errors in the modeling of stellar atmospheres will affect both the calibrations, but the former one directly, for it is based on Kurucz’s models and empirical corrections, and the latter one only indirectly for the standards having been analyzed with classical model atmospheres. The main advantage of the ANN estimator is of a practical nature, as the effort required to tune a complex semi-empirical procedure such as that employed by Beers et al. is generally greater than that needed to properly train the ANNs.
## 4 The Effective Yield of the INT Metal-Poor Star Follow-Up
We have applied the trained ANNs to estimate \[Fe/H\] and $`(BV)_0`$ for the 731 late-type stars in the INT follow-up. The normalized metallicity distribution of the program stars, as derived from the ANN, is shown in Figure 11 (solid histogram). Comparison with the metallicity distribution of the metal-poor candidates from the southern HK survey observed by BPS (dashed histogram) indicates that the INT search has been far less effective in finding very metal-poor stars than the southern program. The peak of the metallicity distribution for stars studied by BPS is centered around \[Fe/H\] $`=2.2`$, while the INT distribution is centered at approximately \[Fe/H\] $`=0.8`$. There are several reasons for this difference. Firstly, the identification of candidates from the plates of the southern HK survey was, intentionally, much more selective than that applied to the plates of the northern hemisphere. In the north, this choice was made in order to obtain larger samples of the kinematically interesting metal-weak thick-disk stars. Secondly, the southern survey benefited from a partial photometric “pre-filtering” of candidates based on $`UBV`$ photometry – so that a greater fraction of the truly metal-poor stars were observed spectroscopically, and the clear “mistake” stars, with colors either too blue or too red to be of interest, were eliminated prior to spectroscopic follow-up. No such photometric pre-filtering step was taken for the INT program stars. Furthermore, the regions of high Galactic latitude (above $`|b|=60\mathrm{deg}`$) were given higher priority in the southern HK survey follow-up, but this was not so for the INT follow-up. Figure 12 also shows that there also exists a tendency of the INT sample towards brighter stars ($``$V(INT)$``$=13.5 $`\pm `$ 0.8 mag, while $``$V(BPS)$``$=14.2 $`\pm `$ 0.6 mag), which would emphasize thick disk stars relative to the generally fainter halo objects.
As a result of the above, a large proportion of the INT program stars of intermediate metallicity are expected to belong to the thick-disk population, while only the stars in the tail of the distribution at very low metallicity are expected to be dominated by members of the Galactic halo population. Figure 13, a comparison of the distribution of heliocentric radial velocities for the INT sample (solid histogram) with the southern stars in the BPS sample (dashed histogram), also bears this out.
The 731 stars analyzed in this work were identified from among the 1203 objects observed spectroscopically with the INT. The survey effort has identified 195 stars with \[Fe/H\] $`1.0`$, 67 stars with \[Fe/H\] $`2.0`$, and 12 new stars with \[Fe/H\] $`3.0`$. Thus, using the definition of Beers (2000), the effective yield (EY) of the INT spectroscopic follow-up for stars below \[Fe/H\] $`=2.0`$ is only on the order of 6%. Had a photometric pre-filtering been done first, at a minimum the EY could have been boosted to roughly 10%, by elimination of the stars that are either too hot or too cool to be of interest for this survey. If complete $`UBV`$ information were in hand before the spectroscopic follow-up commenced, the EY would likely be on the order of 20%-30%, thus commensurate with that found by other collaborations that did perform a photometric pre-filtering (see Table 1 of Beers 1999). Of course, the great advantage of proceeding with the spectroscopic follow-up first is that the observations were conducted over only two observing seasons, as opposed to many years of effort required for a complete photometric pre-filtering, combined with later spectroscopy.
## 5 Summary and Conclusions
We have carried out a search for metal-poor stars in the northern hemisphere using the Isaac Newton Telescope. Medium-resolution spectra of 1203 candidates identified from the HK survey were obtained, from which 731 stars were selected for further analysis. Artificial Neural Networks were trained and tested with previously available line indices for the large sample of standard stars described by Beers et al. (1999), then were used to obtain estimates of \[Fe/H\] and the intrinsic $`(BV)_0`$ color index for the INT program stars. This survey effort has identified 195 stars with \[Fe/H\] $`1.0`$, 67 stars with \[Fe/H\] $`2.0`$, and 12 new stars with \[Fe/H\] $`3.0`$. Although the EY of metal-poor stars obtained was very low, the number of metal-poor stars identified per observing hour was quite high, in a domain were multi-object spectrographs on similar size telescopes cannot presently compete due to the low areal density of targets, $`24`$ per square degree.
We conclude that ANNs are able to provide an accuracy of better than 0.3 dex in the determination of stellar metallicity, and roughly 0.1 magnitudes in estimates of broadband colors from intermediate-resolution spectra. These levels of accuracy are equivalent to those achieved by a much more elaborate methodology that makes use of model atmospheres and empirical corrections. ANNs can be trained rapidly, and can provide very fast classifications of stellar spectra. We emphasize that the procedure is entirely empirical, the errors being driven mainly by the accuracy and consistency with which we can provide metallicities for the training set of standard stars. From results previously obtained feeding the entire stellar spectrum to the ANNs (as opposed to line indices) for determination of spectral classes (Vieira & Ponz 1995; Bailer-Jones 1997; Bailer-Jones et al. 1998), and recent tests on samples spanning a large range of metallicities and temperatures (Qu et al. 1998; Snider et al. 2000), the application of the full-spectrum technique to carry out multi-dimensional classification of stars appears very promising.
We sincerely acknowledge the efficient and professional help of the staff at the Observatorio del Roque de los Muchachos, especially Palmira Arenaz, Marco Azzaro, Miriam Centurión, Carlos Martín, Neil O’Mahony, and Juerg Rey in the operation of the INT and the IDS spectrograph. David Lambert is thanked for his participation in one of the observing runs and for interesting discussions. We are indebted to the referee for suggestions that helped to clarify the presentation. This research has made use of the SIMBAD database, operated at CDS, Strasbourg (France), the NASA’s ADS Abstract Service, and has been partially funded by the Spanish DGES under projects PB95-1132-C02-01 and PB98-0531-C02-02. This work also received partial support from NATO grant 950875, and extensions. TCB acknowledges partial support from grant AST 95-29454 awarded by the US National Science Foundation. SR acknowledges partial support from the Brazilian agencies CNPq (grant 200068/95-4) and FAPESP (98/02706-6).
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# Fuchsian Affine Actions of Surface Groups
## 1 Introduction
Let $`\mathrm{\Gamma }`$ be the fundamental group of a compact surface. Let $`\lambda _n`$ be the standard $`n`$-dimensional representation of $`SL(2,)`$ in $`SL(n,)`$. We shall say a representation $`\rho `$ of $`\mathrm{\Gamma }`$ in $`SL(n,)`$ is Fuchsian ( or $`n`$-Fuchsian ) if $`\rho =\lambda _n\iota `$, where $`\iota `$ is a discrete faithful representation of $`\mathrm{\Gamma }`$ in $`SL(2,)`$. We shall also say by extension the image of $`\rho `$ is Fuchsian, and that an affine action of a surface group is Fuchsian, if its linear part is Fuchsian.
Our main result is the following theorem
###### Theorem 1.1
A finite dimensional affine Fuchsian action of the fundamental group of a compact surface is not proper.
In even dimensions, this is a trivial remark. For dimension $`4p+1`$, this theorem follows from previous results of G. Margulis concerning actions of free groups by using the Margulis invariant and lemma 4.1, also due to Margulis. Therefore, our proof shall concentrate on dimensions $`4p+3`$.
This case bears special features: one should notice that G. Margulis has exhibited proper actions of free group (with two generators) on $`^3`$ . Therefore, surface groups behave differently than free groups in these dimensions.
When $`dim(E)=3`$, our result is a theorem of G. Mess , for which G. Margulis and W. Goldman have obtained a different proof using Margulis invariant and Teichmüller theory. Our proof is based on similar ideas, but uses instead of Teichmüller theory a result on Anosov flows and a holomorphic interpretation of Margulis invariant, hence generalizing to higher dimensions.
It is a pleasure to thank M. Babillot, W. Goldman, G. Margulis for helpful conversations.
## 2 Representations of $`SL(2,)`$, surfaces and connections
In this section, we describe the irreducible representation of $`SL(2,)`$ of dimension $`2n+1`$ as the holonomy of a flat connection.
It is well known that in dimension 3, the irreducible representation of $`SL(2,)`$ is associated with the Minkowski model of the hyperbolic plane $`^2`$. More precisely, there exists a flat connection on $`E=T^2`$, such that the action of $`SL(2,)`$ lifts to a connection preserving action on this bundle. Hence, we obtain a 3-dimensional representation of $`SL(2,)`$. Furthermore, the Minkowski model is obtained using the section $`(1,0)`$ of $`E`$.
We will now be more precise and explain this construction in more details in higher dimensions.
### 2.1 A flat connection
Let $`^2`$ be the hyperbolic plane with its complex structure. Let $`L_k`$ be the complex line bundle over $`^2`$ defined by
$$L_k=(T^2)^_{}^k.$$
Let
$$E=L_1\mathrm{}L_n.$$
Notice now that $`SL(2,)`$ acts on all $`L_k`$, hence on $`E`$, by bundle automorphisms.
If $`Y`$ is a section of $`E`$, $`Y_0`$ will denote its component on the factor $``$, and $`Y_k`$ its component on $`L_k`$. The space of sections of the bundle $`V`$ will be denoted $`\mathrm{\Gamma }(V)`$. The metric on $`L_i`$, induced from the Riemanian metric on $`^2`$ will be denoted $`,`$. By definition, if $`YL_k`$, $`XL_1`$, then $`i_XY`$ is the element of $`L_{k1}`$ such that
$$ZL_k,i_XY,Z=Y,XZ.$$
Let $`\overline{}`$ be the Levi-Civita connection on $`L_1`$, and, by extension, the induced connection on $`L_k`$. We introduce the following connection $``$ on $`E`$, defined if $`XT^2`$, $`Y\mathrm{\Gamma }(E)`$ by
$$\{\begin{array}{ccc}\hfill (_XY)_0& =& L_XY_0+\frac{1}{2}(n+1)X,Y_1\hfill \\ \hfill k>0,(_XY)_k& =& (nk+1)XY_{k1}+\overline{}_XY_k\hfill \\ & & +\frac{1}{4}(n+k+1)i_XY_{k+1}.\hfill \end{array}$$
Consider the family or real numbers, defined for $`k\{0,n1\}`$,
$$a_0=1,a_{k+1}=\frac{1}{2^{2k+1}}\underset{j=0}{\overset{j=k}{}}(\frac{n+j+1}{nj}).$$
Define a metric of signature $`(n,n+1)`$ on $`E`$ by
$$Y,Z=\underset{k=0}{\overset{k=n}{}}(1)^{k+1}a_kY_k,Z_k.$$
The main result of this section is the following statement
###### Proposition 2.1
The connection $``$ is flat, and preserves the metric $`,`$. Furthermore, the $`SL(2,)`$ action on $`E`$ preserves the metric $`,`$ and the connection $``$. The resulting $`(2n+1)`$-representation of $`SL(2,)`$ is irreducible.
Proof: Let $`X`$, $`Z`$ two commuting vector fields on $`^2`$. Let $`\omega `$ the Kähler form of $`^2`$ defined by $`\omega (Z,X)=JZ,X`$. Let’s first introduce the following notation. If $`f`$ is a function of $`Z`$ and $`X`$ then
$$\underset{¯}{f(Z,X)}=f(Z,X)f(X,Z).$$
With these notations at hands, we have
$`\underset{¯}{ZX,Y}`$ $`=`$ $`\omega (Z,X)JY,`$
$`\underset{¯}{Zi_XY}`$ $`=`$ $`2\omega (Z,X)JY,`$
$`\underset{¯}{i_Z(XY)}`$ $`=`$ $`2\omega (Z,X)JY.`$
Let $`Y`$ be a section of $`E`$. Let $`\overline{R}`$ be the curvature tensor of $`\overline{}`$ and recall that
$$\overline{R}(Z,X)Y_k=k\omega (Z,X)JY.$$
We begin our computations. Let $`R`$ be the curvature tensor of $``$. We first have
$`(_Z_XY)_0=`$
$`L_ZL_XY_0+{\displaystyle \frac{1}{2}}(n+1)L_ZX,Y_1`$
$`+{\displaystyle \frac{1}{2}}(n+1)(n)Z,XY_0+{\displaystyle \frac{1}{2}}(n+1)Z,\overline{}_XY_1`$
$`+{\displaystyle \frac{1}{8}}(n+1)(n+2)Z,i_XY_2.`$
Hence
$$(R(Z,X)Y)_0=0.$$
Next
$`(_Z_XY)_1=`$
$`nZL_XY_0+{\displaystyle \frac{1}{2}}n(n+1)ZX,Y_1+n\overline{}_Z(XY_0)`$
$`+\overline{}_Z\overline{}_XY_1+{\displaystyle \frac{1}{4}}(n+2)\overline{}_Z(i_XY_2)+{\displaystyle \frac{1}{4}}(n+2)(n1)i_Z(XY_1)`$
$`+{\displaystyle \frac{1}{4}}(n+2)i_Z\overline{}_XY_2+{\displaystyle \frac{1}{16}}(n+2)(n+3)i_Zi_XY_3.`$
We get
$`(R(Z,X)Y)_1=`$
$`+{\displaystyle \frac{1}{2}}n(n+1)\underset{¯}{ZX,Y_1}+\overline{R}(Z,X)Y_1`$
$`+{\displaystyle \frac{1}{4}}(n+2)(n1)\underset{¯}{i_Z(XY_1)}+{\displaystyle \frac{1}{16}}(n+2)(n+3)\underset{¯}{i_Zi_XY_3}.`$
Hence
$`(R(Z,X)Y)_1`$
$`=\omega (Z,X)JY\left({\displaystyle \frac{1}{2}}n(n+1)+1+{\displaystyle \frac{2}{4}}(n+2)(n1)\right)`$
$`=0.`$
It remains to consider the case $`k>0`$
$`(_Z_XY)_k=`$
$`(nk+1)Z\overline{}_XY_{k1}+{\displaystyle \frac{1}{4}}(nk+1)(n+k)Zi_XY_k`$
$`+(nk+1)(nk+2)ZXY_{k2}+(nk+1)\overline{}_Z(XY_{k1})`$
$`+\overline{}_Z\overline{}_XY_k+{\displaystyle \frac{1}{4}}(n+k+1)\overline{}_Z(i_XY_{k+1})`$
$`+{\displaystyle \frac{1}{4}}(n+k+1)(nk)i_Z(XY_k)+{\displaystyle \frac{1}{4}}(n+k+1)i_Z\overline{}_XY_{k+1}`$
$`+{\displaystyle \frac{1}{16}}(n+k+1)(n+k+2)i_Zi_XY_{k+2}.`$
We get
$`(R(Z,X)Y)_k=`$
$`+{\displaystyle \frac{1}{4}}(nk+1)(n+k)\underset{¯}{Zi_XY_k}+\overline{R}(Z,X)Y_k`$
$`+{\displaystyle \frac{1}{4}}(n+k+1)(nk)\underset{¯}{i_Z(XY_k)}`$
$`+{\displaystyle \frac{1}{16}}(n+k+1)(n+k+2)\underset{¯}{i_Zi_XY_{k+2}}`$
$`+(nk+1)(nk+2)\underset{¯}{(ZXY_{k2})}.`$
Hence
$`(R(Z,X)Y)_k`$
$`=\omega (Z,X)JY\left({\displaystyle \frac{2}{4}}(nk+1)(n+k)+k+{\displaystyle \frac{2}{4}}(n+k+1)(nk)\right)`$
$`=0.`$
We have just proved the connection $`\overline{}`$ is flat. Now, we show $``$ preserves $`,`$. Let $`Y`$ a section of $`E`$. Then
$`_XY,Y`$
$`=`$ $`{\displaystyle \underset{k=0}{\overset{k=n}{}}}(1)^{k+1}a_k(_X)Y_k,Y_k`$
$`=`$ $`L_XY_0{\displaystyle \frac{1}{2}}(n+1)X,Y_1,Y_0+{\displaystyle \underset{k=1}{\overset{k=n}{}}}(1)^{k+1}a_k\overline{}_XY_k,Y_k`$
$`+{\displaystyle \underset{k=1}{\overset{k=n}{}}}(1)^{k+1}a_k(nk+1)XY_{k1}+{\displaystyle \frac{(n+k+1)}{4}}i_XY_{k+1},Y_k`$
$`=`$ $`L_XY_0,Y_0+{\displaystyle \underset{k=1}{\overset{k=n}{}}}(1)^{k+1}a_k\overline{}_XY_k,Y_k`$
$`{\displaystyle \frac{1}{2}}(n+1)X,Y_1,Y_0`$
$`+{\displaystyle \underset{k=1}{\overset{k=n}{}}}(1)^{k+1}a_k{\displaystyle \frac{(n+k+1)}{4}}i_XY_{k+1},Y_k`$
$`+{\displaystyle \underset{k=1}{\overset{k=n}{}}}(1)^{k+1}a_k(nk+1)XY_{k1},Y_k.`$
We make a change of variables in the very last sum, and get
$`_XY,Y`$
$`=`$ $`L_XY_0,Y_0+{\displaystyle \underset{k=1}{\overset{k=n}{}}}(1)^{k+1}a_k\overline{}_XY_k,Y_k`$
$`{\displaystyle \frac{1}{2}}(n+1)X,Y_1,Y_0`$
$`+{\displaystyle \underset{k=1}{\overset{k=n}{}}}(1)^{k+1}a_k{\displaystyle \frac{(n+k+1)}{4}}i_XY_{k+1},Y_k`$
$`+{\displaystyle \underset{k=0}{\overset{k=n1}{}}}(1)^ka_{k+1}(nk)XY_k,Y_{k+1}`$
Hence
$`_XY,Y`$
$`=`$ $`L_XY,Y`$
$`{\displaystyle \frac{1}{2}}(n+1)X,Y_1,Y_0+na_1XY_0,Y_1`$
$`{\displaystyle \underset{k=1}{\overset{k=n}{}}}(1)^k\left(a_{k+1}(nk)a_k{\displaystyle \frac{(n+k+1)}{4}}\right)XY_k,Y_{k+1}.`$
To conclude, we just have to remark that
$`a_1`$ $`=`$ $`{\displaystyle \frac{n+1}{2n}}`$
$`{\displaystyle \frac{a_{k+1}}{a_k}}`$ $`=`$ $`{\displaystyle \frac{n+k+1}{4(nk)}}.`$
We finally have to check that the corresponding representation of the group $`SL(2,)`$ is irreducible. For that let $`S^1SL(2,)`$, a subgroup isomorphic to the circle fixing a point $`x_0`$. The corresponding action on $`L_k(x_0)`$ is given by
$$e^{i\theta }(u)=e^{ki\theta }u.$$
This shows the representation is the $`2n+1`$ dimensional oneq.e.d.
## 3 Cohomology and $`(n+1)`$-holomorphic differentials
Let $`S=^2/\mathrm{\Gamma }`$ be a compact surface. Let $`\rho `$ be a $`(2n+1)`$-Fuchsian representation of $`\mathrm{\Gamma }`$. In this section, we shall describe the vector space $`H_\rho ^1(\mathrm{\Gamma },^{2n+1})`$ in terms of $`(n+1)`$-holomorphic differentials on $`S`$.
We use the notations of the previous sections. Let $`E_S=E/\mathrm{\Gamma }`$ be the vector bundle over $`S=^2/\mathrm{\Gamma }`$ coming from $`E`$.
Let $`^q`$ the vector space of $`q`$-holomorphic differentials on $`S`$. Let $`\mathrm{\Lambda }^p(E_\mathrm{\Gamma })`$ the vector space of $`p`$-forms on $`S`$ with value in $`E_S`$. The flat connection $``$ gives rise to a complex
$$0\mathrm{\Lambda }^0(E)\stackrel{d^{}}{}\mathrm{\Lambda }^1(E)\stackrel{d^{}}{}\mathrm{\Lambda }^2(E)0.$$
The cohomology of this complex is $`H_\rho ^{}(\mathrm{\Gamma },^{n+1})`$. From the metric on $`^2`$, we deduce an isomorphism $`\omega \stackrel{ˇ}{\omega }`$ of $`L_k^{}`$ with $`L_k`$. We define now a map $`\mathrm{\Phi }`$ by
$$\mathrm{\Phi }:\{\begin{array}{ccc}\hfill ^{2n+1}& & \mathrm{\Lambda }^1(E)\hfill \\ \hfill \omega & & (Xi_X\stackrel{ˇ}{\omega }L_nE)\hfill \end{array}$$
We first prove:
###### Proposition 3.1
For every $`(n+1)`$-holomorphic differential $`\omega `$
$$d^{}(\mathrm{\Phi }(\omega ))=0.$$
Furthermore, if $`\mathrm{\Phi }(\omega )=d^{}u`$, then $`\omega =0`$.
Proof: By definition,
$$d^{}\mathrm{\Phi }(\omega )(X,Y)=_Xi_Y\stackrel{ˇ}{\omega }_Yi_X\stackrel{ˇ}{\omega }i_{[X,Y]}\stackrel{ˇ}{\omega }.$$
Hence, if $`n>1`$
$$d^{}\mathrm{\Phi }(\omega )(X,Y)=\frac{2n}{4}(i_Xi_Y\stackrel{ˇ}{\omega }i_Yi_X\stackrel{ˇ}{\omega })+(\overline{}_Xi_Y\stackrel{ˇ}{\omega }\overline{}_Yi_X\stackrel{ˇ}{\omega }i_{[X,Y]}\stackrel{ˇ}{\omega }).$$
Notice that $`i_Xi_Y\stackrel{ˇ}{\omega }`$ is symmetric in $`X`$ and $`Y`$. Finally, the holomorphicity condition on $`\omega `$ implies
$$\overline{}_Xi_Y\stackrel{ˇ}{\omega }\overline{}_Yi_X\stackrel{ˇ}{\omega }i_{[X,Y]}\stackrel{ˇ}{\omega }=0.$$
A similar proof (but with different constants) yields the result for $`n=1`$.
Next, assume $`\mathrm{\Phi }(\omega )=d^{}u`$. The (non Riemaniann) metric on $`E`$ and the Riemannian metric on $`^2`$ induce a metric on $`\mathrm{\Lambda }^{}(E)`$, which we denote $`,_\mathrm{\Lambda }`$. One should notice here that even though this metric is neither positive nor negative, since $`\mathrm{\Phi }(\omega )`$ is a section of a bundle on which the metric is either positive or negative, we have
$$\mathrm{\Phi }(\omega ),\mathrm{\Phi }(\omega )_\mathrm{\Lambda }=0\mathrm{\Phi }(\omega )=0\omega =0$$
Let $`(d^{})^{}`$ be the adjoint of $`d^{}`$. One has, if $`(X_1,X_2)`$ is a basis of $`T^2`$,
$$(d^{})^{}(\varphi (\omega ))=\underset{k=1}{\overset{k=2}{}}_{X_k}(i_{X_k}\stackrel{ˇ}{\omega }).$$
A short calculation shows
$$(d^{})^{}(\varphi (\omega ))=\underset{k=1}{\overset{k=2}{}}\overline{}_{X_k}(i_{X_k}\stackrel{ˇ}{\omega }),$$
and this last term is $`0`$ by holomorphicity. We have just proved that
$$(d^{})^{}\mathrm{\Phi }(\omega )=0.$$
Hence, $`\mathrm{\Phi }(\omega )=d^{}u`$ implies
$$\mathrm{\Phi }(\omega ),\mathrm{\Phi }(\omega )_\mathrm{\Lambda }=(d^{})^{}\mathrm{\Phi }(\omega ),u_\mathrm{\Lambda }=0.$$
This ends the proof q.e.d.
It follows from the previous proposition that $`\mathrm{\Phi }`$ gives rise to a map (also denoted $`\mathrm{\Phi }`$) from $`^{n+1}`$ to the space $`H_\rho ^1(\mathrm{\Gamma },^{2n+1})`$. We have:
###### Corollary 3.2
The map $`\mathrm{\Phi }`$ is an isomorphism from $`^{n+1}`$ to the space $`H_\rho ^1(\mathrm{\Gamma },^{2n+1})`$.
Proof: Indeed, we have just proved that $`\mathrm{\Phi }`$ is injective. Furthermore, if $`\chi (S)`$ is the Euler characteristic of $`S`$, we have
$$dim(H_\rho ^1(\mathrm{\Gamma },^{2n+1}))(2n+1)\chi (S).$$
But, by Riemann-Roch,
$$dim(^{n+1})=(2n+1)\chi (S).$$
Hence, the corollary follows q.e.d.
## 4 A de Rham interpretation of Margulis invariant
The irreducible representation of $`SL(2,)`$ of dimension $`2n+1`$ preserves a metric $`,`$ of signature $`(n,n+1)`$.
### 4.1 Loxodromic elements
We define a loxodromic element in $`SO(n,n+1)`$ to be $``$-split and in the interior of a Weyl chamber. This just means all eigenvalues are real and have multiplicity 1. Recall that 1 allways belong to the spectrum of a loxodromic element. Notice that all the elements, except the identity, of a $`(2n+1)`$-Fuchsian surface group are loxodromic.
### 4.2 The invariant vector of a loxodromic element
Chose now, once and for all, an orientation on $`^{2n+1}`$. The light cone - without the origin - has two components. Let’s also choose one of these components.
Let $`\gamma `$ be a loxodromic element. It follows from the previous choices that we have a well defined eigenvector, the invariant vector, denoted $`v_\gamma `$, associated to the eigenvalue 1.
Indeed, all the other eigenvectors are lightlike. We order all the eigenvaluesdifferent than 1, in such a way that $`\lambda _i<\lambda _{i+1}`$. Thanks to our choices, we may pick one eigenvector $`e_i`$ in the preferred component of the light cone for all the eigenvalues $`\lambda _i`$ different than 1. We now choose $`v_\gamma `$ of norm 1, such that $`(v_\gamma ,e_1,\mathrm{},e_{2n})`$ is positively oriented.
### 4.3 Margulis invariant
Let now $`Iso(n,n+1)=^{2n+1}SO(n,n+1)`$ be the group of orientation preserving isometries of $`^{2n+1}`$ as an affine space. For $`\gamma `$ in $`Iso(n,n+1)`$, $`\widehat{\gamma }`$ denotes its linear part. We shall say an element of $`Iso(n,n+1)`$ is loxodromic if its linear part is a loxodromic element of $`SO(n,n+1)`$.
The Margulis invariant ( ) of a loxodromic element $`\gamma `$ of $`Iso(n,n+1)`$ is
$$\mu (\gamma )=\gamma (x)x,v_{\widehat{\gamma }},$$
where $`x`$ is an element of $`^{2n+1}`$. A quick check shows $`\mu (\gamma )`$ does not depend on $`x`$.
### 4.4 Margulis invariant and properness of an affine action
Let $`\gamma _1`$ and $`\gamma _2`$ be two loxodromic elements. Let $`E_i^+`$ (resp. $`E_i^{}`$) be the space generated by the eigenvectors of $`\widehat{\gamma }_i`$ corresponding to the eigenvalues of absolute value greater (resp. less) than 1.
We say $`\gamma _1`$ and $`\gamma _2`$ are in general position if the two decompositions
$$.v_{\widehat{\gamma }_i}E_i^+E_i^{},$$
are in general position.
Notice that for a $`(2n+1)`$-Fuchsian group, two (non comensurable) elements are loxodromic and in general position.
In (see also ) G. Margulis has proved the following magic lemma
###### Lemma 4.1
If two loxodromic elements $`\gamma _1`$, $`\gamma _2`$, in general position, are such that $`\mu (\gamma _1)\mu (\gamma _2)0`$, then the group generated by $`\gamma _1`$ and $`\gamma _2`$ does not act properly on $`^{2n+1}`$
### 4.5 An interpretation of Margulis invariant
Let $`\rho `$ a representation of $`\mathrm{\Gamma }`$ in $`Iso(n,n+1)`$, whose linear part, $`\widehat{\rho }`$, is Fuchsian. Let $`E_S=L_1\mathrm{}L_n`$ the flat bundle over $`S`$ described in 2.1 whose holonomy is $`\widehat{\rho }`$.
We describe now $`\rho `$ as an element of $`H_{\widehat{\rho }}^1(\mathrm{\Gamma },^{2n+1})`$.
Let $`\alpha H_\rho ^1(\mathrm{\Gamma },^{2n+1})`$, interpreted as an element of $`\mathrm{\Lambda }^1(E_S)`$. Let $`^\alpha `$ be the flat connection on $`F=E_S`$ defined by
$$_X^\alpha (\lambda ,V)=(L_X\lambda ,\lambda .\alpha (X)+_XV).$$
We claim there exists $`\alpha H_{\widehat{\rho }}^1(\mathrm{\Gamma },^{2n+1})`$ such that the holonomy of $`^\alpha `$ is $`\rho `$. Of course, here, $`^pSL(p,)`$ is identified with a subgroup of $`GL(p+1,)`$.
Let now $`c`$ be a closed curve on $`S`$, represented in homotopy by the conjugacy class of some element $`\gamma `$. Since $`v_{\widehat{\rho }(\gamma )}`$ is invariant under $`\widehat{\rho }(\gamma )`$, it gives rise to a parallel section $`v_c`$ of $`E|_c`$.
We first prove the following statement:
###### Proposition 4.2
Let $`c`$, $`\rho `$, $`\gamma `$, $`\alpha `$ be as above. Then
$$\mu (\rho (\gamma ))=_c\alpha ,v_c.$$
Proof: We shall use the previous notations. We parametrise $`c`$ by the circle of length 1. Let $`\pi `$ be the covering $`^2S`$. Consider a lift $`\stackrel{~}{c}`$ of $`c`$ on the universal cover of $`S`$. The bundle $`\pi ^{}F`$ becomes trivial. The canonical section $`\sigma `$ corresponding to the $``$ factor in $`F`$, gives rise to a map
$$i:^2^{2n+1},$$
taking value in the affine hyperplane
$$P=\{(1,u)^{2n+1}\}.$$
Let $`\overline{c}=i\stackrel{~}{c}`$, and let’s identify $`\rho (\gamma )`$ with $`\gamma `$. By definition now:
$`\mu (\gamma )`$ $`=`$ $`\rho (\gamma )(\overline{c}(0))\overline{c}(0),v_{\widehat{\gamma }}`$
$`=`$ $`\overline{c}(1)\overline{c}(0),v_{\widehat{\gamma }}`$
$`=`$ $`{\displaystyle _0^1}\dot{\overline{c}}(s),v_{\widehat{\gamma }}𝑑s`$
Now, we interpret the last term on $`F`$ and we obtain
$`\mu (\gamma )`$ $`=`$ $`{\displaystyle _0^1}_{\dot{c}(s)}^\alpha \sigma ,v_c(s)𝑑s`$
$`=`$ $`{\displaystyle _0^1}\alpha (\dot{c}(s)),v_c(s)𝑑s`$
$`=`$ $`{\displaystyle _c}\alpha ,v_c.`$
This ends the proof q.e.d.
### 4.6 The invariant vector as a section
In this paragraph, we assume $`n=2p+1`$, so that our representation is of dimension $`4p+3`$.
We use the notations of the previous paragraphs. In particular, let $`\gamma \mathrm{\Gamma }`$. Let $`v=v_{\widehat{\rho }(\gamma )}`$. Let $`c`$ be the closed geodesic (for the hyperbolic metric) corresponding to the element $`\gamma `$.
Recall that $`v_\gamma `$ gives rise to a section $`v_c`$ along the closed geodesic, which is parallel.
In this paragraph, we wish to describe $`v_c`$ explicitely. Let $`J`$ the complex structure of $`S`$. Let’s introduce the following section (along c) defined by
$`(w_c)_{2k}`$ $`=`$ $`0`$
$`(w_c)_{2k+1}`$ $`=`$ $`J(4)^k{\displaystyle \underset{l=1}{\overset{l=k}{}}}({\displaystyle \frac{pl}{p+l+1}})\underset{2k+1}{\underset{}{\dot{c}\mathrm{}\dot{c}}}.`$
###### Proposition 4.3
The section $`w_c`$ of $`E_S`$ is parallel along $`c`$. Furthermore, there exists $`\epsilon \{1,1\}`$ independant of $`c`$ such that
$$v_c=\epsilon \frac{w_c}{\sqrt{w_c,w_c}}.$$
Proof: A straightforawrd computation shows that $`w_c`$ (hence $`v_c`$) is parallel. Furthermore $`w_c`$ is a space like vector, and by construction $`v_c`$ has norm 1.
It remains to prove that $`v_c`$ has the correct orientation. For that consider any geodesic arc $`u`$ in $`^2`$ paramatrised by $`[0,L]`$. We have a basis of $`E|_{u(t)}`$ given by
$$B(t)=(1,\dot{u},J\dot{u},\mathrm{},\underset{𝑛}{\underset{}{\dot{u}\mathrm{}\dot{u}}},J\underset{𝑛}{\underset{}{\dot{u}\mathrm{}\dot{u}}}).$$
We may now consider the isometry $`\gamma (u)`$ sending $`B(0)`$ to $`B(L)`$. This is a loxodromic isometry. Next, consider the following section of $`E`$ along $`u`$ given by
$`(w_u)_{2k}`$ $`=`$ $`0`$
$`(w_u)_{2k+1}`$ $`=`$ $`J(4)^k{\displaystyle \underset{l=1}{\overset{l=k}{}}}({\displaystyle \frac{pl}{p+l+1}})\underset{2k+1}{\underset{}{\dot{u}\mathrm{}\dot{u}}}.`$
This section is parallel along $`u`$ and therefore gives rise to a vector proportional to the invariant vector of $`\gamma (u)`$.
Next, by continuity, this proportional is constant. Applying this remark to a lift in the universal cover of our closed geodesic, this ends the proof q.e.d.
## 5 Main theorem in dimension $`4p+3`$
Again, let’s $`\rho `$ be a representation of a compact surface group $`\mathrm{\Gamma }`$ in the group of affine transformation of an affine space of dimension $`4p+3`$, whose linear part $`\widehat{\rho }`$ is Fuchsian. We assume that $`\rho (\mathrm{\Gamma })`$ acts properly on $`^{4p+3}`$. The representation $`\rho `$ is described from $`\widehat{\rho }`$ as an element $`\alpha `$ of $`H_{\widehat{\rho }}^1(\mathrm{\Gamma },^{4p+3})`$.
According to proposition 3.2, this element $`\alpha `$ is described by a $`(2p+2)`$-holomorphic differential $`\omega `$.
Let $`\gamma \mathrm{\Gamma }`$, and $`c`$ the corresponding closed geodesic. From proposition 4.2, we get
$$\mu (\rho (\gamma ))=_c\alpha ,v_c.$$
From proposition 4.3, we deduce there exist a constant $`K_1`$ just depending on $`p`$ such that
$$\mu (\rho (\gamma ))=K_1_ci_{\dot{c}}\omega ,J\underset{2p+1}{\underset{}{\dot{c}\mathrm{}\dot{c}}}𝑑t.$$
From the constructions explained in the paragraph 2.1, we finally obtain there exist a contant $`K_2`$ just depending on $`p`$ such that
$`\mu (\rho (\gamma ))`$ $`=`$ $`K_2{\displaystyle _c}i_{\dot{c}}\stackrel{ˇ}{\omega },J\underset{2p+1}{\underset{}{\dot{c}\mathrm{}\dot{c}}}𝑑t`$
$`=`$ $`K_2{\displaystyle _c}\mathrm{}(\omega (\underset{2p+2}{\underset{}{\dot{c}\mathrm{}\dot{c}}}))𝑑t.`$
Let $`US`$ be the unit tangent bundle of $`S`$. Let $`f`$ be the function defined on $`US`$ by
$$f(u)=\mathrm{}(\omega (\underset{2p+2}{\underset{}{u\mathrm{}u}})).$$
From lemma 4.1 and the previous computation, we obtain that the integral of $`f`$ along closed orbits of the geodesic flow has a constant sign. On the other hand, let $`\lambda `$ be the Lebesgue measure, we have
$$_{US}f𝑑\lambda =0.$$
Indeed, let $`\beta `$ be a complex number such that $`\beta ^{2p+2}=1`$. This number $`\beta `$ is to be considered as a diffeomorphism of $`US`$, which preserves the orientation and the Lebesgue measure. Lastly $`f\beta =f`$ and this proves the last formula.
The conclusion of the proof follows at once from the following lemma.
###### Lemma 5.1
Let $`M`$ be a compact manifold equipped with an Anosov flow preserving a measure $`\nu `$ which charges open sets. Let $`f`$ be a Hölder function defined on $`M`$ such that its integral on every closed orbit is positive, then the integral of $`f`$ with respect to $`\nu `$ is positive.
Proof: I could not find a proper reference in the litterature of this specific lemma, although lots of versions exist for discrete time transformations. Let sketch a proof by overkilling using hints from a conversation with G. Margulis. Let $`\varphi _t`$ be the flow. It follows from a theorem of M. Ratner that if $`f`$ is a Hölder function whose integral is 0, either it is a cocycle (and in particular, its integral over every closed orbit is 0), or it satisfies the central limit theorem. In particular, this implies that there exists at leaat one point $`x`$ with a dense orbit, such that
$$\underset{t\mathrm{}}{lim}\frac{1}{\sqrt{t}}_0^tf\varphi _s(x)𝑑s<A<0.$$
Let now $`\{t_n\}_n`$ be a sequence of real numbers converging to infinity such that $`\{\varphi _{t_n}(x)\}_n`$ converges to $`x`$. By the closing lemma and classical estimates, there exists $`n_0`$ such that for $`n>n_0`$ we can find a closed geodesic $`c_n`$ such that
$$|_0^{t_n}f\varphi _s(x)𝑑s_{c_n}f𝑑s|<A/2.$$
This implies the lemma.q.e.d.
## 6 Other dimensions
The other dimensions are either trivial (even case) or follows from the immediate use of Margulis invariant ($`4p+1`$ case) as we shall explain now.
Let $`\lambda `$ be the representation of $`SL(2,)`$ of even dimension. Let $`h`$ be a loxodromic element of $`SL(2,)`$. Then $`1`$ will not belong to the spectrum of $`\lambda (h)`$. It follows, that if $`\rho `$ is a representation of $`\mathrm{\Gamma }`$ in even dimension whose linear part is Fuchsian then for all $`\gamma `$ in $`\mathrm{\Gamma }`$ different than the identity then $`\rho (\gamma )`$ does not act properly.
Last, in dimensions $`4p+1`$, the Margulis invariant is such that $`\mu (\gamma ^1)=\mu (\gamma )`$. It follows at once from lemma 4.1, that if $`\rho `$ is a representation of $`\mathrm{\Gamma }`$ in dimension $`4p+1`$ whose linear part is Fuchsian, if $`\gamma _1`$ and $`\gamma _2`$ are non commensurable elements of $`\mathrm{\Gamma }`$ then $`\rho (\gamma _1)`$ and $`\rho (\gamma _2)`$ generate a group that does not act properly on the affine space. Of course, the point in our previous discussion in that in dimension $`4p+3`$ then $`\mu (\gamma ^1)=\mu (\gamma )`$, hence such an argument do not work and actually, free groups (even Fuchsian ones) can act properly, see , and .
François Labourie
Topologie et Dynamique
Université Paris-Sud
F-91405 Orsay (Cedex)
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# Charmonium absorption by nucleons
## I Introduction
Very recently the NA50 Collaboration reported measurements of $`J/\mathrm{\Psi }`$ suppression in $`Pb+Pb`$ collisions at the CERN-SPS. It was claimed in Ref. that these new experimental data ruled out conventional hadronic models of $`J/\mathrm{\Psi }`$ suppression and thus indicated the formation of a deconfined state of quarks and gluons, namely the quark gluon plasma (QGP). Presently, there are calculations using two different methods, namely the cascade and the Glauber type model , that reproduce the new NA50 data reasonably well up to the highest transverse energies based on hadronic $`J/\mathrm{\Psi }`$ dissociation alone. In view of this, the evidence for QGP formation based on anomalous $`J/\mathrm{\Psi }`$ suppression is not obvious and an interpretation of the NA50 data can be further considered in terms of the hadronic dissociation of the $`J/\mathrm{\Psi }`$ meson by the nucleons and comovers.
However, the long term puzzle of the problem is the strength of the hadronic dissociation of the $`J/\mathrm{\Psi }`$ meson. The various calculations of $`J/\mathrm{\Psi }`$ production from heavy ion collisions in practice adopt different dissociation cross sections. For instance, the cross section for $`J/\mathrm{\Psi }+N`$ dissociation used in these models ranges from 3.0 up to 6.7 mb when dealing with the same experimental data on $`J/\mathrm{\Psi }`$ production from $`A+A`$ collisions. Moreover, it is presently accepted that the dissociation cross section does not depend on the $`J/\mathrm{\Psi }`$ energy and in the available calculations it enters as a constant. Furthermore, to a certain extent, both of the cross sections for $`J/\mathrm{\Psi }`$ dissociation on nucleons and comovers were taken as free parameters, finally adjusted to the heavy ion data.
Moreover, the $`J/\mathrm{\Psi }+N`$ cross section was evaluated from experimental data on $`J/\mathrm{\Psi }`$ meson production in $`\gamma +A`$ and $`p+A`$ reactions. The analysis of the $`J/\mathrm{\Psi }`$ photoproduction from nuclei at a mean photon energy of 17 GeV indicates a $`J/\mathrm{\Psi }+N`$ cross section of 3.5$`\pm `$0.9 mb. The combined analysis of experimental data on $`J/\mathrm{\Psi }`$ production from $`p+A`$ collisions at beam energies from 200 to 800 GeV provides a $`J/\mathrm{\Psi }+N`$ cross section of 7.3$`\pm `$0.6 mb. On the other hand, the $`J/\mathrm{\Psi }+N`$ cross section evaluated by the vector dominance model from the $`\gamma +NJ/\mathrm{\Psi }+N`$ data is about 1 mb for $`J/\mathrm{\Psi }`$ energies in the nucleon rest frame from 8 to 250 GeV.
In principle, this ambiguity in the $`J/\mathrm{\Psi }`$ dissociation cross section does not allow one to claim a consistent interpretation of the NA50 results. Although the experimental data might be well reproduced by the more recent calculations considering the hadronic dissociation of the $`J/\mathrm{\Psi }`$ meson, the $`J/\mathrm{\Psi }+N`$ cross section was introduced therein as a free parameter. In Ref. we provided results for $`J/\mathrm{\Psi }`$ dissociation on comovers and motivated the large rate of this process as arising from the in-medium modification of the dissociation amplitude, which is obviously different from that given in free space . Most recent results from the different models on $`J/\mathrm{\Psi }`$ dissosiation on light hadrons are given in Refs. In the present study we apply the hadronic model to $`J/\mathrm{\Psi }+N`$ dissociation.
## II Lagrangian densities, coupling constants and form factors
Within the boson exchange model we consider the reactions depicted in Fig. 1 and use the interaction Lagrangian densities:
$`_{JDD}`$ $`=`$ $`ig_{JDD}J^\mu \left[\overline{D}(_\mu D)(_\mu \overline{D})D\right],`$ (1)
$`_{DN\mathrm{\Lambda }_c}`$ $`=`$ $`ig_{DN\mathrm{\Lambda }_c}\left(\overline{N}\gamma _5\mathrm{\Lambda }_cD+\overline{D}\overline{\mathrm{\Lambda }_c}\gamma _5N\right),`$ (2)
$`_{JD^{}D}`$ $`=`$ $`{\displaystyle \frac{g_{JD^{}D}}{m_J}}\epsilon _{\alpha \beta \mu \nu }(^\alpha J^\beta )`$ (4)
$`\times [(^\mu \overline{D}^\nu )D+\overline{D}(^\mu D^\nu )],`$
$`_{D^{}N\mathrm{\Lambda }_c}`$ $`=`$ $`g_{D^{}N\mathrm{\Lambda }_c}\left(\overline{N}\gamma _\mu \mathrm{\Lambda }_cD^\mu +\overline{D}^\mu \overline{\mathrm{\Lambda }_c}\gamma _\mu N\right),`$ (5)
where, $`N=(\begin{array}{c}p\\ n\end{array})`$, $`\overline{N}=N^{}\gamma _0`$, $`D(\begin{array}{c}D^0\\ D^+\end{array})`$ (creation of the meson states), $`\overline{D}=D^{}`$, and similar notations for the $`D^{}`$ and $`\overline{D}^{}`$ should be understood.
The $`JDD`$ coupling constant was derived in Refs. and in the following calculations we adopt $`g_{JDD}`$=7.64. Following the SU(4) relations from Ref. the $`DN\mathrm{\Lambda }_c`$ coupling constant was taken equal the $`KN\mathrm{\Lambda }`$ coupling constant. The analysis of experimental data provides 13.2$`g_{KN\mathrm{\Lambda }}`$15.7 and in the following we use $`g_{DN\mathrm{\Lambda }_c}`$=14.8. However, we notice, that the calculations by the QCD sum rules suggest a smaller coupling $`g_{DN\mathrm{\Lambda }_c}6.7\pm 2.1`$. The $`D^{}N\mathrm{\Lambda }_c`$ coupling constant might be again related to the $`K^{}N\mathrm{\Lambda }`$ constant. The $`K^{}N\mathrm{\Lambda }`$ constant is given in Ref. within the range –18.8$`g_{K^{}N\mathrm{\Lambda }}`$–21.4, while it varies between $``$–5.9 and –22 in Ref. . In the following calculations we adopt $`g_{D^{}N\mathrm{\Lambda }_c}`$=–19 . Furthermore, we assume that $`g_{JDD}=g_{JDD^{}}`$ .
The form factors associated with the interaction vertices were parameterized in a conventional monopole form
$$F(t)=\frac{\mathrm{\Lambda }^2}{\mathrm{\Lambda }^2t},$$
(6)
where $`t`$ is the four-momentum transfer and $`\mathrm{\Lambda }`$ is the cutoff parameter. The introduction of the form factors is dictated by the extended structure of the hadrons . The form factor may be estimated through the vector dominance model (VDM) . The VDM predicts an increase of the form factor with increasing mass of the vector meson and a purely phenomenological prediction provides $`\mathrm{\Lambda }`$=$`m_J`$. In the following we will explicitly use the cutoff parameter $`\mathrm{\Lambda }`$=3.1 GeV for the form factors at the $`JDD`$ and $`JDD^{}`$ vertices.
Furthermore, one should also introduce form factors at the $`DN\mathrm{\Lambda }_c`$ and $`D^{}N\mathrm{\Lambda }_c`$ vertices. Again, we use the monopole form, but there are no direct ways to evaluate the relevant cutoff parameter. In Refs. the cutoff parameters were adjusted to fit the empirical nucleon-nucleon data and range from 1.3 to 2 GeV depending on the exchange meson coupled to the $`NN`$ system. Based on these results, in the following we use a monopole form factor with $`\mathrm{\Lambda }`$=2 GeV for the $`DN\mathrm{\Lambda }_c`$ and $`D^{}N\mathrm{\Lambda }_c`$ vertices.
In principle, the coupling constants and the cutoff parameters in the form factors discussed in this section can be adjusted to the experimental data on the total $`J/\mathrm{\Psi }+N`$ cross section evaluated from different nuclear reactions, as $`\gamma +AJ/\mathrm{\Psi }+X`$ and $`p+AJ/\mathrm{\Psi }+X`$. On the other hand, these parameters can be also fixed by comparison to the short range QCD or Regge theory calculations at high energies. As will be shown later, the set of coupling constants and cutoff parameters proposed above are rather well fitted to the relevant experimental data and the results from theoretical calculations with other models.
## III The $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ and $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}^{}`$ reactions
The diagrams for $`J/\mathrm{\Psi }`$ dissociation by a nucleon with the production of the $`\mathrm{\Lambda }_c+\overline{D}`$ and $`\mathrm{\Lambda }_c+\overline{D}^{}`$ final states are shown in Fig. 1a)-c). They involve the $`J/\mathrm{\Psi }D+\overline{D}`$ and $`J/\mathrm{\Psi }D+\overline{D}^{}`$ vertices and are OZI allowed. These reactions are endothermic, since the total mass of the final states is larger than the total mass of the initial $`J/\mathrm{\Psi }`$-meson and nucleon.
The $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ reaction via $`D`$ meson exchange is shown in the Fig. 1a) and the corresponding amplitude with an amplitude added in order to preserve gauge invariance in the limit, $`m_J0`$, may be given by:
$`_a`$ $`=`$ $`2ig_{JDD}g_{DN\mathrm{\Lambda }_c}(ϵ_Jp_{\overline{D}})`$ (8)
$`\times \left({\displaystyle \frac{1}{q^2m_D^2}}+{\displaystyle \frac{1}{2p_Jp_{\overline{D}}}}\right)\overline{u}_{\mathrm{\Lambda }_c}(p_{\mathrm{\Lambda }_c})\gamma _5u_N(p_N),`$
where $`q=p_Jp_{\overline{D}}(=p_{\mathrm{\Lambda }_c}p_N)`$, while the amplitude via the $`D^{}`$ meson exchange is shown in Fig. 1b) and can be written as
$`_b={\displaystyle \frac{g_{JD^{}D}g_{D^{}N\mathrm{\Lambda }_c}}{m_J}}{\displaystyle \frac{1}{q^2m_D^{}^2}}`$ (9)
$`\times \epsilon _{\alpha \beta \mu \nu }p_J^\alpha ϵ_J^\beta q^\mu \overline{u}_{\mathrm{\Lambda }_c}(p_{\mathrm{\Lambda }_c})\gamma ^\nu u_N(p_N).`$ (10)
Furthermore, the amplitude for the $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}^{}`$ reaction due to $`D`$ meson exchange is shown in Fig. 1c), and is given by
$`_c={\displaystyle \frac{ig_{JD^{}D}g_{DN\mathrm{\Lambda }_c}}{m_J}}{\displaystyle \frac{1}{q^2m_D^2}}`$ (11)
$`\times \epsilon _{\alpha \beta \mu \nu }p_J^\alpha ϵ_J^\beta p_{\overline{D}^{}}^\mu ϵ_{\overline{D}^{}}^\nu \overline{u}_{\mathrm{\Lambda }_c}(p_{\mathrm{\Lambda }_c})\gamma _5u_N(p_N).`$ (12)
In Eqs. (8)- (12) $`ϵ_J`$ and $`ϵ_{\overline{D}^{}}`$ are respectively the polarization vector of the $`J/\mathrm{\Psi }`$ meson and $`\overline{D}^{}`$ meson, and $`q=p_{\mathrm{\Lambda }_c}p_N`$ denotes the four-momentum transfer. We notice that there arises no interference term among the amplitudes, $`_a`$, $`_b`$ and $`_c`$, because of the Dirac structure, when spin components are not specified and summed over all spins for the $`N`$ and $`\mathrm{\Lambda }_c`$, and the different final states.
In our normalization the corresponding differential cross sections in the center of mass frame of the $`J/\mathrm{\Psi }`$ meson and nucleon system are given by:
$`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}_{a,b}={\displaystyle \frac{1}{64\pi ^2s}}|\overline{_{a,b}}|^2`$ (13)
$`\times `$ $`\left({\displaystyle \frac{[(m_{\mathrm{\Lambda }_c}+m_D)^2s][(m_{\mathrm{\Lambda }_c}m_D)^2s]}{[(m_N+m_J)^2s][(m_Nm_J)^2s]}}\right)^{1/2},`$ (15)
$`{\displaystyle \frac{d\sigma }{d\mathrm{\Omega }}}_c={\displaystyle \frac{1}{64\pi ^2s}}|\overline{_c}|^2`$
$`\times `$ $`\left({\displaystyle \frac{[(m_{\mathrm{\Lambda }_c}+m_D^{})^2s][(m_{\mathrm{\Lambda }_c}m_D^{})^2s]}{[(m_N+m_J)^2s][(m_Nm_J)^2s]}}\right)^{1/2},`$ (16)
where $`s=(p_N+p_J)^2`$ is the squared invariant collision energy and $`|\overline{_{a,b,c}}|^2`$ are the corresponding amplitudes squared, averaged over the initial and summed over the final spins. Explicitly, they are given by:
$`|\overline{_a}|^2={\displaystyle \frac{8g_{JDD}^2g_{DN\mathrm{\Lambda }_c}^2}{3m_J^2}}\left({\displaystyle \frac{1}{q^2m_D^2}}+{\displaystyle \frac{1}{2p_Jp_{\overline{D}}}}\right)^2`$ (17)
$`\times `$ $`(p_Np_{\mathrm{\Lambda }_c}m_Nm_{\mathrm{\Lambda }_c})[(p_Jp_{\overline{D}})^2m_J^2m_D^2],`$ (19)
$`|\overline{_b}|^2={\displaystyle \frac{g_{JD^{}D}^2g_{D^{}N\mathrm{\Lambda }_c}^2}{3m_J^2}}{\displaystyle \frac{1}{(q^2m_D^{}^2)^2}}`$
$`\times `$ $`[m_J^2(p^2q^2(m_{\mathrm{\Lambda }_c}^2m_N^2)^2)`$ (24)
$`+2(p_Jp)(p_Jq)(m_{\mathrm{\Lambda }_c}^2m_N^2)`$
$`p^2(p_Jq)^2q^2(p_Jp)^2`$
$`4(p_Np_{\mathrm{\Lambda }_c}m_Nm_{\mathrm{\Lambda }_c})(m_J^2q^2(p_Jq)^2)],`$
$`|\overline{_c}|^2={\displaystyle \frac{4g_{JD^{}D}^2g_{DN\mathrm{\Lambda }_c}^2}{3m_J^2}}{\displaystyle \frac{1}{(q^2m_D^2)^2}}`$
$`\times `$ $`(p_Np_{\mathrm{\Lambda }_c}m_Nm_{\mathrm{\Lambda }_c})[(p_Jp_{\overline{D}^{}})^2m_J^2m_D^{}^2],`$ (25)
with $`pp_{\mathrm{\Lambda }_c}+p_N`$.
Finally, Fig. 2a) shows the $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ cross section as a function of the invariant collision energy $`\sqrt{s}`$ calculated with $`D`$ (solid line) and $`D^{}`$ meson (dashed line) exchanges and with form factors at the interaction vertices. The $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}^{}`$ cross section is shown in Fig. 2b). Both reactions substantially contribute at low energies, close to the respective thresholds, while their contribution to the $`J/\mathrm{\Psi }+N`$ dissociation cross section decreases with increasing $`\sqrt{s}`$.
## IV The $`J/\mathrm{\Psi }+NN+D+\overline{D}`$ reaction
The diagram for $`J/\mathrm{\Psi }`$ dissociation by a nucleon with the production of the $`D+\overline{D}`$ state is shown in Fig. 1d). The dissociation involves the $`J/\mathrm{\Psi }D+\overline{D}`$ vertex and therefore it is an OZI allowed process. The total mass of the produced particles is larger than the total mass of the initial $`J/\mathrm{\Psi }`$-meson and nucleon and thus these reactions are endothermic. Furthermore, the processes with $`D+N`$ and $`\overline{D}+N`$ scattering are different and both contribute to $`J/\mathrm{\Psi }`$ dissociation on the nucleon.
The amplitude for the reaction $`J/\mathrm{\Psi }+NN+D+\overline{D}`$ can be parametrized as
$$_d=2g_{JDD}\frac{(ϵ_Jp_D)}{tm_D^2}F(t)_{𝒟𝒩},$$
(26)
where $`p_D`$ and $`m_D`$ are the $`D`$-meson four-momentum and mass, respectively, $`ϵ_J`$ denotes the polarization vector of $`J/\mathrm{\Psi }`$-meson, $`t`$ is squared four-momentum transfered from initial $`J/\mathrm{\Psi }`$ to final $`D`$-meson, while $`g_{JDD}`$ and $`F(t)`$ are the coupling constant and the form factor at the $`J/\mathrm{\Psi }D\overline{D}`$ vertex, respectively. Furthermore, $`_{𝒟𝒩}`$ is the amplitude for $`D+N`$ or $`\overline{D}+N`$ scattering, which is related to the physical cross section as
$$|_{𝒟𝒩}|^2=16\pi s_1\sigma _{DN}(s_1),$$
(27)
where $`s_1`$ is the squared invariant mass of the $`D+N`$ or $`\overline{D}+N`$ subsystem.
The double differential cross section for the reaction $`J/\mathrm{\Psi }+NN+D+\overline{D}`$ can be written as
$`{\displaystyle \frac{d^2\sigma _{JN}}{dtds_1}}={\displaystyle \frac{g_{JDD}^2}{96\pi ^2q_J^2s}}q_D\sqrt{s_1}{\displaystyle \frac{F^2(t)}{(tm_D)^2}}`$ (28)
$`\times {\displaystyle \frac{[(m_J+m_D)^2t][(m_Jm_D)^2t]}{m_J^2}}\sigma _{DN}(s_1),`$ (29)
where $`q_D`$ is given by
$$q_D^2=\frac{[(m_N+m_D)^2s_1][(m_Nm_D)^2s_1]}{4s_1},$$
(30)
and $`m_J`$ and $`m_N`$ denote the $`J/\mathrm{\Psi }`$-meson and nucleon masses, respectively.
As was proposed in Ref. , by replacing the elastic scattering cross section, $`\sigma _{DN}`$, in Eq. (29) with the total cross section, it is possible to account simultaneously for all available final states that can be produced at the relevant vertex. Since we are interested in inclusive $`J/\mathrm{\Psi }`$ dissociation on the nucleon, in the following we use the total $`D+N`$ and $`\overline{D}+N`$ interaction cross sections. However, we denote this inclusive dissociation generically as the $`J/\mathrm{\Psi }+NN+D+\overline{D}`$ reaction.
The total $`D+N`$ and $`\overline{D}+N`$ cross sections were evaluated in our previous study by considering the quark diagrammatic approach. It was proposed that the $`D+N`$ and $`\overline{D}+N`$ cross sections should be equal to the $`\overline{K}+N`$ and $`K+N`$ cross sections, respectively, at the same invariant collision energy and neglecting the contribution from the baryonic resonances coupled to the $`\overline{K}+N`$ system. The total $`\overline{D}+N`$ cross section was taken as a constant, $`\sigma _{\overline{D}N}`$=20 mb, while the total $`D+N`$ cross section can be parameterized as
$`\sigma _{\overline{D}N}(s_1)`$ $`=`$ $`\left({\displaystyle \frac{[(m_{\mathrm{\Lambda }_c}+m_\pi )^2s_1][(m_{\mathrm{\Lambda }_c}m_\pi )s_1]}{[(m_D+m_N)^2s_1][(m_Dm_N)s_1]}}\right)^{1/2}`$ (31)
$`\times `$ $`{\displaystyle \frac{27}{s_1}}+20,`$ (32)
where $`m_L`$ and $`m_\pi `$ are the $`\mathrm{\Lambda }_c`$-hyperon and pion masses, respectively, given in GeV and the cross section is given in mb.
Some of the partial $`D+N`$ cross sections were calculated Ref. by an effective Lagrangian approach, where the identity between the $`\overline{K}+N`$ and $`D+N`$ interaction was imposed by the $`SU(4)`$ relations. Indeed, on the basis of $`SU(4)`$ symmetry the couplings for the vertices containing $`\overline{K}`$ or $`K`$ mesons were taken to be equal to those obtained by replacing $`\overline{K}`$ or $`K`$ with $`D`$ or $`\overline{D}`$, respectively. For instance, $`g_{\pi DD^{}}`$=$`g_{\pi KK^{}}`$, $`g_{DN\mathrm{\Lambda }_c}`$=$`g_{KN\mathrm{\Lambda }}`$, etc. In the more general case, the large variety of possible diagrams describing the $`\overline{K}+N`$ and $`K+N`$ interactions can be identified to those for the $`D+N`$ and $`\overline{D}+N`$ interactions, respectively. Thus, the estimates given in Ref. are supported by the boson exchange formalism. The comparison between the $`DN`$ scattering cross section calculated by an effective Lagrangian and evaluated by the quark diagrammatic approaches is given in Ref. . Furthermore, for the $`J/\mathrm{\Psi }+NN+D+\overline{D}`$ calculations we used coupling constants and form factors similar to those used for the $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ reaction.
Figure 3 shows the cross section for $`J/\mathrm{\Psi }`$ dissociation on the nucleon with the production of $`\overline{D}+D`$ pairs calculated using Eq. (29), with and without form a factor at the $`J/\mathrm{\Psi }\overline{D}D`$ vertex. The lines in Fig. 3 show the results where the processes with $`\overline{D}+N`$ and $`D+N`$ interactions were summed incoherently.
## V Total $`J/\mathrm{\Psi }+N`$ cross section and comparison with QCD calculations
Our results for the total $`J/\mathrm{\Psi }+N`$ cross section are shown in Fig. 4 as the solid line a). Here the total cross section is given as a sum of the partial $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ cross section calculated with $`D`$ and $`D^{}`$-meson exchange, the $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}^{}`$ cross section calculated with $`D`$-meson exchange and the inclusive cross section for $`J/\mathrm{\Psi }`$ dissociation with $`\overline{D}+D`$ production.
We found that the total $`J/\mathrm{\Psi }+N`$ cross section approaches 5.5 mb at high invariant collision energy. We also indicate a partial enhancement of the $`J/\mathrm{\Psi }+N`$ dissociation cross section at low energies due to the contribution from the $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ and $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}^{}`$ reaction channels.
The line c) in Fig. 4 shows the parameterization from Ref. , explicitly given as
$$\sigma _{JN}=2.5\left(1\frac{\lambda }{\lambda _0}\right)^{6.5},$$
(33)
where the cross section is given in mb, $`\lambda _0=m_N+ϵ_0`$ with $`ϵ_0`$ being the Rydberg energy and
$$\lambda =\frac{sm_J^2m_N^2}{2m_J}.$$
(34)
Parameterization (34) provides a substantially smaller $`J/\mathrm{\Psi }`$ dissociation cross section compared to our result.
Furthermore, the $`J/\mathrm{\Psi }+N`$ cross section can be evaluated using short distance QCD methods based on the operator product expansion . Within the first order calculation the cross section is given as
$$\sigma _{JN}=\frac{2^{13}\pi }{3^4\alpha _sm_c^2}\underset{1/\xi }{\overset{1}{}}\frac{(\xi x1)^{3/2}}{(\xi x)^5}\frac{g(x)}{x}𝑑x,$$
(35)
where $`\xi =\lambda /ϵ_0`$, $`m_c`$ is the $`c`$ quark mass, $`\alpha _s`$ is the strong coupling constant and $`g(x)`$ denotes the gluon distribution function, for which we take the form
$$g(x)=2.5(1x)^4.$$
(36)
Calculations with a more realistic gluon distribution function only change the $`J/\mathrm{\Psi }+N`$ cross section slightly at invariant collision energies $`\sqrt{s}<`$20 GeV, as compared to that obtained with function (36). The differences associated with various gluon structure functions can be predominantly observed at high $`\sqrt{s}`$. The $`J/\mathrm{\Psi }+N`$ cross section from the first order calculations performed using short distance QCD is shown by the line b) in Fig. 4.
We note that the QCD results are in good agreement with our hadronic model calculations at high invariant collision energies, but substantially deviate from our predictions near the threshold for endothermic reaction. We ascribe this discrepancy to the contribution from the
$`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ and $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}^{}`$ reactions. The dashed line in Fig. 5 shows the separate contribution to the total $`J/\mathrm{\Psi }+N`$ cross section from the $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ and $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}^{}`$ reactions, while the solid line a) again shows our result for the total $`J/\mathrm{\Psi }`$ dissociation by a nucleon. The solid line b) in Fig. 5 indicates the QCD result given by Eq. (35). It is clear that the main difference between our prediction and the QCD calculation comes from the $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ and $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}^{}`$ reaction channels, which contributes substantially at low $`\sqrt{s}`$.
## VI $`J/\mathrm{\Psi }`$ photoproduction on the nucleon
Within the vector dominance model, the $`J/\mathrm{\Psi }`$ photoproduction invariant amplitude, $`_{\gamma J}`$, on a nucleon can be related to the $`J/\mathrm{\Psi }+N`$ scattering amplitude, $`_{JN}`$, as:
$$_{\gamma J}(s,t)=\frac{\sqrt{\pi \alpha }}{\gamma _J}F(t)_{JN}(s,t),$$
(37)
where $`s`$ is the squared, invariant collision energy, $`t`$ is the squared four-momentum transfer, $`\alpha `$ is the fine-structure constant and $`\gamma _J`$ is the constant for $`J/\mathrm{\Psi }`$ coupling to the photon. In Eq. (37) $`F(t)`$ stands for the form factor at the $`\gamma J/\mathrm{\Psi }`$ vertex, which accounts for the $`c\overline{c}`$ fluctuation of the photon . Within the naive VDM only the hadronic, but not the quark-antiquark, fluctuations of the photon through the $`\gamma `$ mixing with the vector mesons are considered and the form factor $`F(t)`$ in Eq. 37 is therefore neglected. As was calculated in Ref. by both the hadronic and quark representations of the $`c\overline{c}`$ fluctuation of the photon, the form factor $`F=`$0.3 at $`t=0`$.
The invariant amplitudes are normalized so that the $`J/\mathrm{\Psi }+N`$ differential cross section is written as
$$\frac{d\sigma }{dt}=\frac{||^2}{16\pi [(sm_N^2m_J^2)^24m_N^2m_J^2]},$$
(38)
and similarly for the $`J/\mathrm{\Psi }`$ photoproduction cross section.
Taking the relation of Eq. (37) at $`t=0`$ and applying the optical theorem for the imaginary part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude, one can express the total cross section for $`J/\mathrm{\Psi }`$ dissociation by a nucleon $`\sigma _{JN}`$ in terms of the cross section for $`J/\mathrm{\Psi }`$ photoproduction on the nucleon at $`t=0`$ as
$`{\displaystyle \frac{d\sigma _{\gamma NJN}}{dt}}|_{t=0}={\displaystyle \frac{\alpha }{16\gamma _J^2}}(1+\alpha _{JN}^2)\sigma _{JN}^2F^2(0)`$ (39)
$`\times {\displaystyle \frac{(sm_N^2m_J^2)^24m_N^2m_J^2}{(sm_N^2)^2}},`$ (40)
where $`\alpha _{JN}`$ is the ratio of the real to imaginary part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude at $`t=0`$.
The $`\gamma _J`$ coupling constant can be directly determined from the $`J/\mathrm{\Psi }`$ meson decay into leptons
$$\mathrm{\Gamma }(J/\mathrm{\Psi }l^+l^{})=\frac{\pi \alpha ^2}{3\gamma _J^2}\sqrt{m_J^24m_l^2}\left[1+\frac{2m_l^2}{m_J^2}\right].$$
(41)
The ratio $`\alpha _{JN}`$ is unknown and as a first approximation we put it equal to zero. This approximation is supported from two sides.
At high energies Regge theory dictates that the amplitude for hadron scattering on a nucleon is dominated by pomeron exchange and is therefore purely imaginary. This expectation is strongly supported by the experimental data available for $`p`$, $`\overline{p}`$ $`\pi `$, $`K`$ and $`\overline{K}`$ scattering by a nucleon.
At low energies the real part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude at $`t=0`$ can be estimated from the theoretical calculations of the $`J/\mathrm{\Psi }`$-meson mass shift in nuclear matter. As was predicted by the calculations by the operator product expansion , QCD van der Waals potential and QCD sum rules , the mass of the $`J/\mathrm{\Psi }`$ should only be changed a tiny amount in nuclear matter. However, a partial deviation of our results from the $`J/\mathrm{\Psi }`$ photoproduction data might be actually addressed to the contribution from a non vanishing ratio $`\alpha _{JN}`$. More detailed discussion concerning an evaluation of the real part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude will be given in the following.
The $`J/\mathrm{\Psi }`$ photoproduction cross section at $`t=0`$ is shown in Fig. 6 as a function of the invariant collision energy. The experimental data were taken from Refs. . The solid line shows the results calculated with the total $`J/\mathrm{\Psi }+N`$ cross section given by the boson exchange model. We found that within the experimental uncertainties we can reproduce the $`J/\mathrm{\Psi }`$ photoproduction data at $`t=0`$ at low energies quite well. Let us recall that this is possible because of the contribution from the $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ and $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}^{}`$ reaction channels.
The dashed line in Fig. 6 indicates the calculations using Eq. (35), which reasonably describe the data on photoproduction cross section at $`\sqrt{s}>`$10 GeV. We notice a reasonable agreement between our calculations and the results from short distance QCD at invariant collision energies $`\sqrt{s}>`$12 GeV. We also recall that in Ref. the discrepancy between the QCD results and the photoproduction data at low $`\sqrt{s}`$ was attributed to a large ratio $`\alpha _{JN}`$. Our result does not support this suggestion.
## VII Comparison with Regge theory
It is of some interest to compare our results with the predictions from Regge theory given at high energies. Furthermore, for illustrative purpose we demonstrate the comparison in terms of the cross section at $`t=0`$ for the $`\gamma +NJ/\mathrm{\Psi }+N`$ reaction, since the Regge model parameters were originally fitted to the photoproduction data.
Taking into account the contribution from soft and hard pomeron exchanges alone, the exclusive $`J/\mathrm{\Psi }`$ photoproduction cross section at $`t=0`$ can be explicitly written as
$$\frac{d\sigma }{dt}|_{t=0}=23.15s^{0.16}+0.034s^{0.88}+1.49s^{0.52},$$
(42)
where the cross section is given in nb and the squared invariant collision energy, $`s`$, is in GeV<sup>2</sup>. The first term of Eq. (42) comes from the soft pomeron contribution, the second one is due to the hard pomeron, while the last term stems from their interference. The parameters for the both pomeron trajectories were taken from the most recent fit to the $`J/\mathrm{\Psi }`$ photoproduction data from H1 and ZEUS experiments.
The prediction from Regge theory is shown by the dotted line in Fig. 7 together with experimental data on $`J/\mathrm{\Psi }`$ photoproduction cross section at $`t=0`$ as a function of the invariant collision energy $`\sqrt{s}`$. The solid line in Fig. 7 shows our calculations using the boson exchange model with form factors. We notice, that in Refs. the Regge model fit to the data on $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section at $`t=0`$ was not explicitly included in an evaluation of the coefficients of Eq. 42, which might partially explain some systematic deviation of the Regge calculations from the photoproduction data.
In order to illustrate the compatibility of the Regge model prediction with our calculations we indicate by the dashed line in Fig. 7 the result from Eq. (42), renormalized by a factor 1.3. We note a reasonable agreement between our calculations and the Regge theory within the short range of energies from $`\sqrt{s}`$10 up to 30 GeV.
However, considering the result given in this section one should keep in mind the following critical arguments.
First, we performed the calculations using a boson exchange model, which has some restrictions in its application at high energies. For instance, the form factors introduced at the interaction vertices suppress the contribution at large 4-momentum transfer $`t`$, which allows one to avoid the divergence of the total cross section at large collision energies where the range of the available $`t`$ becomes extremely large . A more accurate way to resolve such a divergence at high energies is to use the Reggeized boson exchange model .
Second, the Regge theory has limitations for applications at low energies . Here we adopt the parameters for soft and hard pomeron trajectories that were originally fixed using the large set of data at $`\sqrt{s}`$40 GeV and the very reasonable agreement between the Regge model calculations with the data, even at lower energies, might in some sense be viewed as surprising.
On the other hand, the $`D`$ and $`D^{}`$ meson exchanges used in our calculations cannot be related to pomeron exchange, rather they can be related to the Regge trajectories, whose contribution to $`J/\mathrm{\Psi }`$ photoproduction at high energies was assumed to be negligible. This does not contradict the results shown in Fig. 7, since our calculations are substantially below the data at high energies.
## VIII Evaluation of the real part of $`J/\mathrm{\Psi }+N`$ scattering amplitude
One of the crucial ways to test the coherence of our calculations consists of an evaluation of the real part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude at $`t=0`$ and a comparison to the various model predictions for $`\mathrm{}f(0)`$ at zero $`J/\mathrm{\Psi }`$ momentum. The latter is of course related to the $`J/\mathrm{\Psi }`$-meson mass shift or the real part of its potential in nuclear matter.
Within the low density theorem the $`J/\mathrm{\Psi }`$ meson mass shift $`\mathrm{\Delta }m_J`$ in nuclear matter at baryon density $`\rho _B`$ can be related to the real part of the scattering amplitude $`\mathrm{}f(0)`$ as
$$\mathrm{\Delta }m_J=2\pi \frac{m_N+m_J}{m_Nm_J}\rho _B\mathrm{}f(0).$$
(43)
The recent QCD sum rule analysis with the operator product expansion predicts attractive mass shifts of about –4$`÷`$–10 MeV for $`J/\mathrm{\Psi }`$ meson in nuclear matter, which corresponds to small $`\mathrm{}f(0)`$ about 0.1$`÷`$0.2 fm. This result is very close to the predictions from the operator product expansion and the calculations with QCD van der Waals potential .
There are two ways to link these results for the effective $`J/\mathrm{\Psi }`$ mass in nuclear matter or the $`J/\mathrm{\Psi }`$-nucleon scattering length with our calculations. Both of the methods applied below contain a number of uncertainties and should be carefully considered.
First, we address the partial discrepancy between our calculations of the $`J/\mathrm{\Psi }`$ photoproduction cross section at $`t=0`$ and data at low energies due to the nonzero ratio $`\alpha _{JN}`$ appearing in Eq. (40) and evaluate it from the experimental measurements as
$$\alpha _{JN}^2=\frac{d\sigma ^{exp}}{dt}|_{t=0}\times \left[\frac{d\sigma ^{th}}{dt}|_{t=0}\right]^11,$$
(44)
where the indices $`exp`$ and $`th`$ denote the experimental and theoretical results for the photoproduction cross section at $`t=0`$, respectively.
The modulus of the ratio $`\alpha _{JN}`$ evaluated from the data is shown in Fig. 8 as a function of the $`J/\mathrm{\Psi }`$ momentum, $`p_J`$, in the nucleon rest frame. It is clear that the sign of the ratio cannot be fixed by Eq. (44), but some insight might be gained from the Regge theory calculation, which is shown in Fig. 8 by the dashed line. Here again we use the parameters of the pomeron trajectories from Ref. . Now one can conclude that the $`\mathrm{}f(0)`$ is positive at low $`J/\mathrm{\Psi }`$ momenta, which agrees with the predictions on the reduction of the $`J/\mathrm{\Psi }`$ meson mass in nuclear matter. However, this speculation is true only if the real part of the $`J/\mathrm{\Psi }`$-nucleon scattering amplitude at $`t=0`$ does not change its sign at some moderate momentum.
Now, $`\mathrm{}f(0)`$ is given as a product of the ratio $`\alpha _{JN}`$ extracted from the photoproduction data and the imaginary part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude at $`t=0`$, which is related to the total $`J/\mathrm{\Psi }+N`$ cross section by the optical theorem (45) as
$$\mathrm{}f(0)=\frac{p_J}{4\pi }\sigma _{JN},$$
(45)
Let us recall that in our calculations the total $`J/\mathrm{\Psi }+N`$ cross section is given only for OZI allowed processes, i.e. at energies $`\sqrt{s}m_{\mathrm{\Lambda }_c}+m_D`$. This threshold corresponds to the $`J/\mathrm{\Psi }`$ meson momentum of$``$1.88 GeV/c. Thus, within the present method we cannot provide the real part of the scattering amplitude at $`p_J=0`$ and make a direct comparison with the absolute value of the $`\mathrm{}f(0)`$ predicted in Refs. . However, the estimate of $`\mathrm{}f(0)`$ at minimal $`J/\mathrm{\Psi }`$ momenta allowed by OZI reactions might be done.
The $`\mathrm{}f(0)`$ extracted from the experimental data on $`J/\mathrm{\Psi }`$ meson photoproduction cross section on a nucleon at $`t=0`$ is shown in Fig. 9 as a function of $`J/\mathrm{\Psi }`$ momentum in the nucleon rest frame. The shadowed area shows the predictions at $`p_J=0`$. Notice, that the shadowed area is extended to higher momenta only for orientation. Within the uncertainties of the method we could not detect any discrepancy with the absolute value of the $`\mathrm{}f(0)`$ given in Refs. . However, we emphasize again, that this conclusion should be accepted very carefully.
A different way to evaluate the real part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude at $`t=0`$ involves the dispersion relation that is given as
$`\mathrm{}f(\omega )=Ref(\omega _0)+{\displaystyle \frac{2(\omega ^2\omega _0^2)}{\pi }}`$ (46)
$`\times P{\displaystyle \underset{\omega _{min}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\omega ^{}\mathrm{}f(\omega ^{})}{(\omega _{}^{}{}_{}{}^{2}\omega _0^2)(\omega _{}^{}{}_{}{}^{2}\omega ^2)}}𝑑\omega ^{},`$ (47)
where $`\omega `$ stands for the $`J/\mathrm{\Psi }`$ total energy, $`\omega _{min}`$ is the threshold energy and $`\mathrm{}f(\omega )`$ denotes the imaginary part of the $`J/\mathrm{\Psi }`$-nucleon scattering amplitude at $`t=0`$, as a function of energy $`\omega `$, which can be evaluated from Eq. (45). Furthermore, $`Ref(\omega _0)`$ is the subtraction constant taken at $`J/\mathrm{\Psi }`$ energy $`\omega _0`$.
Being more fundamental, this method involves substantial uncertainties for the following reasons. First, we should specify the subtraction constant. By taking it at $`p_J=0`$ from Eq. (43) we could not further verify the consistency of our calculations with the predictions from Refs. . Second, to evaluate the principal value of the integral of Eq. 47 one should know the imaginary part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude at $`t=0`$ or the total cross section up to infinite energy.
Since our model could not provide the high energy behavior of $`\sigma _{JN}`$ one should make an extrapolation to high energies by using either short distance QCD or Regge theory. As was shown in Ref. the total $`J/\mathrm{\Psi }+N`$ cross section at high energies is dependent on the gluon distribution function, $`g(x)`$, appearing in Eq. (35). Furthermore, both distributions given by Eq. (36) and taken from Ref. were unable to reproduce the experimental data on the $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section at $`\sqrt{s}>`$70 GeV. Thus we could not extrapolate our results to high energies by the short distance QCD calculations.
As shown in Fig. 7, the Regge theory reproduces the $`J/\mathrm{\Psi }`$ photoproduction data at $`t=0`$ quite well, when renormalized by a factor of 1.3. This meets our criteria since it also fits our calculations at $`\sqrt{s}`$20 GeV. Thus we adopt the Regge model fit for the extrapolation of our results to high energies in order to evaluate the dispersion relation. It is important to notice, that, in principle, there is no need to address $`J/\mathrm{\Psi }`$ photoproduction in order to extrapolate our results to infinite energy using Regge theory. One can adopt the energy dependence of the dissociation cross section given by soft and hard pomeron exchanges and make the absolute normalization by our calculations $`\sqrt{s}>`$20 GeV. This is reasonable, because the absolute normalization from the Regge model can be considered to some extent as a free parameter, not fixed by the theory.
Moreover, we take the subtraction constant $`\mathrm{}f(\omega _0)`$ at high energy, $`\omega _0`$=10<sup>3</sup> GeV, and fix it by Regge theory with the predicted ratio $`\alpha _{JN}`$. Thus we can independently use an evaluated real part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude at $`t=0`$ for comparison with the predictions given at $`p_J`$=0.
Our results for the ratio $`\alpha _{JN}`$ evaluated by the dispersion relation (47) are shown in Fig. 8 by the solid line. Obviously, the calculations match the Regge model predictions since they were used to fixed both the high energy behavior of the total $`J/\mathrm{\Psi }+N`$ cross section and the subtraction constant for the real part of the scattering amplitude at $`t=0`$. Furthermore, the dispersion calculations give a surprisingly good fit to the ratios $`\alpha _{JN}`$ extracted from the experimental data on the cross section at $`t=0`$ of the $`\gamma +NJ/\mathrm{\Psi }+N`$ reaction.
The real part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude calculated by dispersion relation is shown by the solid line in Fig. 9 as a function of $`J/\mathrm{\Psi }`$ meson momentum in the nucleon rest frame. Our calculations describes reasonably well the data evaluated from the $`J/\mathrm{\Psi }`$ photoproduction at $`t=0`$, that are shown by the solid circles in Fig. 9. It is important that our results, normalized by a subtraction constant at high energies, $`p_J=10^3`$ GeV, simultaneously match the predictions given at $`p_J`$=0.
Concluding this section let us make a few remarks. The two methods applied to the evaluation of the real part of the $`J/\mathrm{\Psi }`$-nucleon scattering amplitude at $`t=0`$ are actually different, but nevertheless provide almost identical results for the $`\mathrm{}f(0)`$ over a large range of $`J/\mathrm{\Psi }`$ momenta. The consistency of our approach can be proved by agreement between our calculations simultaneously with the Regge model predictions for $`\mathrm{}f(0)`$ at high energy and the predictions from various models for $`\mathrm{}f(0)`$ at $`p_J=0`$. ¿From these results we confirm that below the momentum of $`p_J`$1.88 GeV, which corresponds to the OZI allowed $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ threshold, the $`J/\mathrm{\Psi }`$ interactions with a nucleon are purely elastic. The $`J/\mathrm{\Psi }`$ meson does not undergo significant dissociation on absorption at momenta $`p_J<1.88`$ GeV.
## IX Estimate of the elastic $`J/\mathrm{\Psi }+N`$ cross section
By definition, the total elastic $`J/\mathrm{\Psi }+NJ/\mathrm{\Psi }+N`$ cross section, $`\sigma _{el}`$, is related to the elastic differential cross section as
$$\sigma _{el}=\frac{d\sigma }{dt}|_{t=0}\mathrm{exp}(bt),$$
(48)
where $`b`$ is an exponential slope of the differential cross section and the cross section at the optical point, $`t=`$0, is related to the total $`J/\mathrm{\Psi }+N`$ cross section and the ratio $`\alpha _{JN}`$ of the real and imaginary part of the scattering amplitude as
$$\frac{d\sigma }{dt}|_{t=0}=\frac{1}{16\pi }(1+\alpha _{JN}^2)\sigma _{JN}^2.$$
(49)
The slope $`b`$ can be estimated from the differential $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section and was fitted in Ref. to available experimental data as $`b=1.64+0.83\mathrm{ln}s`$.
Finally, the elastic $`J/\mathrm{\Psi }+N`$ cross section is given by
$$\sigma _{el}=\frac{1}{16\pi b}(1+\alpha _{JN}^2)\sigma _{JN}^2,$$
(50)
and is shown in Fig. 10 as a function of $`J/\mathrm{\Psi }`$ meson momentum. The dashed area in Fig. 10 indicates the elastic cross section at $`p_J=0`$ given within the scattering length approximation in Refs. as
$$\sigma _{el}=4\pi |\mathrm{}f(0)|^2,$$
(51)
that is valid as far as $`\mathrm{}f(0)`$=0 at $`p_J`$=0. Note, that Eq. 50 could not provide the estimate for elastic cross section at $`p_J<`$1.88 GeV, i.e. below the $`\mathrm{\Lambda }_c+\overline{D}`$ reaction threshold, where the dissociation cross section $`\sigma _{JN}`$=0. In principle, one can interpolate $`\sigma _{el}`$ from zero $`J/\mathrm{\Psi }`$ momentum to our results.
As we already discussed, the uncertainty in the elastic cross section at $`p_J`$=0, given in Refs. , is very large and to make a less unambiguous interpolation let us to use only the most recent results from the QCD sum rule calculations . Let us recall that the QCD sum rule calculations from Ref. indicate a shift of the $`J/\mathrm{\Psi }`$ mass of about –5$`÷`$–10 MeV, while in Ref. the mass shift was reported to be in the range –$`7\mathrm{\Delta }m_J`$–4 MeV. The calculations with higher dimension operators in the QCD sum rule provide $`\mathrm{\Delta }m_J`$=–4 MeV. Thus, we take an average value which corresponds to $`\sigma _{el}`$ =1.6 mb at $`p_J`$=0. Now the dashed line in Fig. 10 indicates our interpolation of the elastic cross section at low momenta.
Furthermore, the vector dominance model can be used to relate the elastic $`J/\mathrm{\Psi }+NJ/\mathrm{\Psi }+N`$ cross section $`\sigma _{el}`$ and the total $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section $`\sigma _{\gamma J}`$. Applying Eq. (37) and neglecting the form factor this relation can be written as
$$\sigma _{el}=\frac{\gamma _J^2}{\pi \alpha }\frac{(sm_J^2)^2}{(sm_N^2m_J^2)^24m_N^2m_J^2}\sigma _{\gamma J}$$
(52)
Fig 11a) shows the $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section as a function of invariant collision energy, while the elastic cross section $`\sigma _{el}`$ evaluated from experimental photoproduction data using Eq. (52) is shown in Fig 11b). Apart of the absolute value of the $`\sigma _{el}`$, which is substantially below our estimate, we notice that very strong energy dependence of the elastic $`J/\mathrm{\Psi }+NJ/\mathrm{\Psi }+N`$ cross section evaluated using Eq. (52) might alone indicate the inconsistency of the application of naive VDM . Indeed, as we discussed above, the calculations of the $`J/\mathrm{\Psi }`$ mass shift in matter or $`J/\mathrm{\Psi }+N`$ scattering length given in Refs. predicts the elastic cross section in the range $`1.25\sigma _{el}`$5 mb at $`p_J`$=0, while the $`J/\mathrm{\Psi }`$ meson photoproduction data evaluated by Eq. (52) indicate a systematic decrease of $`\sigma _{el}`$ with decrease of $`J/\mathrm{\Psi }`$ meson momentum.
Apparently this inconsistency might be resolved by introducing a form factor $`F(t)`$ at the $`\gamma J/\mathrm{\Psi }`$ vertex , as given by Eq. (37) and applied in our study. The correction to the naive vector dominance is then given as
$$\kappa =\left[\mathrm{exp}(bt)\right]^1F^2(t)\mathrm{exp}(bt),$$
(53)
where $`b`$ is the slope of the differential $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section and the integration in Eq. (53) is performed over the range of the squared transverse momentum $`t`$ available at given invariant collision energy $`\sqrt{s}`$.
By taking, for simplicity, an exponential form of the $`\gamma J/\mathrm{\Psi }`$ form factor with cutoff parameter $`\mathrm{\Lambda }`$=1.7 GeV<sup>-2</sup>, we evaluate the elastic $`J/\mathrm{\Psi }+NJ/\mathrm{\Psi }+N`$ cross section from the data on the $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section and show the result by the circles in Fig. 12, as a function of the $`J/\mathrm{\Psi }`$ meson momentum. The solid line in Fig. 12 indicates our estimate for $`\sigma _{el}`$, which clearly give a reasonable fit to the correctly extracted data over a large range of available momenta $`p_J`$.
It is worthwhile to note that Fig. 12 clearly illustrates that the introduction of a form factor at the $`J/\mathrm{\Psi }N`$ vertex allows one to resolve an inconsistency in the evaluation the elastic $`J/\mathrm{\Psi }+NJ/\mathrm{\Psi }+N`$ cross section from the total $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section.
## X Final results
Finally, Fig. 13 shows the $`J/\mathrm{\Psi }+N`$ dissociation and elastic cross sections over a large range of $`J/\mathrm{\Psi }`$ momenta. Let us recall that the dissociation cross section, $`\sigma _{JN}`$, was calculated using the boson exchange model up to $`p_J200`$ Gev/c and extrapolated to high energies by Regge theory.
Now we compare our results with the $`J/\mathrm{\Psi }+N`$ dissociation cross section evaluated from nuclear reactions. The square in Fig. 13 shows the $`\sigma _{JN}`$ extracted from the $`J/\mathrm{\Psi }`$ photoproduction from nuclei at a mean photon energy of 17 GeV . The $`J/\mathrm{\Psi }`$ meson momentum was not explicitly fixed and is indicated in Fig. 13 only for orientation.
Furthermore, the circle indicates the result from the combined analysis of experimental data on $`J/\mathrm{\Psi }`$ production from $`p+A`$ collisions at beam energies from 200 to 800 GeV. Again the $`J/\mathrm{\Psi }`$ meson momentum was not explicitly fixed by the analysis .
Within the experimental uncertainties our results for the total $`J/\mathrm{\Psi }+N`$ dissociation cross section fit the available experimental data reasonably well.
We notice, that our results slightly ($``$15%) underestimate the $`J/\mathrm{\Psi }+N`$ dissociation cross section evaluated from experimental data on $`J/\mathrm{\Psi }`$ production from $`\gamma +A`$ and $`p+A`$ reactions.
We also note that the NA50 experimental data on $`J/\mathrm{\Psi }`$ production from heavy ion collisions can be reasonably well fitted by the calculations with a $`J/\mathrm{\Psi }+N`$ dissociation cross section $`\sigma _{JN}3÷6.7`$ mb. The most recent calculations on $`J/\mathrm{\Psi }`$ production and comparison to NA50 data indicates a smaller dissociation cross section, in the range from 3 to 4.5 mb. In central $`Pb+Pb`$ collisions at 160 A$``$GeV the $`J/\mathrm{\Psi }`$ meson momenta $`p_J`$ in the nucleon rest frame are distributed over the range from 15 to 70 GeV/c with the maximum at $`p_J`$40 GeV/c . Comparing the heavy ion results with our calculations from Fig. 11 we also find reasonable agreement with $`A+A`$ data.
Furthermore, the nonperturbative QCD calculations provides the $`J/\mathrm{\Psi }+N`$ dissociation cross section of $`4.4\pm 0.6`$ mb at $`J/\mathrm{\Psi }`$ meson momenta above 200 GeV, which is consistent with our results shown by Fig. 13.
## XI Future perspectives
There are a few problems that remain open and need further investigation.
In the present study we considered only OZI allowed processes resulting in an endothermic inelastic reaction, where the total mass of the produced particles is larger than the initial mass, $`m_J+m_N`$. These reactions provide the threshold behavior of the cross section. On the other hand, the OZI suppressed reaction channels such as
$`J/\mathrm{\Psi }+NN+n\pi `$ with the number of the final pions from $`n`$=1 up to $`m_J/m_\pi 22`$ are open even at $`p_J`$=0. These exothermic reactions, where the total final mass of the produced particles is less that the initial mass, provide a $`1/p_J`$ behavior of the $`J/\mathrm{\Psi }+N`$ dissociation cross section.
For instance, the OZI suppressed $`J/\mathrm{\Psi }+NN+\pi `$ cross section was calculated by the $`\rho `$ meson exchange model . It was found that the cross section is negligibly small although it diverges at $`p_J`$ close to zero, because it is exothermic. Furthermore, the $`\rho `$ meson exchange model evaluates the $`J/\mathrm{\Psi }\rho +\pi `$ vertex, which accounts only tiny fraction, $``$1.27%, for the total hadronic decays of the $`J/\mathrm{\Psi }`$ meson. The multi-meson OZI suppressed reactions might play a non negligible role for $`J/\mathrm{\Psi }`$ dissociation on a nucleon. Eventually additional theoretical estimates are necessary in order to make a final conclusion about the possibility of $`J/\mathrm{\Psi }+N`$ dissociation at low momenta.
Furthermore, we found a strong momentum dependence of the $`J/\mathrm{\Psi }+N`$ dissociation cross section, as illustrated by Fig. 13. This momentum dependence of $`\sigma _{JN}`$ might provide a partial explanation of the variation of the slope $`\alpha `$ from the $`A^\alpha `$-dependence with the Feynman variable $`x_F`$. As was observed in proton-nucleus collisions at beam energy of 800 GeV the extracted value for $`\alpha `$ is $``$0.9 at $`x_F`$=0.2, while it decreases $`0.8`$ at $`x_F`$0.7. The Feynman variable $`x_F`$ is proportional to the $`J/\mathrm{\Psi }`$ meson momentum and one may naturally expect from the results shown in Fig. 13 that at large momenta, $`p_J`$, or equivalently large $`x_F`$, the value for $`\alpha `$ extracted from the $`A^\alpha `$-dependence, should be smaller in comparison to small $`p_J`$. Now we can qualitatively predict that the slope $`\alpha `$ decreases with increasing $`x_F`$. However, a quantitatively analysis needs further calculation.
An alternative way is to reanalyze the experimental data on the $`A`$-dependence of $`J/\mathrm{\Psi }`$ production from $`p+A`$ collisions and to evaluate the dissociation cross section as a function of the $`J/\mathrm{\Psi }`$ momentum. This will provide a crucial test of our calculations, since it can be directly compared to results shown in Fig. 13.
Furthermore, the ratio $`\alpha _{JN}`$ of the real to imaginary part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude at $`t=0`$, which is shown by Fig. 8, in principle can be measured using coherent $`J/\mathrm{\Psi }`$ production at diffractive minima . Coherent $`J/\mathrm{\Psi }`$ photoproduction from nuclei seems to be an optimal way for such a measurement, however one needs to perform further investigations in order to fix the sensitivity of the data to the sign and magnitude of the ratio $`\alpha _{JN}`$.
Obviously, the form factor $`F(t)`$ at the $`\gamma J/\mathrm{\Psi }`$ vertex plays a key role for the comparison of our hadronic calculations with the data on $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section. While this form factor is theoretically motivated at $`t=`$0 by the microscopic calculations given in Ref. , the evaluation of its $`t`$-dependence still need further studies. Our results shown in Fig. 12 should be consider as a phenomenological estimate for $`F(t)`$.
In addition, one can address the effect due to in-medium modification of the $`J/\mathrm{\Psi }+N`$ dissociation cross section because of the strong changes of the $`D`$ and $`D^{}`$ meson properties in nuclear matter . As was found in Ref. this modification plays a substantial role for $`J/\mathrm{\Psi }`$ dissociation on comovers. One might expect that the modification of the $`J/\mathrm{\Psi }+N`$ dissociation cross section in nuclear matter might play no role for heavy ion collisions where the available $`J/\mathrm{\Psi }`$ momenta given in the nucleon rest frame are larger than $`p_J>`$13 GeV, as was discussed above. However, to clarify the situation a detailed calculation of the in-medium modification of the $`J/\mathrm{\Psi }+N`$ dissociation amplitude is necessary.
## XII Summary
We have calculated the $`J/\mathrm{\Psi }+N`$ dissociation cross section by a boson exchange model, including $`D`$ and $`D^{}`$ meson exchange and considering the $`\mathrm{\Lambda }_c+\overline{D}`$ and $`N+D+\overline{D}`$ final states. We note that our results are in reasonable agreement with short distance QCD calculations at high energies, while differing from them at low energies because of the $`J/\mathrm{\Psi }+N\mathrm{\Lambda }_c+\overline{D}`$ reaction channel explicitly included in our model.
To compare our results with the data on $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section we introduced form factor at the $`\gamma J/\mathrm{\Psi }`$ vertex proposed in Ref. and found good agreement with the photoproduction data.
Furthermore, we evaluated the real part of the $`J/\mathrm{\Psi }+N`$ scattering amplitude at $`t=0`$, $`f(0)`$, using a dispersion relations and extrapolating our result for the dissociation cross section to infinite energy by the Regge theory . Fixing the subtraction point for $`\mathrm{}f(0)`$ at large energies with Regge theory, we were able to simultaneously saturate the $`\mathrm{}f(0)`$ at zero momentum of $`J/\mathrm{\Psi }`$ meson given by the predictions for the $`J/\mathrm{\Psi }`$ mass shift in nuclear matter and for the $`J/\mathrm{\Psi }+N`$ scattering length .
We estimate the elastic $`J/\mathrm{\Psi }+NJ/\mathrm{\Psi }+N`$ cross section from our calculations and illustrate the compatibility of our results with the data on the total $`\gamma +NJ/\mathrm{\Psi }+N`$ cross section. It is important to note that the $`J/\mathrm{\Psi }`$ photoproduction data were used in our study mostly for illustrative purposes and do not influence the parameters of our calculations or the final results.
Finally, we predict the energy dependence of the $`J/\mathrm{\Psi }+N`$ dissociation, $`\sigma _{JN}`$, and the elastic cross section over a large range of energies, from 1 to 10<sup>3</sup> GeV. In contrast to the usual expectation of a constant dissociation cross section, we found that $`\sigma _{JN}`$ varies strongly over the indicated range of energies. Our results are in agreement with the $`\sigma _{JN}`$ evaluated phenomenologically from the experimental data on $`J/\mathrm{\Psi }`$ meson production in $`\gamma +A`$, $`p+A`$ and $`A+A`$ collisions.
###### Acknowledgements.
A.S would like to acknowledge the warm hospitality and partial support of the CSSM during his visit. The discussions with E. Bratkovskaya, W. Cassing, C. Greiner, D. Kharzeev and J. Speth are appreciated. This work was supported by the Australian Research Council and the Forschungzentrum Jülich.
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# About universality of lifetime statistics in quantum chaotic scattering
<sup>1</sup><sup>1</sup>institutetext: Fachbereich Physik, Universität Kaiserslautern, D-67653 Kaiserslautern, Germany
## Abstract
The statistics of the resonance widths and the behavior of the survival probability is studied in a particular model of quantum chaotic scattering (a particle in a periodic potential subject to static and time-periodic forces) introduced earlier in Ref. PRA1 ; PRL1 . The coarse-grained distribution of the resonance widths is shown to be in good agreement with the prediction of Random Matrix Theory (RMT). The behavior of the survival probability shows, however, some deviation from RMT.
journal: Physica E, in press
1. The abstract Random Matrix Theory (RMT) is known as a powerful tool for analyzing complex quantum systems. During the last two decades the predictions of the hermitian RMT (which is supposed to describe the properties of a closed system) were checked for a large number of physical models and an understanding of the conditions of applicability of RMT was reached. The situation is different, however, for nonhermitian RMT, which is aimed to describe the spectral properties of open systems. Here we have quite a few physical models which can serve a suitable test for a nonhermitian RMT. Up to our knowledge the following systems are mainly under discussion: chaotic 2D billiards with attached leads billiard ; the kicked rotor with an absorbing boundary condition rotor ; simplified models of a dissociating molecule molecule ; scattering on graphs graph ; and the Bloch particle affected by static and time periodic forces (1D model of a crystal electron in dc and ac fields) PRA1 ; PRL1 ; PRE2 . The latter model has a number of nice features, which distinguishes it among the other systems. First, it is simple for numerical analysis. Second, it always realizes the so-called case of perfect coupling (where the predictions of the hermitian and nonhermitian RMT are most different). Third, as a physical model it can be and has been studied under laboratory conditions raizen . This present paper continues our study of the Bloch particle in ac and dc fields in relation to RMT. In particular we discuss here the decay of the probability in the system (probability leakage), which is the simplest quantity measured in the laboratory experiment.
2. We briefly recall some of the results of our preceding papers PRA1 ; PRL1 ; PRE2 . After an appropriate rescaling the Hamiltonian of the system of interest can be presented in the dimensionless form
$$\widehat{H}=\frac{\widehat{p}^2}{2}+\mathrm{cos}x+Fx+F_\omega x\mathrm{cos}(\omega t).$$
(1)
The parameters of the system (1) are the amplitude of the static force $`F`$, the amplitude $`F_\omega `$ and the frequency $`\omega `$ of the time-periodic force, and the scaled Planck constant $`\mathrm{}`$ which enters the momentum operator $`\widehat{p}=i\mathrm{}\mathrm{d}/\mathrm{d}x`$. Another, unitary equivalent, form of the Hamiltonian (1) reads as
$$\widehat{H}=\frac{\widehat{p}^2}{2}+\mathrm{cos}[xϵ\mathrm{cos}(\omega t)]+Fx,ϵ=\frac{F_\omega }{\omega ^2}.$$
(2)
When some condition (based on Chirikov’s overlap criterion) on $`ϵ`$ is satisfied, the classical dynamics of the system (2) is an example of chaotic scattering. In fact, one of the main characteristics of the classical chaotic scattering is the delay or dwell time
$$\tau =\underset{p_0\mathrm{}}{lim}[\tau (p_0p_0)2p_0/F].$$
(3)
(Here $`\tau (p_0p_0)`$ is the time taken by the particle to change its initial momentum $`p_0`$ to the opposite one.) Figure 1 shows the delay time (3) as a function of the initial coordinate $`x`$. The fractal behavior is typical for chaotic scattering.
We proceed with the quantum case. It was shown in Ref. PRL1 ; PRE2 that the dynamics and spectral properties of the system (2) depend crucially on the condition of commensurability between the so-called Bloch period $`T_B=\mathrm{}/F`$ and the period $`T_\omega =2\pi /\omega `$ of the driving force,
$$\frac{T_B}{T_\omega }=\frac{r}{q}.$$
(4)
Providing the condition (4) is satisfied, the complex quasienergy spectrum of the system (the resonances) is given by the eigenvalues of the following nonunitary matrix
$$U_{sys}=\left(\begin{array}{cc}0_{M\times N}& 0_{M\times M}\\ W_{N\times N}& 0_{N\times M}\end{array}\right).$$
(5)
In Eq. (5) $`W_{N\times N}`$ is the unitary matrix with the coefficients
$$W_{n^{},n}=n^{}|\mathrm{exp}(ikx)\widehat{W}\mathrm{exp}(ikx)|n,$$
(6)
$$\widehat{W}=\widehat{\mathrm{exp}}\left\{\frac{i}{\mathrm{}}_0^T\left[\frac{(\widehat{p}Ft+\mathrm{}k)^2}{2}+\widehat{V}\right]dt\right\},$$
$$\widehat{V}=\mathrm{cos}[xϵ\mathrm{cos}(\omega t)],$$
and $`0_{M\times N}`$, $`0_{M\times M}`$, $`0_{N\times M}`$ are blocks of zeros. The nonunitary matrix (5) can be thought of as the truncated \[to the size $`(N+M)\times (N+M)`$\] unitary matrix of the system evolution operator over the common period $`T=qT_B=rT_\omega `$. Then the parameter $`M`$, which plays the role of the number of channels, is identical with the integer $`q`$ in condition (4). The additional parameter $`N`$ measures the number of states supported by the chaotic component of the classical phase space. (Unimportant for our present aim is the quasimomentum $`k`$ which can take any value in the interval $`\pi /qk<\pi /q`$.)
The main conjecture made in Ref. PRE2 is that the spectral statistics of the system (2) \[i.e., the eigenvalues statistics of the matrix $`U_{sys}`$\] is the same as the statistics of the eigenvalues of a random matrix $`U_{ran}`$ of the structure (5) but with the matrix $`W_{N\times N}`$ substituted by a member of the Circular Unitary Ensemble (CUE)
$$W_{N\times N}A_{N\times N},A_{N\times N}\mathrm{CUE}.$$
(7)
In what follows we examine this conjecture in more detail.
3. First we discuss the spectral statistics of the random matrix $`U_{ran}`$. There is strong numerical evidence that the statistics of the eigenvalues $`\mathrm{exp}(i)=\mathrm{exp}(iE\mathrm{\Gamma }/2)`$ of the matrix $`U_{ran}`$ is given by the universal distribution derived in Ref. fyodorov . (An analytical proof of this result is still an open problem. <sup>1</sup><sup>1</sup>1After the paper was submitted we learned that this result has been proved in Ref. zyczkowski .) In particular, the distribution of the scaled resonance width $`\mathrm{\Gamma }_s=\pi \mathrm{\Gamma }/\mathrm{\Delta }`$ ($`\mathrm{\Delta }=2\pi /N`$ is the mean level spacing) obeys the equation
$$\mathrm{\Pi }_M(\mathrm{\Gamma }_s)=\frac{(1)^M\mathrm{\Gamma }_s^{M1}}{(M1)!}\frac{d^M}{d\mathrm{\Gamma }_s^M}\left(e^{\mathrm{\Gamma }_s}\frac{\mathrm{sinh}\mathrm{\Gamma }_s}{\mathrm{\Gamma }_s}\right).$$
(8)
As an example Fig. 2 (adopted from paper PRE2 ) shows the histogram for the distribution of $`\mathrm{\Gamma }_s`$ for $`N=41`$ and $`M=1`$. We note that $`\mathrm{\Pi }_M(\mathrm{\Gamma }_s)M/\mathrm{\Gamma }_s^2`$ for $`\mathrm{\Gamma }_s1`$ and, thus, the notion of mean resonance width is not defined.
Let us discuss now the decay of the probability $`P(t)`$, which is given by the following equations
$$P(t)=|𝐜_t|^2,𝐜_{t+1}=U_{ran}𝐜_t,|𝐜_0|=1.$$
(9)
It is obvious that the dynamics of $`P(t)`$ is determined by the spectrum of the system. Thus the behavior of $`P(t)`$ suggests an additional test of the eigenvalue statistics alhassid . The main advantage of studying $`P(t)`$ is that this quantity is more easily measured in the laboratory experiments fractal .
The problem of probability decay was considered (although within a different RMT model) in paper sokolov . It was proved there that the function $`P(t)`$ has an exponential short-term and an algebraic long-term asymptotic. Adopting these results to the present model (9) we obtain
$$P(t)=\{\begin{array}{cc}\hfill \mathrm{exp}\left(\frac{M}{N}t\right)& \hfill ,tt^{}\\ \hfill \left(\frac{2t}{N}\right)^M& \hfill ,tt^{}\end{array},$$
(10)
where $`t^{}N/2=\pi /\mathrm{\Delta }`$ is of order of the Heisenberg time. The results of a numerical simulation of the dynamics of $`P(t)`$ depicted in Fig. 3 well support the analytical expression (10).
It is interesting to note that Eq. (10) can be obtained by using rather simple arguments. In fact, expanding the initial vector $`𝐜_0`$ over the set of eigenvectors of the matrix $`U_{ran}`$ and following Ref. sokolov using the diagonal approximation we have
$$P(t)=_0^{\mathrm{}}\mathrm{\Pi }_M(\mathrm{\Gamma }_s)\mathrm{exp}(\mathrm{\Gamma }t)d\mathrm{\Gamma }_s.$$
(11)
For the long-term asymptotic only the narrow resonances are of importance. Then, substituting $`\mathrm{\Pi }_M(\mathrm{\Gamma }_s)`$ by its asymptotic expression $`\mathrm{\Pi }_M(\mathrm{\Gamma }_s)\mathrm{\Gamma }_s^{M1}`$, $`\mathrm{\Gamma }_s1`$ \[see Eq. (8)\], we obtain $`P(t)(2t/N)^M`$. This power law decay takes place only for a “coherent” evolution of the initial state. In contrast, the short-term asymptotic of $`P(t)`$ coincides with an “incoherent” evolution, which would take place if one used uncorrelated matrices $`U_{ran}`$ in Eq. (9) at each time step. Then, as follows from the structure of the matrix $`U_{ran}`$, at each time step the state vector decreases its norm by the factor $`N/(N+M)`$ and
$$P(t)=\left(\frac{N}{N+M}\right)^t\mathrm{exp}\left(\frac{M}{N}t\right).$$
(12)
4. We proceed with the statistics of the resonances for the physical model (2). Numerically we find them as the eigenvalues of the matrix (5), where the parameter $`M`$ is identical with the denominator $`q`$ in the condition of commensurability (4).
Figure 4 compares the distribution of the scaled resonance widths $`\mathrm{\Gamma }_s=\pi \mathrm{\Gamma }/\mathrm{\Delta }`$ in the system (2) against the prediction of RMT given by Eq. (8). It seen that the global features of the distribution $`\mathrm{\Pi }(\mathrm{\Gamma }_s)`$ fit well to the result of RMT. (The peak-like peculiarities of $`\mathrm{\Pi }(\mathrm{\Gamma }_s)`$ are due to the resonances associated with the stability islands of the classical phase space. In principle these resonances should be removed from the analyzed data.)
We would like to note that to satisfy the condition (4) we adjusted the amplitude of the static force $`F`$ (the other system parameters are kept fixed). By changing $`F`$ we actually change the classical properties of the system \[in particular, the distribution of the classical delay time (3)\]. Nevertheless, the quantum distribution $`\mathrm{\Pi }(\mathrm{\Gamma }_s)`$ remains practically unchanged (see cases $`F=0.0095`$ and $`F=0.0663`$, for example) and is defined exclusively by the channel number $`q`$. This fact clearly demonstrates the applicability of RMT for the system under study.
We come to the problem of the probability decay. In our numerical analysis of the system we calculated the dynamics of probability $`P(t)`$ by two different methods. The first method utilizes Eq. (9) where the random matrix $`U_{ran}`$ is substituted by the matrix (5). The second method is the direct numerical simulation of the wave-packet dynamics of the system (2). The latter method has the advantage that it allows to study the incommensurate case but it is essentially more time consuming. In the commensurate case $`T_B/T_\omega =r/q`$ (with relatively small $`r`$ and $`q`$) both methods give the same result.
Figure 5 shows the behavior of $`P(t)`$ on a double logarithmic scale for $`r/q=1`$, $`3/2`$, $`4/3`$, and $`5/4`$. It is seen that the survival probability follows asymptotically a power law $`P(t)t^\alpha `$. However, the value of $`\alpha `$ ($`\alpha 1`$, $`4/3`$, $`5/3`$, $`2`$, respectively) differs from that predicted by Eq. (10). This means that for very small resonance width (not resolved on the scale of Fig. 4) the actual distribution $`\mathrm{\Pi }(\mathrm{\Gamma }_s)`$ deviates from the distribution (8). For the moment we have no explanation for this deviation from RMT.
It is also worth to note that the algebraic decay discussed is actually a transient phenomenon for physical systems. The point is that RMT deals with an infinite ensemble while in practice it is always finite (and often consists of a single representative). For a finite ensemble most narrow resonance exists and, thus, a very far asymptotic is again an exponential decay with the increment given by the width of this most narrow resonance.
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# On Modular Termination Proofs of General Logic Programs
## 1 Introduction
It is standard practice to tackle a large proof by decomposing it into more manageable pieces (lemmata or modules) and proving them separately. By composing appropriately these simpler results, one can then obtain the final proof. This methodology has been recognized an important one also when proving termination of logic programs. Moreover most practical logic programs are engineered by assembling different modules and libraries, some of which might be pre-compiled or written in a different programming language. In such a situation, a compositional methodology for proving termination is of crucial importance.
The first approach to modular termination proofs of logic programs has been proposed by Apt and Pedreschi in (?). It extends the seminal work on *acceptable* programs (?) which provides an algebraic characterization of programs terminating under Prolog left-to-right selection rule. The class of acceptable programs contains programs which terminate on ground queries. To prove acceptability one needs to determine a measure on literals (*level mapping*) such that, in any clause, the measure of the head is greater than the measure of each body literal. This implies the decreasing of the measure of the literals resolved during any computation starting from a ground or *bounded* query and hence termination.
The significance of a modular approach to termination of logic programs has been recognized also by other authors; more recent proposals can be found in (?????).
All previous proposals (with the exception of (??)) require the existence of a relation between the level mappings used to prove acceptability of distinct modules. This is not completely satisfactory: it would be nice to be able to put together modules which were independently proved terminating, and be sure that the resulting program is still terminating.
We propose a modular approach to termination which allows one to reason independently on each single module and get a termination result on the whole program. We consider general logic programs, i.e., logic programs with negation, employing SLDNF-resolution together with the leftmost selection rule (also called *LDNF-resolution*) as computational mechanism. We consider programs which can be divided into modules in a hierarchical way, so that each module is an extension of the previous ones. We show that in this context the termination proof of the entire program can be given in terms of separate proofs for each module, which are naturally much simpler than a proof for the whole program. While assuming a hierarchy still allows one to tackle most real-life programs, it leads to termination proofs which, in most cases, are extremely simple.
We characterize the class of queries terminating for the whole program by introducing a new notion of boundedness, namely *strong boundedness*. Intuitively, strong boundedness captures the queries which preserve (standard) boundedness through the computation. By proving acceptability of each module wrt. a level mapping which measures only the predicates defined in that module, we get a termination result for the whole program which is valid for any strongly bounded query. Whenever the original program is decomposed into a hierarchy of small modules, the termination proof can be drastically simplified with respect to previous modular approaches. Moreover strong boundedness can be naturally guaranteed by common persistent properties of programs and queries, namely properties preserved through LDNF-resolution such as *well-modedness* (?) or *well-typedness* (?).
The paper is organized as follows. Section 2 contains some preliminaries. In particular we briefly recall the key concepts of LDNF-resolution, acceptability, boundedness and program extension. Section 3 contains our main results which show how termination proofs of separate programs can be combined to obtain proofs of larger programs. In particular we define the concept of strongly bounded query and we prove that for general programs composed by a hierarchy of $`n`$ modules, each one independently acceptable wrt. its own level mapping, any strongly bounded query terminates. In Section 4 we show how strong boundedness is naturally ensured by some program properties which are preserved through LDNF-resolution such as well-modedness and well-typedness. In Section 5 we show how these properties allow us to apply our general results also for proving termination of modular programs in an iterative way. In Section 6 we compare our work with Apt and Pedreschi’s approach. Other related works and concluding remarks are discussed in Section 7.
## 2 Preliminaries
We use standard notation and terminology of logic programming (???). Just note that general logic programs are called in (?) normal logic programs.
### 2.1 General Programs and LDNF-Resolution
A *general clause* is a construct of the form
$$HL_1,\mathrm{},L_n$$
with $`(n0)`$, where $`H`$ is an atom and $`L_1,\mathrm{},L_n`$ are literals (i.e., either atoms or the negation of atoms). In turn, a *general query* is a possibly empty finite sequence of literals $`L_1,\mathrm{},L_n`$, with ($`n0`$). A *general program* is a finite set of general clauses<sup>1</sup><sup>1</sup>1In the examples through the paper, we will adopt the syntactic conventions of Prolog so that each query and clause ends with the period “.” and “$``$” is omitted in the unit clauses.. Given a query $`Q:=L_1,\mathrm{},L_n`$, a *non-empty prefix of* $`Q`$ is any query $`L_1,\mathrm{},L_i`$ with $`i\{1,\mathrm{},n\}`$. For a literal $`L`$, we denote by $`rel(L)`$ the predicate symbol of $`L`$.
Following the convention adopted in (?), we use bold characters to denote sequences of objects (so that L indicates a sequence of literals $`L_1,\mathrm{},L_n`$, while t indicates a sequence of terms $`t_1,\mathrm{},t_n`$).
For a given program $`P`$, we use the following notations: $`B_P`$ for the Herbrand base of $`P`$, $`ground(P)`$ for the set of all ground instances of clauses from $`P`$, $`comp(P)`$ for the Clark’s completion of $`P`$ (?).
Since in this paper we deal with general queries, clauses and programs, we omit from now on the qualification “general”, unless some confusion might arise.
We consider *LDNF-resolution*, and following Apt and Pedreschi’s approach in studying the termination of general programs (?), we view LDNF-resolution as a top-down interpreter which, given a general program $`P`$ and a general query $`Q`$, attempts to build a search tree for $`P\{Q\}`$ by constructing its branches in parallel. The branches in this tree are called *LDNF-derivations* of $`P\{Q\}`$ and the tree itself is called *LDNF-tree* of $`P\{Q\}`$. Negative literals are resolved using the negation-as-failure rule which calls for the construction of a *subsidiary LDNF-tree*. If during this subsidiary construction the interpreter diverges, the (main) LDNF-derivation is considered to be infinite. An LDNF-derivation is finite also if during its construction the interpreter encounters a query with the first literal being negative and non-ground. In such a case we say that the LDNF-derivation flounders.
By termination of a general program we actually mean termination of the underlying interpreter. Hence in order to ensure termination of a query $`Q`$ in a program $`P`$, we require that all LDNF-derivations of $`P\{Q\}`$ are finite.
By an *LDNF-descendant* of $`P\{Q\}`$ we mean any query occurring during the LDNF-resolution of $`P\{Q\}`$, including $`Q`$ and all the queries occurring during the construction of the subsidiary LDNF-trees for $`P\{Q\}`$.
For a non-empty query $`Q`$, we denote by $`\mathrm{𝑓𝑖𝑟𝑠𝑡}(Q)`$ the first literal of $`Q`$. Moreover we define $`\mathrm{𝐶𝑎𝑙𝑙}_P(Q)=\{\mathrm{𝑓𝑖𝑟𝑠𝑡}(Q^{})|Q^{}\text{ is an LDNF-descendant of }P\{Q\}\}`$. It is worth noting that if $`\neg A\mathrm{𝐶𝑎𝑙𝑙}_P(Q)`$ and $`A`$ is a ground atom, then $`A\mathrm{𝐶𝑎𝑙𝑙}_P(Q)`$ too. Notice that, for definite programs, the set $`\mathrm{𝐶𝑎𝑙𝑙}_P(Q)`$ coincides with the call set $`\mathrm{𝐶𝑎𝑙𝑙}(P,\{Q\})`$ in (??).
The following trivial proposition holds.
###### Proposition 1
Let $`P`$ be a program and $`Q`$ be a query. All LDNF-derivations of $`P\{Q\}`$ are finite iff for all positive literals $`A\mathrm{𝐶𝑎𝑙𝑙}_P(Q)`$, all LDNF-derivations of $`P\{A\}`$ are finite.
### 2.2 Acceptability and Boundedness
The method we are going to use for proving termination of modular programs is based on the concept of acceptable program (?). In order to introduce it, we start by the following definition, originally due to (?) and (?).
###### Definition 2 (Level Mapping)
A *level mapping* for a program $`P`$ is a function $`||:B_P𝐍`$ of ground atoms to natural numbers. By convention, this definition is extended in a natural way to ground literals by putting $`|\neg A|=|A|`$. For a ground literal $`L`$, $`|L|`$ is called the *level* of $`L`$.
We will use the following notations. Let $`P`$ be a program and $`p`$ and $`q`$ be relations. We say that $`p`$ *refers to* $`q`$ if there is a clause in $`P`$ that uses $`p`$ in its head and $`q`$ in its body; $`p`$ *depends on* $`q`$ if $`(p,q)`$ is in the reflexive, transitive closure of the relation *refers to*. We say that $`p`$ and $`q`$ are *mutually recursive* and write $`pq`$, if $`p`$ depends on $`q`$ and $`q`$ depends on $`p`$. We also write $`pq`$, when $`p`$ depends on $`q`$ but $`q`$ does not depend on $`p`$.
We denote by $`Neg_P`$ the set of relations in $`P`$ which occur in a negative literal in a clause of $`P`$ and by $`Neg_P^{}`$ the set of relations in $`P`$ on which the relations in $`Neg_P`$ depend. $`P^{}`$ denotes the set of clauses in $`P`$ defining a relation of $`Neg_P^{}`$.
In the sequel we refer to the standard definition of model of a program and model of the completion of a program, see (??) for details. In particular we need the following notion of complete model for a program.
###### Definition 3 (Complete Model)
A model $`M`$ of a program $`P`$ is called *complete* if its restriction to the relations from $`Neg_P^{}`$ is a model of $`comp(P^{})`$.
Notice that if $`I`$ is a model of $`comp(P)`$ then its restriction to the relations in $`Neg_P^{}`$ is a model of $`comp(P^{})`$; hence $`I`$ is a complete model of $`P`$.
The following notion of acceptable program was introduced in (?). Apt and Pedreschi proved that such a notion fully characterizes left-termination, namely termination wrt. any ground query, both for definite programs and for general programs which have no LDNF-derivations which flounder.
###### Definition 4 (Acceptable Program)
Let $`P`$ be a program, $`||`$ be a level mapping for $`P`$ and $`M`$ be a complete model of $`P`$. $`P`$ is called *acceptable wrt. $`||`$ and $`M`$* if for every clause $`A𝐀,B,𝐁`$ in $`ground(P)`$ the following implication holds:
$$\text{ if }M𝐀\text{ then }|A|>|B|.$$
Note that if $`P`$ is a definite program, then both $`P^{}`$ and $`Neg_P^{}`$ are empty and $`M`$ can be any model of $`P`$.
We also need the notion of bounded atom.
###### Definition 5 (Bounded Atom)
Let $`P`$ be a program and $`||`$ be a level mapping for $`P`$. An atom $`A`$ is called *bounded wrt. $`||`$* if the set of all $`|A^{}|`$, where $`A^{}`$ is a ground instance of $`A`$, is finite. In this case we denote by $`max|A|`$ the maximum value in this set.
Notice that if an atom $`A`$ is bounded then, by definition of level mapping, also the corresponding negative literal, $`\neg A`$, is bounded.
Note also that, for atomic queries, this definition coincides with the definition of bounded query introduced in (?) in order to characterize terminating queries for acceptable programs. In fact, in case of atomic queries the notion of boundedness does not depend on a model.
### 2.3 Extension of a Program
In this paper we consider a hierarchical situation where a program uses another one as a subprogram. The following definition formalizes this situation.
###### Definition 6 (Extension)
Let $`P`$ and $`R`$ be two programs. A relation $`p`$ is *defined in* $`P`$ if $`p`$ occurs in a head of a clause of $`P`$; a literal $`L`$ is *defined in* $`P`$ if $`rel(L)`$ is defined in $`P`$; $`P`$ *extends* $`R`$, denoted $`PR`$, if no relation defined in $`P`$ occurs in $`R`$.
Informally, $`P`$ extends $`R`$ if $`P`$ defines new relations with respect to $`R`$. Note that $`P`$ and $`R`$ are independent if no relation defined in $`P`$ occurs in $`R`$ and no relation defined in $`R`$ occurs in $`P`$, i.e. $`PR`$ and $`RP`$.
In the sequel we will study termination in a hierarchy of programs.
###### Definition 7 (Hierarchy of Programs)
Let $`P_1,\mathrm{},P_n`$ be programs such that for all $`i\{1,\mathrm{},n1\}`$, $`P_{i+1}(P_1\mathrm{}P_i)`$. Then we call $`P_n\mathrm{}P_1`$ a *hierarchy of programs*.
## 3 Hierarchical Termination
This section contains our main results which show how termination proofs of separate programs can be combined to obtain proofs of larger programs. We start with a technical result, dealing with the case in which a program consists of a hierarchical combination of two modules. This is the base both of a generalization to a hierarchy of $`n`$ programs and of an iterative proof method for termination presented in Section 5. Let us first introduce the following notion of $`P`$-closed class of queries.
###### Definition 8 (P-closed Class)
Let $`C`$ be a class of queries and $`P`$ be a program. We say that $`C`$ is $`P`$-closed if it is closed under non-empty prefix (i.e., it contains all the non-empty prefixes of its elements) and for each query $`QC`$, every LDNF-descendant of $`P\{Q\}`$ is contained in $`C`$.
Note that if $`C`$ is $`P`$-closed, then for each query $`QC`$, $`\mathrm{𝐶𝑎𝑙𝑙}_P(Q)C`$.
We can now state our first general theorem. Notice that if $`P`$ extends $`R`$ and $`P`$ is acceptable wrt. some level mapping $`||`$ and model $`M`$, then $`P`$ is acceptable also wrt. the level mapping $`||^{}`$ and $`M`$, where $`||^{}`$ is defined on the Herbrand base of the union of the two programs $`B_{PR}`$ and it takes the value $`0`$ on the literals which are not defined in $`P`$ (and hence, in particular, on the literals which occur in $`P`$ but are defined in $`R`$). This shows that in each module it is sufficient to compare only the level of the literals defined inside it, while we can ignore literals defined outside the module. In the following we make use of this observation in order to associate to each module in a hierarchy a level mapping which is independent from the context.
###### Theorem 9
Let $`P`$ and $`R`$ be two programs such that $`P`$ extends $`R`$, $`M`$ be a complete model of $`PR`$ and $`C`$ be a $`(PR)`$-closed class of queries. Suppose that
* $`P`$ is acceptable wrt. a level mapping $`||`$ and $`M`$,
* for all queries $`QC`$, all LDNF-derivations of $`R\{Q\}`$ are finite,
* for all atoms $`AC`$, if $`A`$ is defined in $`P`$ then $`A`$ is bounded wrt. $`||`$.
Then for all queries $`QC`$, all LDNF-derivations of $`(PR)\{Q\}`$ are finite.
###### Proof 3.10.
By the fact that $`C`$ is $`(PR)`$-closed and Proposition 1, it is sufficient to prove that for all positive literals $`AC`$, all LDNF-derivations of $`(PR)\{A\}`$ are finite. Let us consider an atom $`AC`$.
If $`A`$ is defined in $`R`$, then the thesis trivially holds by hypothesis.
If $`A`$ is defined in $`P`$, $`A`$ is bounded wrt. $`||`$ by hypothesis and thus $`\mathrm{𝑚𝑎𝑥}|A|`$ is defined. The proof proceeds by induction on $`\mathrm{𝑚𝑎𝑥}|A|`$.
*Base*. Let $`\mathrm{𝑚𝑎𝑥}|A|=0`$. In this case, by acceptability of $`P`$, there are no clauses in $`P`$ whose head unifies with $`A`$ and whose body is non-empty. Hence, the thesis holds.
*Induction step*. Let $`\mathrm{𝑚𝑎𝑥}|A|>0`$. It is sufficient to prove that for all direct descendants $`(L_1,\mathrm{},L_n)`$ in the LDNF-tree of $`(PR)\{A\}`$, if $`\theta _i`$ is a computed answer for $`P\{L_1,\mathrm{},L_{i1}\}`$ then all LDNF-derivations of $`(PR)\{L_i\theta _i\}`$ are finite.
Let $`c:H^{}L_1^{},\mathrm{},L_n^{}`$ be a clause of $`P`$ such that $`\sigma =mgu(H^{},A)`$. Let $`H=H^{}\sigma `$ and for all $`i\{1,\mathrm{},n\}`$, let $`L_i=L_i^{}\sigma `$ and $`\theta _i`$ be a substitution such that $`\theta _i`$ is a computed answer of $`L_1,\mathrm{},L_{i1}`$ in $`PR`$.
We distinguish two cases. If $`L_i`$ is defined in $`R`$ then the thesis follows by hypothesis.
Suppose that $`L_i`$ is defined in $`P`$. We prove that $`L_i\theta _i`$ is bounded and $`\mathrm{𝑚𝑎𝑥}|A|>\mathrm{𝑚𝑎𝑥}|L_i\theta _i|`$. The thesis will follow by the induction hypothesis.
Let $`\gamma `$ be a substitution such that $`L_i\theta _i\gamma `$ is ground. By soundness of LDNF-resolution (?), there exists $`\gamma ^{}`$ such that $`M(L_1,\mathrm{},L_{i1})\gamma ^{}`$ and $`c\sigma \gamma ^{}`$ is a ground instance of $`c`$ and $`L_i\gamma ^{}=L_i\theta _i\gamma `$. Therefore
$$\begin{array}{ccccc}|L_i\theta _i\gamma |\hfill & =\hfill & |L_i\gamma ^{}|\hfill & & \\ & =\hfill & |L_i^{}\sigma \gamma ^{}|\hfill & (\text{since }L_i=L_i^{}\sigma )\hfill & \\ & <\hfill & |H^{}\sigma \gamma ^{}|\hfill & (\text{since }P\text{ is acceptable})\hfill & \\ & =\hfill & |A\sigma \gamma ^{}|\hfill & (\text{since }\sigma =mgu(H^{},A)).\hfill & \end{array}$$
Since $`A`$ is bounded, we can conclude that $`L_i\theta _i`$ is bounded and also that $`\mathrm{𝑚𝑎𝑥}|A|>\mathrm{𝑚𝑎𝑥}|L_i\theta _i|`$.
We are going to extend the above theorem in order to handle the presence of more than two modules. We need to introduce more notation. Let us consider the case of a program $`P`$ consisting of a hierarchy $`R_n\mathrm{}R_1`$ of distinct modules, and satisfying the property that each module, $`R_i`$, is acceptable wrt. a distinct level mapping, $`||_i`$, and a complete model, $`M`$, of the whole program. Under these assumptions we identify a specific class of queries which terminate in the whole program. We characterize the class of terminating queries in terms of the following notion of strong boundedness. This class enjoys the property of being $`P`$-closed.
###### Definition 3.11 (Strongly Bounded Query).
Let the program $`P:=R_1\mathrm{}R_n`$ be a hierarchy $`R_n\mathrm{}R_1`$ and $`||_1,\mathrm{},||_n`$ be level mappings for $`R_1,\mathrm{},R_n`$, respectively. A query $`Q`$ is called *strongly bounded wrt. $`P`$ and $`||_1,\mathrm{},||_n`$* if
* for all atoms $`A\mathrm{𝐶𝑎𝑙𝑙}_P(Q)`$, if $`A`$ is defined in $`R_i`$ (with $`i\{1,\mathrm{},n\}`$) then $`A`$ is bounded wrt. $`||_i`$.
Notice that the notion of boundedness for an atom (see Definition 5) does not depend on the choice of a particular model of $`P`$. As a consequence, also the definition of strong boundedness does not refer to any model of $`P`$; however, it refers to the LDNF-derivations of $`P`$. For this reason, a ground atom is always bounded but not necessarily strongly bounded. On the other hand, if $`A`$ is strongly bounded then it is bounded too.
The following remark follows immediately.
###### Remark 3.12.
Let the query $`Q`$ be strongly bounded wrt. $`P`$ and $`||_1,\mathrm{},||_n`$, where $`P`$ is a hierarchy $`R_n\mathrm{}R_1`$ . Let $`i\{1,\mathrm{},n\}`$. If $`Q`$ is defined in $`R_1\mathrm{}R_i`$ then $`Q`$ is strongly bounded wrt. $`R_1\mathrm{}R_i`$ and $`||_1,\mathrm{},||_i`$.
In order to verify whether a query Q is strongly bounded wrt. a given program $`P`$ one can perform a call-pattern analysis (???) which allows us to infer information about the form of the call-patterns, i.e., the atoms that will be possibly called during the execution of $`P\{Q\}`$. However this is not the only way for guaranteeing strong boundedness. There are classes of programs and queries for which strong boundedness can be proved in a straightforward way. This is shown in the following section.
Let us illustrate the notion of strong boundedness through an example.
###### Example 3.13.
Let LIST01 be the following program which defines the proper lists of 0’s and 1’s, i.e. lists containing only 0’s and 1’s and at least two distinct elements, as follows:
| r1: | list01(\[ \],0,0). |
| --- | --- |
| r2: | list01(\[0|Xs\],s(N0),N1) $``$list01(Xs,N0,N1). |
| r3: | list01(\[1|Xs\],N0,s(N1)) $``$list01(Xs,N0,N1). |
| r4: | length(\[ \],0). |
| r5: | length(\[X|Xs\],s(N)) $``$length(Xs,N). |
| r6: | plist01(Ls) $``$list01(Ls,N0,N1), |
| | | $`\neg `$length(Ls,N0), $`\neg `$length(Ls,N1). |
Let us distinguish two modules in LIST01: $`R_1=\{𝚛_\mathrm{𝟷},𝚛_\mathrm{𝟸},𝚛_\mathrm{𝟹},𝚛_\mathrm{𝟺},𝚛_\mathrm{𝟻}\}`$ and $`R_2=\{𝚛_\mathrm{𝟼}\}`$ ($`R_2`$ extends $`R_1`$). Let $`||_1`$ be the natural level mapping for $`R_1`$ defined by:
| $`|\mathrm{𝚕𝚒𝚜𝚝𝟶𝟷}(\mathrm{𝑙𝑠},\mathit{n0},\mathit{n1})|_\mathrm{𝟷}=|\mathrm{𝑙𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ |
| --- |
| $`|\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}(\mathrm{𝑙𝑠},n)|_\mathrm{𝟷}=|n|_{\mathrm{𝚜𝚒𝚣𝚎}}`$ |
where for a term $`t`$, if $`t`$ is a list then $`|t|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ is equal to the length of the list, otherwise it is $`0`$, while $`|t|_{\mathrm{𝚜𝚒𝚣𝚎}}`$ is the number of function symbols occurring in the term $`t`$. Let also $`||_2`$ be the trivial level mapping for $`R_2`$ defined by:
and assume that $`|L|_2=0`$, if $`L`$ is not defined in $`R_2`$.
Let us consider the following sets of atomic queries for $`\mathrm{𝙻𝙸𝚂𝚃𝟶𝟷}:=R_1R_2`$:
| $`Q_1=\{\mathrm{𝚕𝚒𝚜𝚝𝟶𝟷}(\mathrm{𝑙𝑠},\mathit{n0},\mathit{n1})|\mathrm{𝑙𝑠}\text{ is a list, possibly non-ground, of a fixed length}\};`$ |
| --- |
| $`Q_2=\{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}(\mathrm{𝑙𝑠},n)|n\text{ is a ground term of the form either 0 or }\text{s(s(…(0)))}\};`$ |
| $`Q_3=\{\mathrm{𝚙𝚕𝚒𝚜𝚝𝟶𝟷}(\mathrm{𝑙𝑠})|\mathrm{𝑙𝑠}\text{ is a list, possibly non-ground, of a fixed length}\}`$. |
By definition of $`||_1`$, all the atoms in $`Q_1`$ and $`Q_2`$ are bounded wrt. $`||_1`$. Analogously, all the atoms in $`Q_3`$ are bounded wrt. $`||_2`$. Notice that for all atoms $`ACall_P(Q_j)`$, with $`j\{1,2,3\}`$, there exists $`k\{1,2,3\}`$ such that $`AQ_k`$. Hence, if $`A`$ is defined in $`R_i`$ then $`A`$ is bounded wrt. $`||_i`$. This proves that the set of queries $`Q_1`$, $`Q_2`$ and $`Q_3`$ are strongly bounded wrt. $`\mathrm{𝙻𝙸𝚂𝚃𝟶𝟷}`$ and $`||_1`$, $`||_2`$.
Here we introduce our main result.
###### Theorem 3.14.
Let $`P:=R_1\mathrm{}R_n`$ be a program such that $`R_n\mathrm{}R_1`$ is a hierarchy, $`||_1,\mathrm{},||_n`$ be level mappings for $`R_1,\mathrm{},R_n`$, respectively, and $`M`$ be a complete model of $`P`$. Suppose that
* $`R_i`$ is acceptable wrt. $`||_i`$ and $`M`$, for all $`i\{1,\mathrm{},n\}`$.
* $`Q`$ is a query strongly bounded wrt. $`P`$ and $`||_1,\mathrm{},||_n`$.
Then all LDNF-derivations of $`P\{Q\}`$ are finite.
###### Proof 3.15.
Let $`Q`$ be a query strongly bounded wrt. $`P`$ and $`||_1,\mathrm{},||_n`$. We prove the theorem by induction on $`n`$.
Base. Let $`n=1`$. This case follows immediately by Theorem 9, where $`P=R_1`$, $`R`$ is empty and $`C`$ is the class of strongly bounded queries wrt. $`R_1`$ and $`||_1`$, and the fact that a strongly bounded atom is also bounded.
*Induction step*. Let $`n>1`$. Also this case follows by Theorem 9, where $`P=R_n`$, $`R=R_1\mathrm{}R_{n1}`$ and $`C`$ is the class of strongly bounded queries wrt. $`R_1\mathrm{}R_n`$ and $`||_1,\mathrm{},||_n`$. In fact,
* $`R_n`$ is acceptable wrt. $`||_n`$ and $`M`$;
* for all queries $`QC`$, all LDNF-derivations of $`(R_1\mathrm{}R_{n1})\{Q\}`$ are finite, by Remark 3.12 and the inductive hypothesis;
* for all atoms $`AC`$, if $`A`$ is defined in $`R_n`$ then $`A`$ is bounded wrt. $`||_n`$, by definition of strong boundedness.
Here are a few examples applying Theorem 3.14.
###### Example 3.16.
Let us reconsider the program of Example 3.13. In the program LIST01, $`R_1`$ and $`R_2`$ are acceptable wrt. any complete model and the level mappings $`||_1`$ and $`||_2`$, respectively. We already showed that $`Q_1,Q_2`$ and $`Q_3`$ are strongly bounded wrt. $`𝙻IST01`$ and $`||_1`$, $`||_2`$. Hence, by Theorem 3.14, all LDNF-derivations of $`𝙻IST01\{Q\}`$, where $`Q`$ is a query in $`Q_1,Q_2`$ or $`Q_3`$, are finite.
Notice that in the previous example the top module in the hierarchy, $`R_2`$, contains no recursion. Hence it is intuitively clear that any problem for termination cannot depend on it. This is reflected by the fact that the level mapping for $`R_2`$ is completely trivial. This shows how the hierarchical decomposition of the program can simplify the termination proof.
###### Example 3.17.
Consider the sorting program MERGESORT (?):
| c1: | mergesort(\[ \],\[ \]). |
| --- | --- |
| c2: | mergesort(\[X\],\[X\]). |
| c3: | mergesort(\[X,Y|Xs\],Ys) $``$ |
| | | split(\[X,Y|Xs\],X1s,X2s), |
| | | mergesort(X1s,Y1s), |
| | | mergesort(X2s,Y2s), |
| | | merge(Y1s,Y2s,Ys). |
| c4: | split(\[ \],\[ \],\[ \]). |
| c5: | split(\[X|Xs\],\[X|Ys\],Zs) $``$split(Xs,Zs,Ys). |
| c6: | merge(\[ \],Xs,Xs). |
| c7: | merge(Xs,\[ \],Xs). |
| c8: | merge(\[X|Xs\],\[Y|Ys\],\[X|Zs\]) $``$X\<=Y, merge(Xs,\[Y|Ys\],Zs). |
| c9: | merge(\[X|Xs\],\[Y|Ys\],\[Y|Zs\]) $``$X\>Y, merge(\[X|Xs\],Ys,Zs). |
Let us divide the program MERGESORT into three modules, $`R_1,R_2,R_3`$, such that $`R_3R_2R_1`$ as follows:
* $`R_3:=\{\mathrm{𝚌𝟷},\mathrm{𝚌𝟸},\mathrm{𝚌𝟹}\}`$, it defines the relation mergesort,
* $`R_2:=\{\mathrm{𝚌𝟺},\mathrm{𝚌𝟻}\}`$, it defines the relation split,
* $`R_1:=\{\mathrm{𝚌𝟼},\mathrm{𝚌𝟽},\mathrm{𝚌𝟾},\mathrm{𝚌𝟿}\}`$, it defines the relation merge.
Let us consider the natural level mappings
| $`|\mathrm{𝚖𝚎𝚛𝚐𝚎}(\mathrm{𝑥𝑠},\mathrm{𝑦𝑠},\mathrm{𝑧𝑠})|_\mathrm{𝟷}=|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}+|\mathrm{𝑦𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ |
| --- |
| $`|\mathrm{𝚜𝚙𝚕𝚒𝚝}(\mathrm{𝑥𝑠},\mathrm{𝑦𝑠},\mathrm{𝑧𝑠})|_\mathrm{𝟸}=|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ |
| $`|\mathrm{𝚖𝚎𝚛𝚐𝚎𝚜𝚘𝚛𝚝}(\mathrm{𝑥𝑠},\mathrm{𝑦𝑠})|_\mathrm{𝟹}=|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ |
and assume that for all $`i\{1,2,3\}`$, $`|\text{L}|_𝚒=\mathrm{𝟶}\text{ if L is not defined in }R_i`$.
All ground queries are strongly bounded wrt. the program MERGESORT and the level mappings $`||_1,||_2,||_3`$. Moreover, since the program is a definite one, $`R_1`$ and $`R_2`$ are acceptable wrt. any model and the level mappings $`||_1`$ and $`||_2`$, respectively, while $`R_3`$ is acceptable wrt. the level mapping $`||_3`$ and the model $`M`$ below:
| $`M=`$ | $`[\mathrm{𝚖𝚎𝚛𝚐𝚎𝚜𝚘𝚛𝚝}(\mathrm{𝚇𝚜},\mathrm{𝚈𝚜})][\mathrm{𝚖𝚎𝚛𝚐𝚎}(\mathrm{𝚇𝚜},\mathrm{𝚈𝚜},\mathrm{𝚉𝚜})]`$ |
| --- | --- |
| | $`\{\mathrm{𝚜𝚙𝚕𝚒𝚝}([],[],[])\}`$ |
| | $`\{\mathrm{𝚜𝚙𝚕𝚒𝚝}([x],[],[x])|x\text{ is any ground term}\}`$ |
| | $`\{\mathrm{𝚜𝚙𝚕𝚒𝚝}([x],[x],[])|x\text{ is any ground term}\}`$ |
| | $`\{\mathrm{𝚜𝚙𝚕𝚒𝚝}(\mathrm{𝑥𝑠},\mathrm{𝑦𝑠},\mathrm{𝑧𝑠})|\mathrm{𝑥𝑠},\mathrm{𝑦𝑠},\mathrm{𝑧𝑠}\text{ are ground terms and}`$ |
| | $`|\mathrm{𝑥𝑠}|_{length}2,|\mathrm{𝑥𝑠}|_{length}>|\mathrm{𝑦𝑠}|_{length},|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}>|\mathrm{𝑧𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}\}`$ |
where we denote by $`[A]`$ the set of all ground instances of an atom $`A`$.
Hence, by Theorem 3.14, all LDNF-derivations of $`𝙼ERGESORT\{Q\}`$, where $`Q`$ is a ground query, are finite.
Note that by exchanging the roles of $`R_1`$ and $`R_2`$ we would obtain the same result. In fact the definition of $`𝚖erge`$ and $`𝚜plit`$ are independent from each other.
## 4 Well-Behaving Programs
In this section we consider the problem of how to prove that a query is strongly bounded. In fact one could argue that checking strong boundedness is more difficult and less abstract than checking boundedness itself in the sense of (?): we have to refer to all LDNF-derivations instead of referring to a model, which might well look like a step backwards in the proof of termination of a program. This is only partly true: in order to check strong boundedness we can either employ tools based on abstract interpretation or concentrate our attention only on programs which exhibit useful persistence properties wrt. LDNF-resolution.
We now show how the well-established notions of well-moded and well-typed programs can be employed in order to verify strong boundedness and how they can lead to simple termination proofs.
### 4.1 Well-Moded Programs
The concept of a well-moded program is due to (?). The formulation we use here is from (?), and it is equivalent to that in (?). The original definition was given for definite programs (i.e., programs without negation), however it applies to general programs as well, just by considering literals instead of atoms. It relies on the concept of *mode*, which is a function that labels the positions of each predicate in order to indicate how the arguments of a predicate should be used.
###### Definition 4.18 (Mode).
Consider an $`n`$-ary predicate symbol $`p`$. By a *mode* for $`p`$ we mean a function $`m_p`$ from $`\{1,\mathrm{},n\}`$ to the set $`\{+,\}`$. If $`m_p(i)=+`$ then we call $`i`$ an input position of $`p`$; if $`m_p(i)=`$ then we call $`i`$ an output position of $`p`$. By a *moding* we mean a collection of modes, one for each predicate symbol.
In a moded program, we assume that each predicate symbol has a unique mode associated to it. Multiple moding may be obtained by simply renaming the predicates. We use the notation $`p(m_p(1),\mathrm{},m_p(n))`$ to denote the moding associated with a predicate $`p`$ (e.g., $`\mathrm{𝚊𝚙𝚙𝚎𝚗𝚍}(\text{+},\text{+},)`$). Without loss of generality, we assume, when writing a literal as $`p(𝐬,𝐭)`$, that we are indicating with $`𝐬`$ the sequence of terms filling in the input positions of $`p`$ and with $`𝐭`$ the sequence of terms filling in the output positions of $`p`$. Moreover, we adopt the convention that $`p(𝐬,𝐭)`$ could denote both negative and positive literals.
###### Definition 4.19 (Well-Moded).
* A query $`p_1(𝐬_1,𝐭_1),\mathrm{},p_n(𝐬_n,𝐭_n)`$ is called *well-moded* if for all $`i\{1,\mathrm{},n\}`$
$$Var(𝐬_i)\underset{j=1}{\overset{i1}{}}Var(𝐭_j).$$
* A clause $`p(𝐭_0,𝐬_{n+1})p_1(𝐬_1,𝐭_1),\mathrm{},p_n(𝐬_n,𝐭_n)`$ is called *well-moded* if for all $`i\{1,\mathrm{},n+1\}`$
$$Var(𝐬_i)\underset{j=0}{\overset{i1}{}}Var(𝐭_j).$$
* A program is called *well-moded* if all of its clauses are well-moded.
Note that well-modedness can be syntactically checked in a time which is linear wrt. the size of the program (query).
###### Remark 4.20.
If $`Q`$ is a well-moded query then all its prefixes are well-moded.
The following lemma states that well-moded queries are closed under LDNF-resolution. This result has been proved in (?) for LD-derivations and definite programs.
###### Lemma 4.21.
Let $`P`$ and $`Q`$ be a well-moded program and query, respectively. Then all LDNF-descendants of $`P\{Q\}`$ are well-moded.
###### Proof 4.22.
It is sufficient to extend the proof in (?) by showing that if a query $`\neg A,L_1,\mathrm{},L_n`$ is well-moded and $`A`$ is ground then both $`A`$ and $`L_1,\mathrm{},L_n`$ are well-moded. This follows immediately by definition of well-modedness. If $`A`$ is non-ground then the query above has no descendant.
When considering well-moded programs, it is natural to measure atoms only in their input positions (?).
###### Definition 4.23 (Moded Level Mapping).
Let $`P`$ be a moded program. A function $`||`$ is a *moded level mapping for $`P`$* if it is a level mapping for $`P`$ such that
* for any $`𝐬`$, $`𝐭`$ and $`𝐮`$, $`|p(𝐬,𝐭)|=|p(𝐬,𝐮)|`$.
Hence in a moded level mapping the level of an atom is independent from the terms in its output positions.
The following Remark and Proposition allow us to exploit well-modedness for applying Theorem 3.14.
###### Remark 4.24.
Let $`P`$ be a well-moded program. If $`Q`$ is well-moded, then $`\text{first}(Q)`$ is ground in its input position and hence it is bounded wrt. any moded level mapping for $`P`$. Moreover, by Lemma 4.21, every well-moded query is strongly bounded wrt. $`P`$ and any moded level mapping for $`P`$.
###### Proposition 4.25.
Let $`P:=R_1\mathrm{}R_n`$ be a *well-moded* program and $`R_n\mathrm{}R_1`$ a hierarchy, and $`||_1,\mathrm{},||_n`$ be *moded* level mappings for $`R_1,\mathrm{},R_n`$, respectively.
Then every well-moded query is strongly bounded wrt. $`P`$ and $`||_1,\mathrm{},||_n`$.
###### Example 4.26.
Let MOVE be the following program which defines a permutation between two lists such that only one element is moved. We introduce modes and we distinguish the two uses of append by renaming it as append1 and append2.
| mode $`\mathrm{𝚍𝚎𝚕𝚎𝚝𝚎}(+,,)`$. |
| --- |
| mode $`\mathrm{𝚊𝚙𝚙𝚎𝚗𝚍𝟷}(,,+`$). |
| mode $`\mathrm{𝚊𝚙𝚙𝚎𝚗𝚍𝟸}(+,+,)`$. |
| mode $`\mathrm{𝚖𝚘𝚟𝚎}(+,)`$. |
| r1: | delete(\[X|Xs\],X,Xs). |
| --- | --- |
| r2: | delete(\[X|Xs\],Y,\[X|Ys\]) $``$delete(Xs,Y,Ys). |
| r3: | append1(\[ \],Ys,Ys). |
| r4: | append1(\[X|Xs\],Ys,\[X|Zs\]) $``$append1(Xs,Ys,Zs). |
| r5: | append2(\[ \],Ys,Ys). |
| r6: | append2(\[X|Xs\],Ys,\[X|Zs\]) $``$append2(Xs,Ys,Zs). |
| r7: | move(Xs,Ys) $``$append1(X1s,X2s,Xs), |
| | | delete(X1s,X,Y1s), append2(Y1s,\[X|X2s\],Ys). |
Let us partition MOVE into the modules $`R_1=\{𝚛_\mathrm{𝟷},𝚛_\mathrm{𝟸},𝚛_\mathrm{𝟹},𝚛_\mathrm{𝟺},𝚛_\mathrm{𝟻},𝚛_\mathrm{𝟼}\}`$ and $`R_2=\{𝚛_\mathrm{𝟽}\}`$ ($`R_2`$ extends $`R_1`$). Let $`||_1`$ be the natural level mapping for $`R_1`$ defined by:
| $`|\mathrm{𝚊𝚙𝚙𝚎𝚗𝚍𝟷}(\mathrm{𝑥𝑠},\mathrm{𝑦𝑠},\mathrm{𝑧𝑠})|_\mathrm{𝟷}=|\mathrm{𝑧𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ |
| --- |
| $`|\mathrm{𝚊𝚙𝚙𝚎𝚗𝚍𝟸}(\mathrm{𝑥𝑠},\mathrm{𝑦𝑠},\mathrm{𝑧𝑠})|_\mathrm{𝟷}=|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$. |
| $`|\mathrm{𝚍𝚎𝚕𝚎𝚝𝚎}(\mathrm{𝑥𝑠},x,\mathrm{𝑦𝑠})|_\mathrm{𝟷}=|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$. |
$`R_2`$ does not contain any recursive definition hence let $`||_2`$ be the trivial level mapping defined by:
and assume that $`|L|_2=0`$, if $`L`$ is not defined in $`R_2`$.
The program $`𝙼OVE:=R_1R_2`$ is well-moded and hence by Proposition 4.25 every well-moded query is strongly bounded wrt. MOVE and $`||_1`$, $`||_2`$.
###### Example 4.27.
Let $`R_1`$ be the program which defines the relations member and is, $`R_2`$ be the program defining the relation count and $`R_3`$ be the program defining the relation diff with the moding and the definitions below.
| mode $`\mathrm{𝚖𝚎𝚖𝚋𝚎𝚛}(+,+)`$. |
| --- |
| mode $`\mathrm{𝚒𝚜}(,+)`$. |
| mode $`\mathrm{𝚍𝚒𝚏𝚏}(+,+,+,)`$. |
| mode $`\mathrm{𝚌𝚘𝚞𝚗𝚝}(+,+,)`$. |
| r1: | member(X,\[X|Xs\]). |
| r2: | member(X,\[Y|Xs\]) $``$member(X,Xs). |
| r3: | diff(Ls,I1,I2,N) $``$count(Ls,I1,N1), count(Ls,I2,N2), |
| | | N is N1-N2. |
| r4: | count(\[ \],I,0). |
| r5: | count(\[H|Ts\],I,M) $``$member(H,I), count(Ts,I,M1), |
| | | M is M1+1. |
| r6: | count(\[H|Ts\],I,M) $``$$`\neg `$ member(H,I), count(Ts,I,M). |
The relation $`\mathrm{𝚍𝚒𝚏𝚏}(\mathrm{𝑙𝑠},\mathit{i1},\mathit{i2},n)`$, given a list $`\mathrm{𝑙𝑠}`$ and two check-lists $`\mathit{i1}`$ and $`\mathit{i2}`$, defines the difference $`n`$ between the number of elements of $`\mathrm{𝑙𝑠}`$ occurring in $`\mathit{i1}`$ and the number of elements of $`\mathrm{𝑙𝑠}`$ occurring in $`\mathit{i2}`$. Clearly $`R_3R_2R_1`$. It is easy to see that $`R_1`$ is acceptable wrt. any complete model and the moded level mapping
$`R_2`$ is acceptable wrt. any complete model and the moded level mapping:
and $`R_3`$ is acceptable wrt. any complete model and the trivial moded level mapping:
where $`|L|_𝚒=\mathrm{𝟶}`$, if $`L`$ is not defined in $`R_i`$.
The program $`𝙳IFF:=R_1R_2R_3`$ is well-moded. Hence, by Proposition 4.25, every well-moded query is strongly bounded wrt. $`𝙳IFF`$ and $`||_1`$, $`||_2`$, $`||_3`$.
Note that the class of strongly bounded queries is generally larger than the class of well-moded queries. Consider for instance the program MOVE and the query $`Q:=`$ $`\mathrm{𝚖𝚘𝚟𝚎}([\mathrm{𝚇𝟷},\mathrm{𝚇𝟸}],\mathrm{𝚈𝚜}),\mathrm{𝚍𝚎𝚕𝚎𝚝𝚎}(\mathrm{𝚈𝚜},𝚈,\mathrm{𝚉𝚜})`$ which is not well-moded since it is not ground in the input position of the first atom. However $`Q`$ can be easily recognized to be strongly bounded wrt. MOVE and $`||_1`$, $`||_2`$ defined in Example 4.26. We will come back to this query later.
### 4.2 Well-Typed Programs
A more refined well-behavior property of programs, namely well-typedness, can also be useful in order to ensure the strong boundedness property.
The notion of well-typedness relies both on the concepts of *mode* and *type*. The following very general definition of a type is sufficient for our purposes.
###### Definition 4.28 (Type).
A *type* is a set of terms closed under substitution.
Assume as given a specific set of types, denoted by *Types*, which includes $`Any`$, the set of all terms, and $`Ground`$ the set of all ground terms.
###### Definition 4.29 (Type Associated with a Position).
A *type for an $`n`$-ary predicate symbol $`p`$* is a function $`t_p`$ from $`\{1,\mathrm{},n\}`$ to the set *Types*. If $`t_p(i)=T`$, we call $`T`$ *the type associated with the position $`i`$ of $`p`$*. Assuming a type $`t_p`$ for the predicate $`p`$, we say that a literal $`p(s_1,\mathrm{},s_n)`$ is *correctly typed in position $`i`$* if $`s_it_p(i)`$.
In a typed program we assume that every predicate $`p`$ has a fixed mode $`m_p`$ and a fixed type $`t_p`$ associated with it and we denote it by
$$p(m_p(1):t_p(1),\mathrm{},m_p(n):t_p(n)).$$
So, for instance, we write
$$\mathrm{𝚊𝚙𝚙𝚎𝚗𝚍}(+:\mathrm{𝐿𝑖𝑠𝑡},+:\mathrm{𝐿𝑖𝑠𝑡},:\mathrm{𝐿𝑖𝑠𝑡})$$
to denote the moded atom $`\mathrm{𝚊𝚙𝚙𝚎𝚗𝚍}(+,+,)`$ where the type associated with each argument position is $`\mathrm{𝐿𝑖𝑠𝑡}`$, i.e., the set of all lists.
We can then talk about types of input and of output positions of an atom.
The notion of well-typed queries and programs relies on the following concept of type judgement.
###### Definition 4.30 (Type Judgement).
By a *type judgement* we mean a statement of the form $`𝐬:𝐒𝐭:𝐓.`$ We say that a type judgement $`𝐬:𝐒𝐭:𝐓`$ *is true*, and write $`𝐬:𝐒𝐭:𝐓,`$ if for all substitutions $`\theta `$, $`𝐬\theta 𝐒`$ implies $`𝐭\theta 𝐓`$.
For example, the type judgements $`(x:\mathrm{𝑁𝑎𝑡},l:\mathrm{𝐿𝑖𝑠𝑡𝑁𝑎𝑡})([x|l]:\mathrm{𝐿𝑖𝑠𝑡𝑁𝑎𝑡})`$ and $`([x|l]:\mathrm{𝐿𝑖𝑠𝑡𝑁𝑎𝑡})(l:\mathrm{𝐿𝑖𝑠𝑡𝑁𝑎𝑡})`$ are both true.
A notion of well-typed program has been first introduced in (?) and also studied in (?) and in (?). Similarly to well-moding, the notion was developed for definite programs. Here we extend it to general programs.
In the following definition, we assume that $`𝐢_s:𝐈_s`$ is the sequence of typed terms filling in the input positions of $`L_s`$ and $`𝐨_s:𝐎_s`$ is the sequence of typed terms filling in the output positions of $`L_s`$.
###### Definition 4.31 (Well-Typed).
* A query $`L_1,\mathrm{},L_n`$ is called *well-typed* if for all $`j\{1,\mathrm{},n\}`$
$$𝐨_{j_1}:𝐎_{j_1},\mathrm{},𝐨_{j_k}:𝐎_{j_k}𝐢_j:𝐈_j$$
where $`L_{j_1},\mathrm{},L_{j_k}`$ are all the positive literals in $`L_1,\mathrm{},L_{j1}`$.
* A clause $`L_0L_1,\mathrm{},L_n`$ is called *well-typed* if for all $`j\{1,\mathrm{},n\}`$
$$𝐢_0:𝐈_0,𝐨_{j_1}:𝐎_{j_1},\mathrm{},𝐨_{j_k}:𝐎_{j_k}𝐢_j:𝐈_j$$
where $`L_{j_1},\mathrm{},L_{j_k}`$ are all the positive literals in $`L_1,\mathrm{},L_{j1}`$, and
$$𝐢_0:𝐈_0,𝐨_{j_1}:𝐎_{j_1},\mathrm{},𝐨_{j_h}:𝐎_{j_h}𝐨_0:𝐎_0$$
where $`L_{j_1},\mathrm{},L_{j_h}`$ are all the positive literals in $`L_1,\mathrm{},L_n`$.
* A program is called *well-typed* if all of its clauses are well-typed.
Note that an atomic query is well-typed iff it is correctly typed in its input positions and a unit clause $`p(𝐬:𝐒,𝐭:𝐓)`$ is well-typed if $`𝐬:𝐒𝐭:𝐓`$.
The difference between Definition 4.31 and the one usually given for definite programs is that the correctness of the terms filling in the output positions of negative literals cannot be used to deduce the correctness of the terms filling in the input positions of a literal to the right (or the output positions of the head in a clause). The two definitions coincide either for definite programs or for general programs whose negative literals have only input positions.
As an example, let us consider the trivial program
| p($`:\mathrm{𝐿𝑖𝑠𝑡}`$). |
| --- |
| q($`+:\mathrm{𝐿𝑖𝑠𝑡}`$). |
| p(\[\]). |
| q(\[\]). |
By adopting a straightforward extension of well-typedness to normal programs which considers also the outputs of negative literals, we would have that the query $`\neg 𝚙(𝚊),𝚚(𝚊)`$ is well-typed even if $`𝚊`$ is not a list. Moreover well-typedness would not be persistent wrt. LDNF-resolution since $`𝚚(𝚊)`$, which is the first LDNF-resolvent of the previous query, is no more well-typed. Our extended definition and the classical one coincide either for definite programs or for general programs whose negative literals have only input positions.
For definite programs, well-modedness can be viewed as a special case of well-typedness if we consider only one type: $`Ground`$. With our extended definitions of well-moded and well-typed general programs this is no more true. We could have given a more complicated definition for well-typedness in order to capture also well-modedness as a special case. For the sake of simplicity, we prefer to give two distinct and simpler definitions.
###### Remark 4.32.
If $`Q`$ is a well-typed query, then all its non-empty prefixes are well-typed. In particular, $`\mathrm{𝑓𝑖𝑟𝑠𝑡}(Q)`$ is well-typed.
The following Lemma shows that well-typed queries are closed under LDNF-resolution. It has been proved in (?) for definite programs.
###### Lemma 4.33.
Let $`P`$ and $`Q`$ be a well-typed program and query, respectively. Then all LDNF-descendants of $`P\{Q\}`$ are well-typed.
###### Proof 4.34.
Similarly to the case of well-moded programs, to extend the result to general programs it is sufficient to show that if a query $`Q:=\neg A,L_1,\mathrm{},L_n`$ is well-typed then both $`A`$ and $`L_1,\mathrm{},L_n`$ are well-typed. In fact, by Remark 4.32, $`\neg A=\mathrm{𝑓𝑖𝑟𝑠𝑡}(Q)`$ is well-typed and by Definition 4.31, if the first literal in a well-typed query is negative, then it is not used to deduce well-typedness of the rest of the query.
It is now natural to exploit well-typedness in order to check strong boundedness. Analogously to well-moded programs, there are level mappings that are more natural in presence of type information. They are the level mappings for which every well-typed atom is bounded. By Lemma 4.33 we have that a well-typed query $`Q`$ is strongly bounded wrt. a well-typed program $`P`$ and any such level mapping. This is stated by the next proposition.
###### Proposition 4.35.
Let $`P:=R_1\mathrm{}R_n`$ be a *well-typed* program and $`R_n\mathrm{}R_1`$ be a hierarchy, and $`||_1,\mathrm{},||_n`$ be level mappings for $`R_1,\mathrm{},R_n`$, respectively. Suppose that for every well-typed atom $`A`$, if $`A`$ is defined in $`R_i`$ then $`A`$ is bounded wrt. $`||_i`$, for $`i\{1,\mathrm{},n\}`$. Then every well-typed query is strongly bounded wrt. $`P`$ and $`||_1,\mathrm{},||_n`$.
###### Example 4.36.
Let us consider again the modular proof of termination for $`\mathrm{𝙼𝙾𝚅𝙴}:=R_1R_2`$, where $`R_1`$ defines the relations append1, append2 and delete, while $`R_2`$, which extends $`R_1`$, defines the relation move. We consider the moding of Example 4.26 with the following types:
| $`\mathrm{𝚍𝚎𝚕𝚎𝚝𝚎}(+:\mathrm{𝐿𝑖𝑠𝑡},:\mathrm{𝐴𝑛𝑦},:\mathrm{𝐿𝑖𝑠𝑡})`$ |
| --- |
| $`\mathrm{𝚊𝚙𝚙𝚎𝚗𝚍𝟷}(:\mathrm{𝐿𝑖𝑠𝑡},:\mathrm{𝐿𝑖𝑠𝑡},+:\mathrm{𝐿𝑖𝑠𝑡}`$) |
| $`\mathrm{𝚊𝚙𝚙𝚎𝚗𝚍𝟸}(+:\mathrm{𝐿𝑖𝑠𝑡},+:\mathrm{𝐿𝑖𝑠𝑡},:\mathrm{𝐿𝑖𝑠𝑡})`$ |
| $`\mathrm{𝚖𝚘𝚟𝚎}(+:\mathrm{𝐿𝑖𝑠𝑡},:\mathrm{𝐿𝑖𝑠𝑡})`$. |
Program $`\mathrm{𝙼𝙾𝚅𝙴}`$ is *well-typed* in the assumed modes and types.
Let us consider the same level mappings as used in Example 4.26. We have already seen that $`R_2`$ is acceptable wrt. $`||_2`$ and any model, and $`R_1`$ is acceptable wrt. $`||_1`$ and any model. By definition of $`||_2`$ and $`||_1`$, one can easily see that
* every well-typed atom $`A`$ defined in $`R_i`$ is bounded wrt. $`||_i`$.
Hence, by Proposition 4.35,
* every well-typed query is strongly bounded wrt. $`\mathrm{𝙼𝙾𝚅𝙴}`$ and $`||_1`$, $`||_2`$.
Let us consider again the query $`Q:=\mathrm{𝚖𝚘𝚟𝚎}([\mathrm{𝚇𝟷},\mathrm{𝚇𝟸}],\mathrm{𝚈𝚜}),\mathrm{𝚍𝚎𝚕𝚎𝚝𝚎}(\mathrm{𝚈𝚜},𝚈,\mathrm{𝚉𝚜})`$ which is not well-moded but it is well-typed. We have that $`Q`$ is strongly bounded wrt. MOVE and $`||_1`$, $`||_2`$, and consequently, by Theorem 3.14, that all LDNF-derivations of $`𝙼OVE\{Q\}`$ are finite.
###### Example 4.37.
Consider the program COLOR\_MAP from (?) which generates a coloring of a map in such a way that no two neighbors have the same color. The map is represented as a list of regions and colors as a list of available colors. In turn, each region is determined by its name, color and the colors of its neighbors, so it is represented as a term region(name,color,neighbors), where neighbors is a list of colors of the neighboring regions.
| c1: | color\_map(\[ \],Colors). |
| --- | --- |
| c2: | color\_map(\[Region|Regions\],Colors) $``$ |
| | | color\_region(Region,Colors), |
| | | color\_map(Regions,Colors). |
| c3: | color\_region(region(Name,Color,Neighbors),Colors) $``$ |
| | | select(Color,Colors,Colors1) |
| | | subset(Neighbors,Colors1). |
| c4: | select(X,\[X|Xs\],Xs). |
| c5: | select(X,\[Y|Xs\],\[Y|Zs\]) $``$select(X,Xs,Zs). |
| c6: | subset(\[ \],Ys). |
| c7: | subset(\[X|Xs\],Ys) $``$member(X,Ys), subset(Xs,Ys). |
| c8: | member(X,\[X|Xs\]). |
| c9: | member(X,\[Y|Xs\]) $``$member(X,Xs). |
Consider the following modes and types for the program COLOR\_MAP:
| $`\mathrm{𝚌𝚘𝚕𝚘𝚛}\mathrm{\_}\mathrm{𝚖𝚊𝚙}(+:\mathrm{𝐿𝑖𝑠𝑡𝑅𝑒𝑔𝑖𝑜𝑛},+:\mathrm{𝐿𝑖𝑠𝑡})`$ |
| --- |
| $`\mathrm{𝚌𝚘𝚕𝚘𝚛}\mathrm{\_}\mathrm{𝚛𝚎𝚐𝚒𝚘𝚗}(+:\mathrm{𝑅𝑒𝑔𝑖𝑜𝑛},+:\mathrm{𝐿𝑖𝑠𝑡})`$ |
| $`\mathrm{𝚜𝚎𝚕𝚎𝚌𝚝}(+:\mathrm{𝐴𝑛𝑦},+:\mathrm{𝐿𝑖𝑠𝑡},:\mathrm{𝐿𝑖𝑠𝑡})`$ |
| $`\mathrm{𝚜𝚞𝚋𝚜𝚎𝚝}(+:\mathrm{𝐿𝑖𝑠𝑡},+:\mathrm{𝐿𝑖𝑠𝑡})`$ |
| $`\mathrm{𝚖𝚎𝚖𝚋𝚎𝚛}(+:\mathrm{𝐴𝑛𝑦},+:\mathrm{𝐿𝑖𝑠𝑡})`$ |
where
* *Region* is the set of all terms of the form region(name,color,neighbors) with $`\mathrm{𝚗𝚊𝚖𝚎},\mathrm{𝚌𝚘𝚕𝚘𝚛}\mathrm{𝐴𝑛𝑦}`$ and $`\mathrm{𝚗𝚎𝚒𝚐𝚑𝚋𝚘𝚛𝚜}\mathrm{𝐿𝑖𝑠𝑡}`$,
* *ListRegion* is the set of all lists of regions.
We can check that COLOR\_MAP is well-typed in the assumed modes and types.
We can divide the program COLOR\_MAP into four distinct modules, $`R_1,R_2,R_3,R_4`$, in the hierarchy $`R_4R_3R_2R_1`$ as follows:
* $`R_4:=\{\mathrm{𝚌𝟷},\mathrm{𝚌𝟸}\}`$ defines the relation color\_map,
* $`R_3:=\{\mathrm{𝚌𝟹}\}`$ defines the relation color\_region,
* $`R_2:=\{\mathrm{𝚌𝟺},\mathrm{𝚌𝟻},\mathrm{𝚌𝟼},\mathrm{𝚌𝟽}\}`$ defines the relations select and subset,
* $`R_1:=\{\mathrm{𝚌𝟾},\mathrm{𝚌𝟿}\}`$ defines the relation member.
Each $`R_i`$ is trivially acceptable wrt. any model $`M`$ and the simple level mapping $`||_i`$ defined below:
| $`|\mathrm{𝚌𝚘𝚕𝚘𝚛}\mathrm{\_}\mathrm{𝚖𝚊𝚙}(\mathrm{𝑥𝑠},y𝚜)|_\mathrm{𝟺}=|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ |
| --- |
| $`|\mathrm{𝚌𝚘𝚕𝚘𝚛}\mathrm{\_}\mathrm{𝚛𝚎𝚐𝚒𝚘𝚗}(x,\mathrm{𝑥𝑠})|_\mathrm{𝟹}=\mathrm{𝟷}`$ |
| $`|\mathrm{𝚜𝚎𝚕𝚎𝚌𝚝}(x,\mathrm{𝑥𝑠},\mathrm{𝑦𝑠})|_\mathrm{𝟸}=|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ |
| $`|\mathrm{𝚜𝚞𝚋𝚜𝚎𝚝}(\mathrm{𝑥𝑠},\mathrm{𝑦𝑠})|_\mathrm{𝟸}=|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ |
| $`|\mathrm{𝚖𝚎𝚖𝚋𝚎𝚛}(x,\mathrm{𝑥𝑠})|_\mathrm{𝟷}=|\mathrm{𝑥𝑠}|_{\mathrm{𝚕𝚎𝚗𝚐𝚝𝚑}}`$ |
where for all $`i\{1,2,3,4\}`$, $`|L|_i=0`$, if $`L`$ is not defined in $`R_i`$.
Moreover, for every well-typed atom $`A`$ and $`i\{1,2,3,4\}`$, if $`A`$ is defined in $`R_i`$ then $`A`$ is bounded wrt. $`||_i`$. Hence, by Proposition 4.35,
* every well-typed query is strongly bounded wrt. the program $`\mathrm{𝙲𝙾𝙻𝙾𝚁}\mathrm{\_}\mathrm{𝙼𝙰𝙿}`$ and $`||_1,\mathrm{},||_4`$.
This proves that all LDNF-derivations of the program COLOR\_MAP starting in a well-typed query are finite. In particular, all the LDNF-derivations starting in a query of the form $`\mathrm{𝚌𝚘𝚕𝚘𝚛}\mathrm{\_}\mathrm{𝚖𝚊𝚙}(\mathrm{𝑥𝑠},\mathrm{𝑦𝑠})`$, where $`\mathrm{𝑥𝑠}`$ is a list of regions and $`\mathrm{𝑦𝑠}`$ is a list, are finite. Note that in proving termination of such queries the choice of a model is irrelevant. Moreover, since such queries are well-typed, their input arguments are required to have a specified structure, but they are not required to be ground terms as in the case of well-moded queries. Hence, well-typedness allows us to reason about a larger class of queries with respect to well-modedness.
This example is also discussed in (?). In order to prove its termination they define a particular level mapping $`||`$, obtained by combining the level mappings of each module, and a special model $`M`$ wrt. which the whole program COLOR\_MAP is acceptable. Both the level mapping $`||`$ and the model $`M`$ are non-trivial.
## 5 Iterative Proof Method
In the previous section we have seen how we can exploit properties which are preserved by LDNF-resolution, such as well-modedness and well-typedness, for developing a modular proof of termination in a hierarchy of programs. In this section we show how these properties allow us to apply our general result, i.e., Theorem 9, also in an iterative way.
###### Corollary 5.38.
Let $`P`$ and $`R`$ be two programs such that $`PR`$ is well-moded and $`P`$ extends $`R`$, and $`M`$ be a complete model of $`PR`$. Suppose that
* $`P`$ is acceptable wrt. a moded level mapping $`||`$ and $`M`$,
* for all well-moded queries $`Q`$, all LDNF-derivations $`R\{Q\}`$ are finite.
Then for all well-moded queries $`Q`$, all LDNF-derivations of $`(PR)\{Q\}`$ are finite.
###### Proof 5.39.
Let $`C`$ be the class of well-moded queries of $`PR`$. By Remark 4.20 and Lemma 4.21, $`C`$ is $`(PR)`$-closed. Moreover
* $`P`$ is acceptable wrt. a moded level mapping $`||`$ and $`M`$, by hypothesis;
* for all well-moded queries $`Q`$, all LDNF-derivations of $`R\{Q\}`$ are finite, by hypothesis;
* for all well-moded atoms $`A`$, if $`A`$ is defined in $`P`$ then $`A`$ is bounded wrt. $`||`$, by Remark 4.24, since $`||`$ is a moded level mapping.
Hence by Theorem 9 we get the thesis.
Note that this result allows one to incrementally prove well-termination for general programs thus extending the result given in (?) for definite programs.
A similar result can be stated also for well-typed programs and queries, provided that there exists a level mapping for $`P`$ implying boundedness of atomic well-typed queries.
###### Corollary 5.40.
Let $`P`$ and $`R`$ be two programs such that $`PR`$ is well-typed and $`P`$ extends $`R`$, and $`M`$ be a complete model of $`PR`$. Suppose that
* $`P`$ is acceptable wrt. a level mapping $`||`$ and $`M`$,
* every well-typed atom defined in $`P`$ is bounded wrt. $`||`$,
* for all well-typed queries $`Q`$, all LDNF-derivations of $`R\{Q\}`$ are finite.
Then for all well-typed queries $`Q`$, all LDNF-derivations of $`(PR)\{Q\}`$ are finite.
###### Proof 5.41.
Let $`C`$ be the class of well-typed queries of $`PR`$. By Remark 4.32 and Lemma 4.33, $`C`$ is $`(PR)`$-closed. Moreover
* $`P`$ is acceptable wrt. a level mapping $`||`$ and $`M`$, by hypothesis;
* for all well-typed queries $`Q`$, all LDNF-derivations of $`R\{Q\}`$ are finite, by hypothesis;
* for all well-typed atoms $`A`$, if $`A`$ is defined in $`P`$ then $`A`$ is bounded wrt. $`||`$, by hypothesis.
Hence by Theorem 9 we have the thesis.
###### Example 5.42.
Let us consider again the program COLOR\_MAP with the same modes and types as in Example 4.37. We apply the iterative termination proof given by Corollary 5.40 to COLOR\_MAP.
First step. We can consider at first two trivial modules, $`R_1:=\{\mathrm{𝚌𝟾},\mathrm{𝚌𝟿}\}`$ which defines the relation member, and $`R_0:=\mathrm{}`$. We already know that
* $`R_1`$ is acceptable wrt. any model $`M`$ and the level mapping $`||_1`$ already defined;
* all well-typed atoms $`A`$, defined in $`R_1`$, are bounded wrt. $`||_1`$;
* for all well-typed queries $`Q`$, all LDNF-derivations of $`R_0\{Q\}`$ are trivially finite.
Hence, by Corollary 5.40, for all well-typed queries $`Q`$, all LDNF-derivations of $`(R_1R_0)\{Q\}`$ are finite.
Second step. We can now iterate the process one level up. Let us consider the two modules, $`R_2:=\{\mathrm{𝚌𝟺},\mathrm{𝚌𝟻},\mathrm{𝚌𝟼},\mathrm{𝚌𝟽}\}`$ which defines the relations select and subset, and $`R_1:=\{\mathrm{𝚌𝟾},\mathrm{𝚌𝟿}\}`$ which defines the relation member and it is equal to $`(R_1R_0)`$ of the previous step. We already showed in Example 4.37 that
* $`R_2`$ is acceptable wrt. any model $`M`$ and the level mapping $`||_2`$ already defined;
* all well-typed atoms $`A`$, defined in $`R_2`$, are bounded wrt. $`||_2`$;
* for all well-typed queries $`Q`$, all LDNF-derivations of $`R_1\{Q\}`$ are finite.
Hence, by Corollary 5.40, for all well-typed queries $`Q`$, all LDNF-derivations of $`(R_2R_1)\{Q\}`$ are finite.
By iterating the same reasoning for two steps more, we can prove that all LDNF-derivations of the program COLOR\_MAP starting in a well-typed query are finite.
Our iterative method applies to a hierarchy of programs where on the lowest module, $`R`$, we require termination wrt. a particular class of queries. This can be a weaker requirement on $`R`$ than acceptability as shown in the following contrived example.
###### Example 5.43.
Let $`R`$ define the predicate lcount which counts the number of natural numbers in a list.
| $`\mathrm{𝚕𝚌𝚘𝚞𝚗𝚝}(+:\mathrm{𝐿𝑖𝑠𝑡},:\mathrm{𝑁𝑎𝑡})`$ |
| --- |
| $`\mathrm{𝚗𝚊𝚝}(+:\mathrm{𝐴𝑛𝑦})`$. |
| r1: | lcount(\[ \],0). |
| r2: | lcount(\[X|Xs\],s(N)) $``$nat(X), lcount(Xs,N). |
| r3: | lcount(\[X|Xs\],N) $``$$`\neg `$ nat(X), lcount(Xs,N). |
| r4: | lcount(0,N) $``$lcount(0,s(N)). |
| r5: | nat(0). |
| r6: | nat(s(N)) $``$nat(N). |
$`R`$ is well-typed wrt. the specified modes and types. Note that $`R`$ cannot be acceptable due to the presence of clause r4. On the other hand, the program terminates for all well-typed queries.
Consider now the following program $`P`$ which extends $`R`$. The predicate split, given a list of lists, separates the list elements containing more than max natural numbers from the other lists:
| $`\mathrm{𝚜𝚙𝚕𝚒𝚝}(+:\mathrm{𝐿𝑖𝑠𝑡𝐿𝑖𝑠𝑡},:\mathrm{𝐿𝑖𝑠𝑡𝐿𝑖𝑠𝑡},:\mathrm{𝐿𝑖𝑠𝑡𝐿𝑖𝑠𝑡})`$ |
| --- |
| \>$`(+:\mathrm{𝑁𝑎𝑡},+:\mathrm{𝑁𝑎𝑡})`$ |
| \<=$`(+:\mathrm{𝑁𝑎𝑡},+:\mathrm{𝑁𝑎𝑡})`$ |
| p1: | split(\[ \],\[ \],\[ \]). |
| p2: | split(\[L|Ls\],\[L|L1\],L2) $``$lcount(L,N), N \> max, |
| | | split(Ls,L1,L2). |
| p3: | split(\[L|Ls\],L1,\[L|L2\]) $``$lcount(L,N), N \<= max, |
| | | split(Ls,L1,L2). |
where ListList denotes the set of all lists of lists, and max is a natural number. The program $`PR`$ is well-typed. Let us consider the simple level mapping $`||`$ for $`P`$ defined by:
which assigns level $`0`$ to any literal not defined in $`P`$. Note that
* $`P`$ is acceptable wrt. the level mapping $`||`$ and any complete model $`M`$,
* all well-typed atoms defined in $`P`$ are bounded wrt. $`||`$,
* for all well-typed queries $`Q`$, all LDNF-derivations of $`R\{Q\}`$ are finite.
Hence, by Corollary 5.40, for all well-typed queries $`Q`$, all LDNF-derivations of $`(PR)\{Q\}`$ are finite.
This example shows that well-typedness could be useful to exclude what might be called “dead code”.
## 6 Comparing with Apt and Pedreschi’s Approach
Our work can be seen as an extension of a proposal in (?). Hence we devote this section to a comparison with their approach.
On one hand, since our approach applies to general programs, it clearly covers cases which cannot be treated with the method proposed in (?), which was developed for definite programs. On the other hand, for definite programs the classes of queries and programs which can be treated by Apt and Pedreschi’s approach are properly included in those which can be treated by our method as we show in this section.
We first recall the notions of *semi-acceptability* and *bounded query* used in (?).
###### Definition 6.44 (Semi-acceptable Program).
Let $`P`$ be a definite program, $`||`$ be a level mapping for $`P`$ and $`M`$ be a model of $`P`$. $`P`$ is called *semi-acceptable wrt. $`||`$ and $`M`$* if for every clause $`A𝐀,B,𝐁`$ in $`ground(P)`$ such that $`M𝐀`$
* $`|A|>|B|,\text{ if }rel(A)rel(B)`$,
* $`|A||B|,\text{ if }rel(A)rel(B)`$.
###### Definition 6.45 (Bounded Query).
Let $`P`$ be a definite program, $`||`$ be a level mapping for $`P`$, and $`M`$ be a model of $`P`$.
* With each query $`Q:=L_1,\mathrm{},L_n`$ we associate $`n`$ sets of natural numbers defined as follows: For $`i\{1,\mathrm{},n\}`$,
$$|Q|_i^M=\{|L_i^{}||L_1^{},\mathrm{},L_n^{}\text{ is a ground instance of }Q\text{ and }ML_1^{},\mathrm{},L_{i1}^{}\}.$$
* A query $`Q`$ is called *bounded wrt. $`||`$ and $`M`$* if $`|Q|_i^M`$ is finite (i.e., if $`|Q|_i^M`$ has a maximum in $`𝐍`$) for all $`i\{1,\mathrm{},n\}`$.
###### Lemma 6.46.
Let $`P`$ be a definite program which is semi-acceptable wrt. $`||`$ and $`M`$. If $`Q`$ is a query bounded wrt. $`||`$ and $`M`$ then all LD-descendants of $`P\{Q\}`$ are bounded wrt. $`||`$ and $`M`$.
###### Proof 6.47.
It is a consequence of Lemma 3.6 in (?) and (the proof of) Lemma 5.4 in (?).
We can always decompose a definite program $`P`$ into a hierarchy of $`n1`$ programs $`P:=R_1\mathrm{}R_n`$, where $`R_n\mathrm{}R_1`$ in such a way that for every $`i\{1,\mathrm{},n\}`$ if the predicate symbols $`p_i`$ and $`q_i`$ are both defined in $`R_i`$ then neither $`p_iq_i`$ nor $`q_ip_i`$ (either they are mutually recursive or independent). We call such a hierarchy a *finest decomposition* of $`P`$.
The following property has two main applications. First it allows us to compare our approach with (?), then it provides an extension of Theorem 3.14 to hierarchies of semi-acceptable programs.
###### Proposition 6.48.
Let $`P`$ be a semi-acceptable program wrt. a level mapping $`||`$ and a model $`M`$ and $`Q`$ be a query strongly bounded wrt. $`P`$ and $`||`$. Let $`P:=R_1\mathrm{}R_n`$ be a finest decomposition of $`P`$ into a hierarchy of modules. Let $`||_i`$, with $`i\{1,\mathrm{},n\}`$, be defined in the following way: if $`A`$ is defined in $`R_i`$ then $`|A|_i=|A|`$ else $`|A|_i=0`$. Then
* every $`R_i`$ is acceptable wrt. $`||_i`$ and $`M`$ (with $`i\{1,\mathrm{},n\}`$),
* $`Q`$ is strongly bounded wrt. $`R_1\mathrm{}R_n`$ and $`||_1,\mathrm{},||_n`$.
###### Proof 6.49.
Immediate by the definitions of semi-acceptability and strongly boundedness, since we are considering a finest decomposition.
In order to compare our approach to the one presented in (?) we consider only Theorem 5.8 in (?), since this is their most general result which implies the other ones, namely Theorem 5.6 and Theorem 5.7.
###### Theorem 6.50 (Theorem 5.8 in (?)).
Let $`P`$ and $`R`$ be two definite programs such that $`P`$ extends $`R`$, and let $`M`$ be a model of $`PR`$. Suppose that
* $`R`$ is semi-acceptable wrt. $`||_R`$ and $`MB_R`$,
* $`P`$ is semi-acceptable wrt. $`||_P`$ and $`M`$,
* there exists a level mapping $`||||_P`$ such that for every ground instance of a clause from $`P`$, $`A𝐀,B,𝐁`$, such that $`M𝐀`$
+ $`A_PB_P`$, if $`\mathrm{𝑟𝑒𝑙}(B)`$ is defined in $`P`$,
+ $`A_P|B|_R`$, if $`\mathrm{𝑟𝑒𝑙}(B)`$ is defined in $`R`$.
Then $`PR`$ is semi-acceptable wrt. $`||`$ and $`M`$, where $`||`$ is defined as follows:
* $`|A|=|A|_P+A_P`$, if $`\mathrm{𝑟𝑒𝑙}(A)`$ is defined in $`P`$,
* $`|A|=|A|_R`$, if $`\mathrm{𝑟𝑒𝑙}(A)`$ is defined in $`R`$.
The following remark follows from Lemma 5.4 in (?) and Corollary 3.7 in (?). Together with Theorem 6.50, it implies termination of bounded queries in (?).
###### Remark 6.51.
If $`PR`$ is semi-acceptable wrt. $`||`$ and $`M`$ and $`Q`$ is bounded wrt. $`||`$ and $`M`$ then all LD-derivations of $`(PR)\{Q\}`$ are finite.
We now show that whenever Theorem 6.50 can be applied to prove termination of all the queries bounded wrt. $`||`$ and $`M`$, then also our method can be used to prove termination of the same class of queries with no need of $`||||_P`$ for relating the proofs of the two modules.
In the following theorem for the sake of simplicity we assume that $`PR`$ is a finest decomposition of $`PR`$. We discuss later how to extend the result to the general case.
###### Theorem 6.52.
Let $`P`$ and $`R`$ be two programs such that $`P`$ extends $`R`$, and let $`M`$ be a model of $`PR`$. Suppose that
* $`R`$ is semi-acceptable wrt. $`||_R`$ and $`MB_R`$,
* $`P`$ is semi-acceptable wrt. $`||_P`$ and $`M`$,
* there exists a level mapping $`||||_P`$ defined as in Theorem 6.50.
Let $`||`$ be the level mapping defined by Theorem 6.50. Moreover, suppose $`PR`$ is a finest decomposition of $`PR`$. If $`Q`$ is bounded wrt. $`||`$, then $`Q`$ is strongly bounded wrt. $`PR`$ and $`||_P`$ and $`||_R`$.
###### Proof 6.53.
Since we are considering a finest decomposition of $`PR`$, by Proposition 6.48, $`R`$ is acceptable wrt. $`||_R`$, while $`P`$ is acceptable wrt. $`||_P^{}`$ such that if $`A`$ is defined in $`P`$ then $`|A|_P^{}=|A|_P`$ else $`|A|_P^{}=0`$.
By Lemma 6.46 all LD-descendants of $`(PR)\{Q\}`$ are bounded wrt. $`||`$ and $`M`$. By definition of boundedness, for all LD-descendants $`Q^{}`$ of $`(PR)\{Q\}`$, $`first(Q^{})`$ is bounded wrt. $`||`$. By definition of $`||`$, for all atoms $`A`$ bounded wrt. $`||`$ we have that: if $`A`$ is defined in $`R`$ then $`A`$ is bounded wrt. $`||_R`$, while if $`A`$ is defined in $`P`$ then $`A`$ is bounded wrt. $`||_P`$ and hence wrt. $`||_P^{}`$ (since $`|A|_P^{}=|A|_P`$). Hence the thesis follows.
If the hierarchy $`PR`$ is not a finest one and $`||_P`$ and $`||_R`$ are the level mappings corresponding to $`P`$ and $`R`$ respectively, then we can decompose $`P`$ into a finest decomposition, $`P:=P_n\mathrm{}P_1`$ , and consider instead of $`||_P`$ the derived level mappings $`||_{P_i}`$ defined in the following way: if $`A`$ is defined in $`P_i`$ then $`|A|_{P_i}=|A|_P`$ else $`|A|_{P_i}=0`$. Similarly we can decompose $`R:=R_n\mathrm{}R_1`$ and define the corresponding level mappings. The derived level mappings satisfy all the properties we need for proving that if $`Q`$ is bounded wrt. $`||`$, then $`Q`$ is strongly bounded wrt. $`PR`$ and $`||_{P_1},\mathrm{},||_{P_n},||_{R_1},\mathrm{},||_{R_n}`$.
To complete the comparison with (?), we can observe that our method is applicable also for proving termination of queries in modular programs which are not (semi-)acceptable. Such programs clearly cannot be dealt with Apt and Pedreschi’s method. The program of Example 5.43 is a non-acceptable program for which we proved termination of all well-typed queries by applying Corollary 5.40. The following is a simple example of a non-acceptable program to which we can apply the general Theorem 3.14.
###### Example 6.54.
Let $`R`$ be the following trivial program:
| r1: | q(0). |
| --- | --- |
| r2: | q(s(Y)) $``$q(Y). |
The program $`R`$ is acceptable wrt. the following natural level mapping $`||_R`$ and any model $`M`$:
Let $`P`$ be a program, which extends $`R`$, defined as follows:
| p1: | r(0,0). |
| --- | --- |
| p2: | r(s(X),Y). |
| p3: | p(X) $``$r(X,Y), q(Y). |
The program $`P`$ is acceptable wrt. the following trivial level mapping $`||_P`$ and any model $`M`$:
| $`|𝚚(y)|_P=0`$, |
| --- |
| $`|𝚛(x,y)|_P=0`$, |
| $`|𝚙(x)|_P=1`$. |
Note that, even if each module is acceptable, $`PR`$ cannot be acceptable wrt. any level mapping and model. In fact $`PR`$ is not left-terminating: for example it does not terminate for the ground query p(s(0)). As a consequence Apt and Pedreschi’s method does not apply to $`PR`$. On the other hand, there are ground queries, such as p(0), which terminate in $`PR`$. We can prove it as follows.
* By Theorem 3.14, for all strongly bounded queries $`Q`$ wrt. $`PR`$ and $`||_R`$, $`||_P`$, all LD-derivations of $`(PR)\{Q\}`$ are finite.
* p(0) is strongly bounded wrt. $`PR`$ and $`||_R`$, $`||_P`$. In fact, $`\mathrm{𝐶𝑎𝑙𝑙}_{(PR)}(\text{p(0)})=\{\text{p(0)},\text{r(0,Y)},\text{q(0)}\}`$ and all these atoms are bounded wrt. their corresponding level mapping.
## 7 Conclusions
In this paper we propose a modular approach to termination proofs of general programs by following the proof style introduced by Apt and Pedreschi. Our technique allows one to give simple proofs in hierarchically structured programs, namely programs which can be partitioned into $`n`$ modules, $`R_1\mathrm{}R_n`$, such that for all $`i\{1,\mathrm{},n1\}`$, $`R_{i+1}`$ extends $`R_1\mathrm{}R_i`$.
We supply the general Theorem 9 which can be iteratively applied to a hierarchy of two programs and a class of queries enjoying persistence properties through LDNF-resolution. We then use such a result to deal with a general hierarchy of acceptable programs, by introducing an extension of the concept of boundedness for hierarchical programs, namely strong boundedness. Strong boundedness is a property on queries which can be easily ensured for hierarchies of programs behaving well, such as well-moded or well-typed programs. We show how specific and simple hierarchical termination proofs can be derived for such classes of programs and queries. We believe this is a valuable proof technique since realistic programs are typically well-moded and well-typed.
The simplifications in the termination proof derive from the fact that for proving the termination of a modular program, we simply prove acceptability of each module by choosing a level mapping which focuses only on the predicates defined in it, with no concern of the module context. Generally this can be done by using very simple and natural level mappings which are completely independent from one module to another. A complicated level mapping is generally required when we prove the termination of a program as a whole and we have to consider a level mapping which appropriately relates all the predicates defined in the program. Hence the finer the modularization of the program the simpler the level mappings. Obviously we cannot completely ignore how predicates defined in different modules relate to each other. On one hand, when we prove acceptability for each module, we consider a model for the whole program. This guarantees the compatibility among the definitions in the hierarchy. On the other hand, for queries we use the notion of strong boundedness. The intuition is that we consider only what may influence the evaluation of queries in the considered class.
The proof method of Theorem 9 can be applied also to programs which are not acceptable. In fact, the condition on the lower module is just that it terminates on all the queries in the considered class and not on all ground queries as required for acceptable programs. From Theorem 9 we could also derive a method to deal with pre-compiled modules (or even modules written in a different language) provided that we already know termination properties and we have a complete specification.
For sake of simplicity, in the first part of the paper we consider the notion of acceptability instead of the less requiring notion of semi-acceptability. This choice makes proofs of our results much simpler. On the other hand, as we show in Section 6, our results can be applied also to hierarchies of semi-acceptable programs.
We have compared our proposal with the one in (?). They propose a modular approach to left-termination proofs in a hierarchy of two definite programs $`PR`$. They require both the (semi)-acceptability of the two modules $`R`$ and $`P`$ wrt. their respective level mappings and a condition relating the two level mappings which is meant to connect the two termination proofs.
Our method is more powerful both because we consider also general programs and because we capture definite programs and queries which cannot be treated by the method developed in (?). In fact there are non-acceptable programs for which we can single out a class of terminating queries.
For the previous reasons our method improves also with respect to (??) where hierarchies of modules are considered. In (??) a unifying framework for the verification of total correctness of logic programs is provided. The authors consider modular termination by following the approach in (?).
In (?) a methodology for proving termination of general logic programs is proposed which is based on modularization. In this approach, the acyclic modules, namely modules that terminate independently from the selection rule, play a distinctive role. For such modules, the termination proof does not require a model. In combination with appropriate notions of up-acceptability and low-acceptability for the modules which are not acyclic, this provides a practical technique for proving termination of the whole program. Analogously to (?), also in (?) a relation between the level mappings of all modules is required. It is interesting to note that the idea of exploiting acyclicity is completely orthogonal to our approach: we could integrate it into our framework.
Another related work is (?), even if it does not aim explicitly at modularity. In fact they propose a technique for automatic termination analysis of definite programs which is highly efficient also because they use a rather operational notion of acceptability with respect to a set of queries, where decreasing levels are required only on (mutually) recursive calls as in (?). Effectively, this corresponds to considering a finest decomposition of the program and having independent level mappings for each module. However, their notion of acceptability is defined and verified on call-patterns instead of program clauses. In a sense, such an acceptability with respect to a set of queries combines the concepts of strongly boundedness and (standard) acceptability. They start from a class of queries and try to derive automatically a termination proof for such a class, while we start from the program and derive a class of queries for which it terminates.
In (?) termination in the context of tabled execution is considered. Also in this case modular results are inspired by (?) by adapting the notion of acceptability wrt. call-patterns to tabled executions. This work is further developed in (?) where their modular termination conditions are refined following the approach by (?).
In (?) a method for modular termination proofs for well-moded definite programs is proposed. Our present work generalizes such result to general programs.
Our method may help in designing more powerful automatic systems for verifying termination (????). We see two directions which could be pursued for a fruitful integration with existing automatic tools. The first one exploits the fact that in each single module it is sufficient to synthesize a level mapping which does not need to measure atoms defined in other modules. The second one concerns tools based on call-patterns analysis (???). They can take advantage of the concept of strong boundedness which, as we show, can be implied by well-behavior of programs (??).
#### Acknowledgements.
This work has been partially supported by MURST with the National Research Project “Certificazione automatica di programmi mediante interpretazione astratta”.
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# An Exploration of the Paradigm for the 2-3 Hour Period Gap in Cataclysmic Variables
## 1 Introduction
Cataclysmic variables (CV) are short period binary systems consisting of a white dwarf that accretes matter via Roche-lobe overflow from a low-mass companion star. These objects exhibit a wide range of phenomenology including optical flickering in nova-like systems, dwarf nova eruptions which are thought to be caused by thermal instabilities in the accretion disks, and classical nova explosions which are thermonuclear runaways of the accreted matter on the white dwarf (see, e.g., Warner 1995). The range of observed phenomena depends on the mass transfer rate, the mass ratio of the stellar components, and the magnetic field strength of the accreting white dwarf. The orbital periods of the majority of CVs range from 8 hours down to about 78 minutes, but both longer and shorter period systems are known. In the former case, the donor stars are typically somewhat evolved, while in the latter case, the donor stars are hydrogen exhausted. In this paper we focus on the gap that exists in the orbital period distribution of CVs in the range of $`23`$ hr (see, e.g., Warner 1976; Rappaport, Verbunt, & Joss 1983, hereafter RVJ; Spruit & Ritter 1983; Hameury, et al. 1988a; Warner 1995).
The overall evolution of CV binaries is thought to be fairly well understood. The widely accepted explanation for the period gap rests on a mechanism for extracting angular momentum from the binary orbit (e.g., via magnetic braking of the secondary) for periods down to $`3`$ hours, followed by a relatively substantial decrease in the angular momentum loss rate<sup>1</sup><sup>1</sup>1 We note that this scenario does not require the angular momentum loss rate to drop suddenly. Instead, it requires only that the timescale over which the angular momentum loss rate decreases must be shorter than the thermal timescale of the donor.. The donor star, which had been thermally “bloated” in response to the mass loss driven by the systemic angular momentum losses, is then able to relax inside of its Roche lobe and mass transfer ceases. The donor star is then thought to reestablish Roche-lobe contact by the time the orbital period has decreased to about 2 hr, after which mass transfer resumes. In this paper we critically examine this paradigm for the creation of the period gap. While most workers believe in the existence of the so-called “2-3 hr period gap”, a few (e.g., Wickramasinghe & Wu 1994; Verbunt 1997, but see also Warner 1995; Wheatley 1995) have questioned its reality, especially when all types of CV are considered; however, we adopt the view that the period gap is a real feature of the CV population as a whole and, as such, requires a theoretical explanation (with observational tests) within the context of their binary evolution. Finally in this regard we note a suggestion by Clemens et al. (1998) that the period gap results from a “kink” in the radius-mass relation for main-sequence stars at a mass of about $`0.25M_{}`$ (but see the rebuttal by Kolb, King, & Ritter 1998).
In §2 we describe the conventional picture of the evolution of a typical CV, including the period gap, and show some illustrative examples of binary evolution calculations for individual systems. In §3 we explore how the binary evolution alters the relations among mass, radius, and orbital period of the secondary star. Specifically we discuss how the main-sequence radius-mass relationship must be modified to include the addition of a “bloating factor” that accounts for the changes caused by departures from thermal equilibrium of the mass-losing secondary star. We derive semi-analytic mass-period and radius-period relationships for CV secondaries. In §4 we describe our population synthesis and binary evolution codes, while in §5 we present results from our population synthesis study of CVs in which the binary parameters of the CVs at all phases of their evolution are explored. In §6 we show how assumptions that the donor star has a main-sequence radius-mass relation can lead to large errors in the assignment of the constituent stellar masses, most notably within the orbital period range of 3-5 hr. This period range should encompass the maximum bloating exhibited by a CV secondary compared to a main-sequence star of the same mass. Also in §6 we discuss some specific observational implications resulting from our theoretical work. In particular, a specific test for CVs just above the period gap which will enable us, in principle, to distinguish unambiguously among different possible explanations for the period gap is presented. Finally, we present our summary and conclusions in §7.
## 2 Standard Evolutionary Scenario for CVs
In the conventional picture of CV evolution (see, e.g., Faulkner 1971; Paczyński & Sienkiewicz 1981; Rappaport, Joss, & Webbink 1982, hereafter RJW; RVJ; Spruit & Ritter 1983; Hameury et al. 1988a; Kolb 1993), the early phases are expected to be dominated by angular momentum losses due to magnetic braking via a magnetically constrained stellar wind from the donor star (see, e.g., Verbunt & Zwaan 1981; RVJ). In these early phases, mass transfer rates are typically $`10^9`$ to $`10^8M_{}\mathrm{yr}^1`$, and orbital periods range from $``$ 8 hr to $``$ 3 hr, just at the upper edge of the period gap. At some point in the evolution, the secondary becomes completely convective (at $`0.23M_{}`$) and, in the currently accepted view, magnetic braking is assumed to be greatly reduced. The near cessation of magnetic braking reduces the mass transfer rate, and allows the secondary to shrink toward its thermal equilibrium radius. This causes a period of detachment, during which $`\dot{M}`$ drops to essentially zero, and which lasts until the Roche lobe shrinks sufficiently to bring the secondary back into contact with it, at an orbital period of $`2`$ hr. This is the commonly accepted explanation for the observed period gap between 2-3 hrs in CVs (RVJ; Spruit & Ritter 1983).
When mass transfer recommences at $`P_{orb}`$ $``$ 2 hrs, it is then driven largely by gravitational radiation losses at rates of $`10^{10}M_{}\mathrm{yr}^1`$. As the orbit shrinks and the mass of the donor star decreases, the mass-loss timescale increases, but the thermal timescale, $`\tau _{KH}`$, increases much faster, due to its approximate $`M^2`$ dependence. Therefore, at some point the thermal timescale grows larger than the mass-transfer timescale. When this occurs, the donor star is unable to adjust to the mass loss on its thermal timescale, and it therefore starts to expand upon further mass loss, in accordance with its adiabatic response; i.e., $`[\mathrm{dln}(\mathrm{R})/\mathrm{dln}(\mathrm{M})]_{ad}<0`$. The orbital period at this point is typically $`80`$ min and the mass of the donor star is $`0.06M_{}`$. From this point on, the mass of the donor star will continue to decrease (but with longer and longer $`\dot{M}`$ timescales), the orbital period will increase back toward periods approaching $``$2 hr (within a Hubble time), and the interior of the donor star will become increasingly electron degenerate. A discussion of this later stage of CV evolution is presented by Howell, Rappaport, & Politano 1997; hereafter HRP).
To make these evolutionary descriptions somewhat more quantitative, we show in Figure 1 the secular evolution of several model CVs under the influence of magnetic braking and gravitational radiation. The evolution code used to generate these results is a descendant of the one used by RVJ, and is described in §4.2 along with recent improvements to the code. The two panels on the left side of Fig. 1 show the evolution with time of a CV binary with initial constituent masses of $`M_2=0.9M_{}`$ and $`M_{WD}=1.1M_{}`$, where $`M_2`$ and $`M_{WD}`$ are the masses of the donor star and white dwarf, respectively. Other parameters used in the calculation are for our “Standard Model” (see Table 1 for definitions). The top and bottom panels show the evolution with time of the mass transfer rate and orbital period, respectively, for an assumed donor star with solar composition. The calculations have been carried out to approximately the age of the Galaxy. The evolutionary phases and features discussed above are present in Fig. 1, including the interval where mass transfer is driven by magnetic braking ($`10^{7.3}10^{8.4}`$ yr), the period gap ($`10^{8.4}10^{8.8}`$ yr), the interval where $`\dot{M}`$ is driven by gravitational radiation losses ($`>10^{8.8}`$ yr), the period minimum at $`10^{9.4}`$ yr, the subsequent increase in $`P_{orb}`$ back up to $`2`$ hrs, and the sharp falloff in $`\dot{M}`$ after orbital period minimum.
On the right side of Fig. 1 the temporal evolution of four other illustrative model CV binaries are shown. The following discussion contains descriptions of the period gap which develops in these systems; these are easier to visualize by looking also at Fig. 2. The initial masses ($`M_2,M_{WD}`$) of these systems are (0.2, 0.4), (0.35, 0.35), (0.3, 0.6), and (0.65, 0.7), all in units of $`M_{}`$. For the system with initial masses (0.2, 0.4) the binary comes into Roche-lobe contact for the first time at an orbital period below the gap, i.e., at $`P_{orb}=2`$ hr (solid curve in Figs. 1 c and d). Note the enhanced mass transfer rate at $``$ 30 Myr after Roche-lobe contact is made. The subsequent evolution is not dissimilar to the one shown in the left panels. For the system with initial masses (0.35, 0.35), the donor star commences mass transfer at a period of 3.3 hr, with magnetic braking still operative (dotted curve in Figs. 1 c and d). Because the two masses are the same when the donor star first fills its Roche lobe, the mass transfer is only marginally stable (see the discussion in §4.2 below). Therefore, $`\dot{M}`$ is initially very high and the system is quickly driven out of thermal equilibrium, causing the orbit to expand. This system comes out of contact (i.e., enters the period gap) at an orbital period of 3.5 hr. The system with initial masses (0.3, 0.6) is an example of one that commences mass transfer in the period gap. Lastly, the system with initial masses (0.65, 0.7) is another example of a system which exhibits the “usual” 2-3 hr period gap, but commences mass transfer at $`P_{orb}=5`$ hr.
In Figure 2 the same evolutions shown in Fig. 1 are again presented, but this time the binary parameters are displayed as a function of the evolving orbital period. As in Fig. 1, the left panels are for initial masses ($`M_2,M_{WD}`$) of (0.9, 1.1), while the right panels are for initial masses of (0.2, 0.4), (0.35, 0.35), (0.3, 0.6), and (0.65, 0.7). The top, middle, and bottom panels show the evolution of $`\dot{M}`$, $`M_2`$, and $`R_2`$, respectively. As mentioned above, the period gap is more evident in Fig. 2 than it is in Fig. 1. We note here several unique features associated with the evolution of individual CVs; an understanding of these features will aid our interpretation of the results obtained for an entire population of evolving CV systems (see §5). For example, $`\dot{M}`$ typically exhibits a sharp spike at the onset of mass transfer, (see also RVJ and Hameury et al. 1988b); this behavior will appear in all of the two-dimensional “images” we produce from the population synthesis calculations in §5. The mass of the donor stays constant during its evolution through the period gap since there is no mass transfer taking place at that time - this is indicated by the horizontal lines in the middle panels. The abrupt shift in location between the $`M_2P_{orb}`$ track above the period gap and below the gap will be dramatically apparent in the population synthesis results, and will have important consequences that are discussed below. Finally, the radius of the donor star decreases sharply after the system enters the period gap; in fact, it is the shrinking of the donor inside of its Roche lobe when the magnetic braking ceases that is the putative cause of the period gap. Again, the abrupt shift between the $`R_2P_{orb}`$ track above and below the period gap will be very pronounced in the population synthesis results.
A noteworthy feature of Figures 1a, 1c, 2a, and 2d mentioned above is the sharp rise in $`\dot{M}`$ whenever mass transfer has just commenced, including the first time that the donor star fills its Roche lobe, and after the resumption of mass transfer below the period gap. This results from the fact that when a low-mass star (i.e., $`0.5M_{}`$) is in thermal equilibrium (i.e., the nuclear luminosity, $`L_{nuc}`$, equals the bolometric luminosity, $`L_{opt}`$), the sudden onset of mass transfer forces the star to expand because its adiabatic index is negative (discussed above). This expansion can cause a temporarily anomalously high rate of mass transfer, viz, the episodes of high $`\dot{M}`$ seen in Figures 1a, 1c, 2a, and 2d. However, as soon as the donor star expands, its core temperature drops slightly, and $`L_{nuc}`$, which is a highly sensitive function of temperature, drops dramatically. This leads to a luminosity deficit wherein $`L_{nuc}<L_{opt}`$. The star can then lose a net amount of energy, shrink, and approach its new thermal equilibrium radius (appropriate to its lower mass) on a Kelvin-Helmholtz (i.e., thermal) timescale. During the mass loss process, true thermal equilibrium is never reached, and the luminosity deficit attains a value which is adequate to allow the star to shrink continuously. The above discussion explains the transient episodes of higher transfer rates at the start of mass transfer epochs, and the “outlying” lower probability CV states we shall encounter in §5. It also explains the thermal “bloating” of the donor star which is discussed in §3, and which will play a key role in the observational test we propose in §6. (For earlier discussions of some of these basic effects, see RJW and Hameury et al. 1988a.)
The five individual evolutions shown in Figs. 1 and 2 serve to illustrate the range of interesting possibilities for CVs which commence mass transfer with different mass ratios. The population synthesis study described in §5, explores these various possibilities in a more systematic and complete way.
## 3 Quantitative Effects of Thermal Bloating of the Secondary Star
We start with the assumption that during mass transfer in a CV the Roche lobe of the donor star is located within its atmosphere, i.e., the donor star is “filling” its Roche lobe (see Howell et al. 2000). We then take the Roche-lobe radius of the secondary star to be given by the simple expression of Kopal (1959):
$`R_20.46a\left({\displaystyle \frac{M_2}{M_2+M_{WD}}}\right)^{1/3}.`$ (1)
This can be combined with Kepler’s 3rd law to yield the well-known relationship among the mass, radius, and orbital period of the donor star:
$`P_{orb}(M_2,R_2)9M_{2}^{}{}_{}{}^{1/2}R_{2}^{}{}_{}{}^{3/2},`$ (2)
where $`M_2,R_2`$ and $`P_{orb}`$ are expressed in units of $`M_{}`$, $`R_{}`$, and hours, respectively. If we now assume that the radius of the donor star is some factor $`f`$ times the radius it would have if it were a normal main-sequence star, we can write
$`R_2=faM_{2}^{}{}_{}{}^{b},`$ (3)
where we approximate the radius-mass relation for stars on the lower main sequence (i.e., G to M stars) by $`R_2=aM_{2}^{}{}_{}{}^{b}`$ where $`a`$ and $`b`$ are constants, and we refer to $`f`$ as the “bloating factor”. This bloating factor $`f`$ is simply a measure of how much larger the radius of a CV secondary is than that of a single, main-sequence star of the same mass due to the departure from thermal equilibrium. We can now combine equations (2) and (3) to derive relations for the mass and radius of CV secondaries as a function of the binary orbital period.
$`M_29^{2/(3b1)}P_{orb}^{}{}_{}{}^{2/(3b1)}\left(af\right)^{3/(3b1)}`$ (4)
$`R_29^{2b/(3b1)}P_{orb}^{}{}_{}{}^{2b/(3b1)}\left(af\right)^{1/(3b1)}.`$ (5)
For the purposes of this exercise, we take $`a`$ = 0.85 and $`b`$ = 0.85, which we find by fitting a power law to the main-sequence models of Dorman, Nelson, & Chau (1989; hereafter DNC). With these values for the constants $`a`$ and $`b`$, the above equations simplify to:
$`M_2(P_{orb})0.08f^{1.95}P_{orb}^{1.3}`$ (6)
$`R_2(P_{orb})0.10f^{0.65}P_{orb}^{1.1}.`$ (7)
where, again, $`M_2`$ and $`R_2`$ are in solar units and $`P_{orb}`$ is in hours. The conclusions drawn from these expressions are somewhat counterintuitive in that, for a CV at a given orbital period, if the donor star is bloated, the proper radius and mass that should be inferred from the orbital period are smaller than the values that would be inferred if the star were on the main sequence (see also Beuermann et al. 1998). In §6 we derive polynomial fits for $`M_2(P_{orb})`$ and $`R_2(P_{orb})`$ from our population synthesis study; the analytic expressions given by equations (6) and (7) serve mainly to demonstrate how these quantities scale with the bloating factor $`f`$.
## 4 Population Synthesis Study
The individual binary evolution runs shown in Figs. 1 and 2 for several different combinations of initial constituent masses are instructive, but do not (i) adequately sample the full range of possible initial masses, nor (ii) provide us with the distributions of CV binary properties at the current epoch. We have therefore undertaken a population synthesis study of CVs which consists of two parts. In the first part, we utilize a Monte Carlo approach to select a large number ($`3\times 10^6`$) of primordial binaries, and follow the evolution of these systems to see which ones undergo a common envelope phase. In such events, the envelope of the giant star engulfs the secondary, leading to a spiral-in episode which leaves the secondary in a close orbit with a white dwarf (i.e., the core of the primary star; see, e.g., (Paczyński 1976; Webbink 1979; see also §4.1 for details). Primordial binaries which are too wide will not undergo any significant mass transfer and will not lead to the formation of CV systems - the evolution of such wide binaries is not followed in the present study. Successful systems which emerge from the first part of our population synthesis caculations are those which do undergo a common envelope phase and yield a close binary consisting of a white dwarf and low mass ($`1M_{}`$) companion. The second part of the population synthesis considers those white-dwarf main-sequence binaries for which systemic angular momentum losses, or a modest amount of evolution by the normal companion star, can initiate Roche-lobe contact within a Hubble time. Each of these systems is then evolved in detail through the mass-transfer phase (CV phase) until the donor star has been reduced to a negligible mass (typically 0.03 $`M_{}`$).
A number of prior population synthesis studies of cataclysmic variables have been carried out. These include work by Politano (1988; 1996); de Kool (1992); Kolb (1993); Di Stefano & Rappaport (1994, for CVs in globular clusters); and HRP (emphasizing systems which evolve beyond the orbital period minimum). The current study has several new features and advantages over the previous studies. First, we compute probability density functions in two parameters, e.g., $`\dot{M}:P_{orb}`$ and $`M_2:P_{orb}`$ (see §4.3). This way of studying and evaluating the results of population synthesis calculations has a distinct advantage over producing distributions of a single parameter. For example, we are able to quantitatively evaluate phases of the evolution that are short lived or represent unusual evolutionary states, (e.g., whenever Roche-lobe contact has just been established or when the initial binary mass ratio is near unity). Other examples include the ability to discern the spread in $`\dot{M}`$ at a given orbital period, the distinction between systems with He and CO white dwarfs, and the pronounced depression in secondary mass at a given orbital period (for systems just above the period gap). A second advantage is that our code for evolving the donor stars was originally developed to evolve brown dwarfs of very low mass and to very old ages. The code has been well “calibrated” against other more sophisticated ones that have been used for the purpose of evolving brown dwarfs (see §4.1). Finally, our population synthesis code which is used to generate the zero-age CVs that are input to the binary evolution code provides an independent check on previous work, and tests the sensitivity of our conclusions to various uncertainties in the physics, initial conditions, and other input parameters.
Finally, we mention a possibly important limitation on the study we present here, which also applies to most other prior work in this area. We have considered only donor stars with $`M_21M_{}`$, and do not allow for the chemical (nuclear) evolution the donor. The latter approximation is realistic if the donor star commences mass transfer within $``$3$`\times `$10<sup>9</sup> years of the common envelope event, or if the donor has a mass of $`0.7M_{}`$. These conditions apply to most of the systems that successfully evolve through the CV phase in our calculations. Furthermore, we find that only $``$ 5% of all the stable mass-transferring, zero-age CVs in our population synthesis study have secondaries which are older than 1/3 of their main-sequence lifetime prior to the start of mass transfer. Theoretically, there should indeed be some CVs which evolve from donor stars that are initially more massive than $`1M_{}`$, and they should be followed in future population synthesis studies. For the present study, we simply assume that such systems, with donors whose initial mass exceeds 1 $`M_{}`$, do not contribute substantially to the CV population and, above all, would not affect our conclusions concerning systems near the period gap.
In this regard, recent work by Beuermann et al. (1998) examines the properties of the secondary star in CVs in an effort to determine if they are indeed similar to normal main-sequence stars. They show that, in the spectral type-$`P_{orb}`$ plane, the $`50`$ CVs with measured spectral types lie below the expected relation for main-sequence stars (i.e., they are cooler at a fixed value of $`P_{orb}`$). Beuermann et al. conclude that for systems with $`P_{orb}<6`$ hr this effect could result from mass loss (see §5), but that for longer orbital periods this effect suggests chemical evolution of the donor star. This is a potentially important finding for systems with orbital periods longer than we consider here, and could also possibly impact the shorter period systems as well. ¿From our population synthesis results, we find that $`10\%`$ of CVs could potentially form with progenitors whose mass is initially sufficiently high (i.e., $`1M_{}`$) that chemical evolution of the donor star might indeed be significant. Systems with such donor stars are not followed in the present study. If, for some as yet unknown reason, the more massive donor stars have a greater efficiency for producing CVs than their lower-mass counterparts, then chemical evolution may indeed prove influential in the evolution of CVs. These possibilities should be examined in future population synthesis studies.
### 4.1 Choosing the Zero-Age CVs
The properties of the primordial binary systems are chosen via Monte Carlo techniques as follows. The primary mass is picked from Eggleton’s (2000) Monte Carlo representation of the Miller & Scalo (1979) IMF,
$`M_1(x)=0.19x[(1x)^{3/4}+0.032(1x)^{1/4}]^1,`$ (8)
where $`x`$ is a uniformly distributed random number. This distribution flattens out toward lower masses, in contrast with a Salpeter-type power-law IMF (1955). We considered primary stars whose mass is in the range of $`0.8<M_1<8M_{}`$. Next, the mass of the secondary, $`M_2`$, is chosen from the probability distribution, $`f(q)=5/4q^{1/4}`$, where $`qM_2/M_1`$ (Abt & Levy 1978; but also see Halbwachs 1987). This mass ratio distribution is, at best, poorly known empirically. Our adopted distribution has the property that the mass of the secondary is correlated with the mass of the primary, but is not strongly peaked toward $`q=1`$. We find that our results are not very sensitive to the choice of $`f(q)`$, unless an extreme is adopted such as the assumption that the two masses are to be chosen completely independently of one another (see, e.g., Rappaport, Di Stefano, & Smith 1994; hereafter RDS - their Table 2 and Fig. 4). Secondary masses as small as $`0.09M_{}`$ are chosen (we wanted to ensure that only stars with masses clearly above the minimum main-sequence mass are included). To choose an initial orbital period, a distribution that is uniform in log$`(P)`$ over the period range of $`1`$ day to $`10^6`$ years is used (see, e.g., Abt & Levy 1978; Duquennoy & Mayor 1991). Since we consider only circular orbits the adopted orbital period distribution more properly pertains to the tidally circularized orbits than to the initial orbits of the primordial binaries.
After the masses and orbital period are chosen, the orbital separation is calculated using Kepler’s 3rd law. We utilize an analytic expression for the relation among the core mass, the radius, and the total mass of the primary to estimate the mass of the degenerate core, $`M_{WD}`$, when the primary fills its Roche lobe. The expression we used for this purpose (see RDS) was designed to reproduce the features of Figure III.2 of Politano (1988) and Figure 1 of de Kool (1992), except that the core-mass radius relation for stars with mass $`2M_{}`$ was renormalized to match the fitting formula of Eggleton (2000; see eq. of Joss, Rappaport, & Lewis 1987). Mass loss via a stellar wind prior to the start of the first mass-transfer phase was computed via an analytic expression derived by M. Politano (1999, private communication). In practice, the inclusion of this wind mass loss does not significantly affect the results.
In order to select only systems which undergo a common envelope phase we require that the radius of the Roche lobe of the primary be larger than the radius of a star of mass $`M_1`$ at the base of the giant branch (see, e.g., Paczyński 1965; Webbink 1979, 1985, 1992, de Kool 1992, and references therein). This ensures that unstable mass transfer will occur on a timescale that is substantially shorter than a thermal time, and should lead to a common envelope phase. Once mass transfer from the primary to the secondary commences, we assume that a common envelope phase occurs and compute the final spiral-in separation based on simple energetic considerations (see, e.g., Taam, Bodenheimer, & Ostriker 1978; Meyer & Meyer-Hofmeister 1979; Livio & Soker 1988; Webbink 1992; RDS; Taam & Sandquist 1998). The expression we use for determining $`a_f`$, the final orbital separation after spiral-in, is given by:
$`{\displaystyle \frac{ϵGM_2}{2}}\left({\displaystyle \frac{M_{core}}{a_f}}{\displaystyle \frac{M_1}{a_i}}\right)={\displaystyle \frac{GM_{env}(M_{env}+3M_{core})}{R_1}}`$ (9)
where $`M_{core}M_{WD}`$ and $`M_{env}`$ are the core and envelope masses of the primary, $`R_1`$ is the radius of the primary, $`a_i`$ is the initial orbital separation, and $`ϵ`$ is the energy efficiency factor for ejecting the envelope. We take $`ϵ`$ to have a value of 1.0 in our standard model. The two terms in parentheses on the right hand side of equation (9) represent the binding energy of the envelope of the primary to itself and to its core. The dimensionless coefficients multiplying each term were computed for an assumed polytropic envelope structure with polytropic index $`n=3.5`$ (RDS). For other similar values of $`n`$ the ratio of $`3:1`$ between the two coefficients is roughly the same. We assume that the duration of the spiral-in is sufficiently short ($`<10^4`$ yr; see above references) that the mass of the secondary does not change significantly during the common envelope phase.
After the spiral-in episode, the separation, the white dwarf mass, the secondary mass, and the corresponding Roche-lobe radius of the secondary are known. If at the end of the common envelope phase the secondary would already be overfilling its Roche lobe, then we eliminate the system. (In most cases, this circumstance would be expected to lead to a merger of the secondary star with the degenerate core of the primary, which presumably would result in the formation of a giant star.) In practice, if the Roche lobe is larger than $`20R_{}`$, then neither magnetic braking nor gravitational radiation would bring the system into Roche-lobe contact before the secondary would evolve past the base of the giant branch. We can also eliminate these systems since the ensuing mass transfer would either be dynamically unstable (see, e.g., Paczyński 1967; Kippenhahn, Kohl, & Weigert 1967; Webbink 1979, 1992) or lead to an even wider orbit; either case would not produce a CV of the ordinary kind.
We typically start with $`35\times 10^6`$ primordial binaries and end up with $`15,000`$ pre-CVs to evolve through the mass-transfer phase with the bipolytrope evolution code described in the next section. The computational time for this first portion of the calculations is negligibly short.
### 4.2 Evolving the CVs Through Their Mass-Transfer Phase
As mentioned earlier, the evolutionary tracks of CV systems are calculated using a version of the code that was first developed by RVJ (see also RJW) to explore the effects of the parameterized Verbunt & Zwaan (1981) magnetic braking law on the evolutionary properties of cataclysmic variables. According to their algorithm, the mass losing donor is approximated by a bipolytrope wherein the convective envelope is represented by an n=3/2 polytrope and the radiative core by an n=3 polytrope. One of the advantages of this code is that it allows for the rapid computation of a large number of evolutionary tracks and provides a more physically intuitive interpretation of the results. The original version of the code has been modified substantially to allow for improvements to the input physics, and to ensure that the conditions near the surface (atmosphere) are more physically realistic. A number of these changes have been discussed in previous papers.
The most significant of these modifications and updates are described by Nelson, Rappaport, & Joss (1986a, 1986b; 1993) who used a single polytrope model to follow the evolution of fully convective low-mass stars and brown dwarfs. The results of the brown dwarf cooling evolutions and the calculation of ZAMS models of low-mass stars are in excellent agreement with those calculated using more sophisticated techniques (see, e.g., DNC; Burrows et al. 1993; Burrows et al. 1997; Baraffe et al. 1998, and references therein). Specifically, coulombic corrections to the pressure equation of state were incorporated and an updated version of the Alexander, Johnson, & Rypma (1983; Alexander ) low-temperature, radiative (surface) opacities was used. The molecular hydrogen partition function was also calculated more accurately. Most importantly, the specific entropy at the surface was matched directly to the specific entropy in the interior; i.e., at the interface between the radiative core and the convective envelope.<sup>2</sup><sup>2</sup>2A small entropy mismatch was introduced to correct for thin regions of superadiabatic convection/radiative transport that exist beneath the photosphere of the more massive stars in our mass range. These corrections depend on the assumed value of the mixing length parameter and were chosen so as to provide the best possible representations of ZAMS stars. They were largest for the $`1.0M_{}`$ model ($`5\%`$ of the specific entropy), decreasing to zero for fully convective stars.
In addition to these changes, the atmospheric pressure boundary condition was modified so as to approximate more closely the scaled solar T-$`\tau `$ (Krishna-Swamy 1972) relation. The radiative surface opacities did not include the effects of grain formation. Since grains can only form in the atmospheres of very low-temperature stars ($``$ 1500 K), this should affect mostly the evolution of those CVs that have evolved beyond the orbital period minimum. However, we have found that the evolution of CVs through and beyond the period minimum, is not particularly sensitive to this omission.
The overall result of all of these changes is that the theoretical radius-mass relation for our ZAMS models with masses $`1.0M_{}`$ is now in substantial agreement with other theoretical calculations as well as with observational studies of low-mass stars (see DNC). For similar abundances of hydrogen and for stars of approximately solar metallicity, we find that the radii of our new models compared with other theoretical models (and the DNC results) typically agree to within an rms error of $`3\%`$ ($`M1.0M_{}`$). Deviations among the theoretical models are greatest for the higher mass stars due to uncertainties in the mixing length parameter and the treatment of inefficient (superadiabatic) convection. When observations of double stars are considered, we believe that our ZAMS radii are accurate to within $`5\%`$. Our ZAMS models become fully convective at a mass of $`0.34M_{}`$. This is considerably smaller than the value given in RVJ but agrees well with the DNC results (as well as with newer generations of models).
Mass transfer in CVs is driven by angular momentum losses due to gravitational radiation (Landau & Lifshitz, 1962) and other systemic angular momentum losses such as “magnetic braking”. The magnetic braking law that we utilize is that of Verbunt & Zwaan (1981) and parameterized by RVJ. The magnetic braking parameters were chosen so as to best reproduce the observed period gap. According to the parameterization described in RVJ, we took $`\gamma =3`$ and did not adjust the multiplicative constant (defined here as $`C_{MB}`$) used in the RVJ prescription. We also “shut off” magnetic braking when the radiative core had been reduced to less than 15% of the mass of the donor. Magnetic braking is assumed to be greatly reduced as a result of the restructuring of the magnetic field of the donor star when it becomes nearly fully convective. This reduction in the angular momentum loss rate gives the donor an opportunity to shrink inside of its Roche lobe on a thermal timescale. Further angular momentum losses due to gravitational radiation cause mass transfer to recommence once the binary system is brought back into a state of semi-detachment (see, e.g., RVJ; Spruit & Ritter 1983; Hameury et al. 1988a for a more detailed explanation). As pointed out in several places in this work, the actual mechanism that produces the bloating of the donor, and the means by which mass transfer is interrupted, are not central to the conclusions drawn in this paper. What is important in this regard is that the bloating be sufficiently large as to produce the observed width of the period gap. For our standard evolutionary model the period gap covers the range of 2.1 $`<P_{orb}<`$ 2.85 hr. According to Warner (1995), this synthetic gap approximates the observed one very well.
We assume that mass and orbital angular momentum are lost as a result of nova explosions on the surface of the white dwarf accretor. For our standard model we assumed that all of the mass that is accreted by the white dwarf is lost with the same specific angular momentum as the white dwarf itself (see Schenker, Kolb & Ritter 1992). Given the relatively low mass transfer rates, it is likely that the nova events are extremely hydrodynamic, and thus it is unlikely that any of the accreted mass actually contributes to increasing the mass of the white dwarf (see, e.g., Prialnik & Kovetz 1995, Starrfield 1998, and references therein).
After a potential cataclysmic variable system has been generated with the population synthesis code, the two detached components are given the opportunity to come into contact, via magnetic braking, within the age of the Galaxy (minus the CV formation time). However, the initial mass transfer may actually be unstable, thereby leading to a common envelope phase (and the ultimate demise of the binary system). As derived by RJW, the expression for the long-term mean mass transfer rate in a CV is given by $`|\dot{M}|/M`$ = $`N/D`$, where the numerator, $`N`$, contains the drivers of mass transfer, e.g., systemic angular momentum losses, and the thermal expansion/contraction of the donor star (see eq. in RVJ). The denominator is given by
$`D=\left[\left({\displaystyle \frac{5}{6}}+{\displaystyle \frac{\xi _{ad}}{2}}\right){\displaystyle \frac{(1\beta )q}{3(1+q)}}(1\beta )\alpha (1+q)\beta q\right]`$ (10)
where $`qM_2/M_{WD}`$ (note that this is the inverse of the definition used in RVJ), $`\beta `$ is the fraction of the mass lost by the donor star that is ultimately retained by the white dwarf, $`\alpha `$ is the specific angular momentum carried away by matter ejected from the binary system in units of the binary angular momentum per unit reduced mass, and $`\xi _{ad}`$ is the adiabatic index of the donor star, i.e., $`[\mathrm{dln}(\mathrm{R})/\mathrm{dln}(\mathrm{M})]_{ad}`$. For our Standard Model (see Table 1), we take $`\beta =0`$ (i.e., all the mass accreted by the white dwarf is eventually ejected in nova explosions<sup>3</sup><sup>3</sup>3See Schenker et al. (1998) for a justification as to why it is valid to approximate the ejection of mass in a series of nova explosions with a constant value of $`\beta =0`$.), and $`\alpha =M_{2}^{}{}_{}{}^{2}/\left(M_2+M_{WD}\right)^2`$. With these definitions, the above equation reduces to
$`D={\displaystyle \frac{5}{6}}+{\displaystyle \frac{\xi _{ad}}{2}}{\displaystyle \frac{q(1+3q)}{3(1+q)}}`$ (11)
As discussed by RJW, stable mass transfer requires $`N>0`$ and $`D>0`$. As an example, consider donor stars with $`M_2<0.3M_{}`$ and $`\xi _{ad}`$ = -1/3. In this case, stability (based on eq. ) requires that $`M_2<M_{WD}`$. This allows for considerably larger values of $`M_2`$ than the more conventional limit for conservative transfer where $`M_2<2/3M_{WD}`$ is required for stable mass transfer (with low-mass unevolved donors). Thus, the mass ratios that appear in our population synthesis can often approach unity or exceed it.
### 4.3 Generating the Population Synthesis Tracks
We define a birth rate function, $`BRF(t)`$, for the progenitor primordial binaries, where $`t`$ is the elapsed time between the formation of the Galaxy and the birth of the primordial binary. If a binary is born at time $`t`$, then an additional time $`\tau _{\mathrm{prim}}`$ must elapse before the primary evolves to the point where a common envelope phase may occur (see §4.1). We define this time ($`t+\tau _{\mathrm{prim}}`$) to be the birth time of the incipient CV. The resultant zero-age CV is then evolved in the binary evolution code for a total time $`t_{max}`$ = ($`10^{10}`$ \- $`\tau _{\mathrm{prim}}`$) yr, which is the maximum time any CV that is descended from a similar primordial binary could evolve before the current epoch. (The binary evolution code starts with the white dwarf and companion star as they emerge from the common envelope, so the elapsed time, $`t_{ev}`$, includes the interval before the donor star fills its Roche lobe.) At each step in the evolution code, specified by time $`t_{ev}`$ (with respect to the first time step in the code), we sum in discrete binned arrays for various combinations of $`P_{orb}`$, $`M_2`$, $`M_{\mathrm{WD}}`$, $`q`$, $`\dot{M}`$, $`T_{eff}`$, and $`L_2`$, the following quantity, $`\mathrm{\Delta }Q`$:
$`\mathrm{\Delta }Q={\displaystyle \frac{\mathrm{\Delta }t\times BRF(10^{10}\tau _{\mathrm{prim}}t_{ev})}{N}}.`$ (12)
In this expression, the argument of $`BRF`$ is the time that the primordial binary was born with respect to the formation of the Galaxy, $`\mathrm{\Delta }t`$ is the time interval for that particular step in the evolution run, and $`N`$ is the total number of systems that are selected to start the population synthesis run. For all of the population synthesis runs in this study, the $`BRF`$ was taken to be constant in time. Even though we have adopted a constant stellar birth rate per unit time, the method we use for generating the CV population at the current epoch is completely general (see also Kolb 1993).
The net result of this procedure is that the sum of the $`\mathrm{\Delta }Q`$s at the end of the population synthesis run, in any particular bin, represents the number of CVs at the current epoch with that particular parameter value.
## 5 Population Synthesis Results
The computed population of current-epoch CVs as generated by the above techniques is displayed as a sequence of color images in Figures 3 through 7. In Figure 3 we show the model CV population in the $`\dot{M}P_{orb}`$ plane for our standard model (cf. Fig. 2a). The image is generated in such a way that the color reflects the logarithm of the number of current-epoch CVs at a particular location in the $`\dot{M}P_{orb}`$ plane. In each of the images the color scale is located on the right side. The image in Fig. 3 is comprised of 100 pixels per hour interval in $`P_{orb}`$, and 100 pixels per decade in $`\dot{M}`$.
The most noteworthy features in Figure 3 include the distinct groups of systems located above and below the period gap. Note the substantial difference in $`\dot{M}`$ for systems above and below the period gap; for the latter systems only gravitational radiation losses drive mass transfer. The minimum orbital period ($`P_{min}`$65 min) is also clearly evident, as are systems that have evolved well past the minimum period back up to values of $`P_{orb}2`$ hr. It has been proposed that these latter systems may be related to the so-called TOADs (“Tremendous Outburst Amplitude Dwarf Novae”; see, e.g., Howell et al. 1995, HRP). In the systems above the gap, there is a central band of evolutionary tracks (blue and green) where a typical CV is most likely to be found at a particular point in time during its evolution. One also notices very short lived episodes (red and yellow) of high mass transfer rates. These occur for individual systems as the donor star first fills its Roche lobe and commences mass transfer, but before it can come into a quasi-steady state of mass transfer (see discussion in §2). The same type of behavior is seen (green structure) for systems that have come into contact for the first time below the period gap, i.e., with initially very low-mass donor stars. The two main tracks (purple) evident in the systems below the period gap are for He (lower) and CO (upper) white dwarfs, respectively. We also call attention to the small vertical (blue) feature at $`P_{orb}2`$ hr. This may be related to the statistically significant larger number of CVs with periods in the range of $`110120`$ minutes (first pointed out by Hameury et al. 1988b). Finally, we note that there are systems found within the period gap, though fewer per period interval than for systems below the gap. Systems found within the period gap are typically ones which had initial donor masses of $``$0.22-0.34 M and commenced Roche-lobe overflow at orbital periods in the range of $``$2-3 hr (see also Fig. 2 and its associated discussion).
Further, in regard to the $`\dot{M}P_{orb}`$ plane shown in Figure 3, we point out that for systems above the period gap, the width of the distribution in $`\dot{M}`$ at any fixed value of $`P_{orb}`$ is only about a factor of $``$2 (we define the “width” as containing $`80\%`$ of the systems). This is in contrast with the observed spread in $`\dot{M}`$ for CVs which is closer to an order of magnitude (see, e.g., Patterson 1984; Warner 1995). One cause of this spread may be the inherent uncertainty in translating observed parameters into accurate estimates of $`\dot{M}`$. Additionally, some of this discrepancy might be resolved by the inclusion of the effects of nova explosions in CVs which, on a quasi-regular basis, slightly increase (or perhaps even decrease) the orbital separation (e.g., by $`\delta a/a10^4`$) which is sufficient to change $`\dot{M}`$ appreciably for some interval of time (see, Shara et al. 1986; Schenker et al. 1998; Kolb et al. 2000). However, we note that Schenker et al. (1998) showed that, except for extreme model parameters, the occurrence of the nova explosions generally does not substantially affect the overall secular evolution of the CVs. Therefore the main results and conclusions presented in this work should be robust even without the inclusion of orbital perturbations due to nova explosions (we do, in fact, take into account the mass and angular momentum lost in such events).
The population of current-epoch CVs in the $`R_2P_{orb}`$ plane for our standard model is shown in Figure 4. The shape traced out in this figure represents a statistical ensemble of the type of evolutions graphed in Figs. 2c and 2f. The usual features of the “upper branch” of systems above the period gap, systems in the “lower branch” below the gap, and the minimum orbital period are all represented in this figure. Again, as in Fig. 3, we see that some systems are formed within the period gap. It is difficult from this image to judge quantitatively how many systems are in the gap, versus the density of points just below the gap. This is quantified later in this section (see Fig. 9). Note that an extrapolation of the “upper branch” to shorter orbital periods would undershoot the “lower branch”, on which the stars are close to thermal equilibrium. As discussed in §2, this undershooting actually (counterintuitively) results from the thermal bloating of the donor star when it has a higher mass loss rate which is driven by magnetic braking. In particular, see equation (7) where we show that $`R_2`$ scales as $`f^{0.65}`$, where $`f`$ is the bloating factor. The low-density features (yellow) just above the main tracks through the “upper branch” are systems that have just come into Roche-lobe contact for the first time and have not yet established a quasi-steady state of mass transfer. The blue-green “thumb” feature just below the main track of the “upper branch” near the top edge of the period gap, represents systems with He white dwarfs and donor stars of comparable mass that have just come into contact. Their mass transfer rates are higher than normal for these orbital periods; thus the bloating factors for these donor stars are significantly larger than for systems on the main track (cf. Fig. 2f).
Perhaps the most dramatic demonstration of the effects of thermal bloating of the donor star can be seen in Figure 5 which shows the population of current-epoch CVs in the $`M_2P_{orb}`$ plane for our standard model. The shape traced out in this figure represents a statistical ensemble of the type of evolutions graphed in Figs. 2b and 2e. All of the features that appear in the $`R_2P_{orb}`$ image (Fig. 4), also appear in this $`M_2P_{orb}`$ image, except in a more exaggerated form. This is a direct result of the simple scaling argument summarized in equation (6) in §2, which indicates that the bloating effect on the masses just above the period gap scales as $`M_2f^{1.95}`$. A casual inspection of Figure 5 shows that the masses of the donor stars in CVs with periods just above the period gap are fully $`40\%`$ lower than would be expected if their radius-mass relation followed that of main-sequence stars. It is this effect that we propose be used to discriminate between the currently held explanation for the period gap and alternate scenarios. We return to a quantitative discussion of this issue in the next section.
Lastly in regard to the color images of the $`R_2P_{orb}`$ and $`M_2P_{orb}`$ planes (Figs. 4 and 5), we comment on the relatively large spreads in $`R_2`$ and $`M_2`$ for systems above the period gap in contrast with those below the gap. As we showed in equations (6) and (7), for a fixed value of the bloating parameter $`f`$, both $`M_2`$ and $`R_2`$ are unique functions of the orbital period (which would imply narrow tracks). For systems well above the period gap, the Kelvin timescale, $`\tau _{KH}`$, is shorter than the mass-loss timescale, $`\tau _{\dot{M}}M/\dot{M}`$, but, as the orbital period decreases and approaches the period gap, the two timescales become more comparable. Thus, as discussed in §2, the donor star must become ever more bloated so as to establish a luminosity deficit, which in turn enables the donor to contract inside its ever shrinking Roche lobe. Additionally, the adiabatic stellar index is changing from positive to negative, and this tends to make the star expand even further as it loses mass (see also Beuermann et al. 1998). These two effects lead to the bloating behavior that is seen in Figs. 4 and 5. The actual amount of bloating depends upon the absolute values of the two constituent masses as well as on the thermal history of the donor; therefore, it is to be expected that $`f`$ may vary from one donor star to another. As a result, we not only see enhanced bloating as systems approach the period gap, but a relatively wider and wider spread in the values of $`M_2`$ and $`R_2`$ for these systems, especially for $`P_{orb}`$ in the range of 3-5 hr. By contrast, for systems just below the period gap, $`\tau _{\dot{M}}`$ increases abruptly - by about an order of magnitude - because the mass transfer is then driven only by gravitational radiation losses (at least according to our model), and therefore the donor stars can remain much closer to thermal equilibrium. This allows the systems right below the gap to establish a nearly main-sequence radius-mass relation (i.e., $`f1`$), thereby leading to a relatively narrow set of evolution tracks. However, as the secondary’s mass approaches the minimum main-sequence mass (before the orbital period minimum), $`\tau _{KH}`$ becomes very long (due to a sharp decrease in the secondary’s nuclear luminosity), thereby causing $`\tau _{KH}`$ and $`\tau _{\dot{M}}`$ to again become approximately equal. Thus, the width of the tracks broadens somewhat near the orbital period minimum. For systems beyond the orbital period minimum, the interiors become increasingly electron degenerate. This leads to a nearly unique mass-radius relationship $`R_2M_{2}^{}{}_{}{}^{1/3}`$ which, in turn, leads to entirely different $`R_2(P_{orb})`$ and $`M_2(P_{orb})`$ relations than are given by equations (6) and (7). Nonetheless, they are unique relations (easily derivable from eqs. and ) which also lead to a very narrow set of tracks in Figs. 4 and 5.
The distribution of expected mass ratios, $`q`$, in CVs at the current epoch is shown as a function of orbital period in Figure 6. At any given orbital period the range of $`q`$ values is considerably broader than the distribution of values of $`R_2`$ or $`M_2`$ as can be seen by comparison with Figs. 4 and 5. The reason for this is straightforward. Equations (6) and (7) indicate that, as long as the bloating factor $`f`$ depends largely on the orbital period of a CV, then both the radius and mass are nearly unique functions of the orbital period. Thus, the much broader distribution of $`q`$ in Fig. 6 is due largely to the substantial range of masses that the white dwarf may have, which is much less constrained by the orbital period than is $`M_2`$. For both the systems above and below the period gap, the upper set of tracks corresponds to He white dwarfs, while the lower tracks are for CO white dwarfs. The period gap is especially conspicuous in this figure, especially for systems with CO white dwarfs. Note that some of the mass ratios extend up to values of unity and, in some cases, above unity. The stability of mass transfer in these systems was discussed in §4.2 \[see equation (11)\].
The evolution of our model population of CVs in the $`T_{eff}P_{orb}`$ and luminosity$`P_{orb}`$ planes is shown in Figure 7. The left panel displays the effective temperature, $`T_{eff}`$ of the donor star, while the right panel has a superposition of the optical (bolometric) luminosity, $`L_{opt}`$ and nuclear luminosity $`L_{nuc}`$. Where the two sets of luminosity tracks cross (e.g., near $`P_{orb}`$ = 1 hr) or overlap (e.g., to a minor extent between 3 and 6 hrs), the default is to display $`L_{opt}`$. With regard to the $`T_{eff}`$ curves, we first note that the absolute temperature scale for our main sequence stars (based on our bipolytrope code) is somewhat shifted from that produced with more sophisticated codes, e.g., our bipolytrope main-sequence models are $``$200 K higher than the DNC models over the mass range of $`0.850.1M_{}`$. However, aside from this small quantitative difference we are confident that the overall qualitative trends and shapes of these tracks are highly indicative of the behavior and properties of the donor through its evolutionary history. Note that for values of $`P_{orb}`$ below the gap as well as above $``$5 hr, the $`T_{eff}`$ tracks are quite narrow, in analogy with the tracks in the $`M_2P_{orb}`$ and $`R_2P_{orb}`$ planes, since the donor stars are typically quite close to thermal equilibrium. By contrast, within the period range of 3-5 hr, $`T_{eff}`$ of the donors is systematically lowered by up to 250 K compared with $`T_{eff}`$ of main sequence stars at the same $`P_{orb}`$. This lower temperature amounts to a change to a later spectral type (at a given $`P_{orb}`$) of $`24`$ in decimal subclass. Additionally, we can see from Fig. 7 that, over this same period range, the use of temperature (or spectral type) to determine the mass of the secondary star would require very precise measurements, since the expected $`T_{eff}P_{orb}`$ distribution is relatively flat. We draw two conclusions from this figure: (1) an observationally produced $`T_{eff}P_{orb}`$ or spectral type-$`P_{orb}`$ relation for CVs should indeed yield a fairly simple shape (especially when smoothed out by uncertainties in the measurements), and (ii) the use of $`T_{eff}`$ or spectral type in the 3-5 hr period range will not yield reliable indications of the mass of the donor star.
Recently, Smith & Dhillon (1988)<sup>4</sup><sup>4</sup>4The sample used in Smith & Dhillon consisted of what are believed to be 55 reliable spectral types and 14 reliable secondary star masses. All systems in their sample have $`P_{orb}>`$ 90 min and $`V_{min}`$ brighter than $``$17th magnitude, thus Smith & Dhillon’s conclusions about finding no evidence for post-period minimum systems or very low-mass brown dwarf-like secondaries cannot be drawn from the sample they used. published results which took a critical look at the relation between orbital period, spectral type, and secondary mass based on the best observational data available in the literature. They presented a relatively smooth spectral type-orbital period relation but concluded that one cannot reliably estimate $`M_2`$ in any given CV based solely on its spectral type. Our theoretical results are quite consistent with their conclusion.
The image in the right panel of Figure 7 displays a superposition of the evolutionary tracks for $`L_{opt}`$ and $`L_{nuc}`$ as functions of $`P_{orb}`$. For systems above the gap, the highest luminosity track corresponds to $`L_{opt}`$, while the two prominent lower (green) tracks are for $`L_{nuc}`$ and are related to the corresponding features in the $`M_2P_{orb}`$ image. These two lower luminosity tracks are for systems with CO (upper) and He (lower) white dwarf accretors (see the discussion of Fig. 5). The large luminosity deficit in the $`P_{orb}`$ range of 3-5 hr, already discussed in §2, shows up quite dramatically in this Fig. 7. The group of systems with the highest luminosity deficit (with He white dwarf accretors) has the largest values of $`\dot{M}`$ and the donors are the most out of thermal equilibrium (largest bloating factor). In spite of the relatively low values of $`L_{nuc}`$ in this period range, the bolometric luminosity $`L_{opt}`$ is depressed only modestly (e.g., by factors of $``$ 2) over main-sequence stars at the same orbital period. For systems below the period gap, both luminosities fall off dramatically, especially for donor masses below $`0.050.08M_{}`$, where the donors are already below the hydrogen-burning main sequence, and are cooling toward their ultimate degenerate state. The higher track for all points below the period gap corresponds to $`L_{opt}`$, the lower one to $`L_{nuc}`$. While it is formally true that the $`L_{opt}`$ and $`L_{nuc}`$ tracks “cross” at $`10^4L_{}`$, the two luminosities are never equal in this part of the diagram; they reach the crossing point at very different times.
Finally, with regard to the color image representations of CV populations in parameter space, we note that in Figs. 3 through 7, the color represents the logarithm of the numbers of systems expected at the current epoch. As can be seen in any of these figures (but especially Figs. 3 and 6), the number of systems below the period gap outweighs the number above the period gap by a large margin (by about 100:1; see the more quantitative discussion below; see also de Kool 1992, Kolb 1993). However, due to observational selection effects, the systems with the shorter orbital periods, lower values of $`\dot{M}`$, and generally longer intervals between dwarf-nova outbursts, are more difficult to discover. The exact factors that go into the observational selection effects are complex, especially since some CVs are discovered via their dwarf-nova outbursts, others (e.g., longer period CVs) by their blue colors or flickering behavior, and still others by their nova outbursts. Some of these issues are discussed in RJW and Kolb (1993). For purposes of the present work we will indicate only qualitatively how the numbers of observationally known CVs might be expected to be distributed by orbital period. We adopt two crude detectability factors which scale simply as $`\dot{M}^{3/2}`$ and as $`\dot{M}`$. The first of these is appropriate to steady-state accretion luminosities, that are proportional to $`\dot{M}`$ in the optical bandpass, which give rise to a bolometrically flux-limited detectability proportional to $`\dot{M}^{3/2}`$ (analogous to the 3/2 slope of a log(N)-log(S) curve for isotropically distributed sources). The other scaling, appropriate to the case where the flux in the optical bandpass is proportional to $`\dot{M}^{2/3}`$ (see Lynden-Bell & Pringle 1974; RJW; Webbink et al. 1987), leads to a flux-limited detectability proportional to $`\dot{M}`$. Other factors leading to the discovery of CVs, beyond the simple consideration of flux limited samples, in particular the detection of dwarf nova outbursts, would undoubtedly substantially modify the rudimentary dependences on $`\dot{M}`$ that we use here for purposes of illustration. In Figure 8 we redisplay Fig. 3, but this time rescaled by a factor of $`\dot{M}`$. It is clear from a casual inspection of Fig. 8 that the number of “detectable” systems above the period gap is now at least as great as for those below the gap. The actual quantitative values for this simple scaling are presented below. We again caution, however, that either an $`\dot{M}^{3/2}`$ or $`\dot{M}`$ scaling is oversimplified.
The color images of parameter space shown in Figures 3 – 7 can be displayed in a somewhat more quantitative fashion by projecting the numbers of systems onto the various axes and plotting the results as simple histograms. For example, the data used to produce any of the images, can be projected onto the $`P_{orb}`$ axis to yield the orbital period distribution. The results are shown in Figure 9. The solid histogram in Fig. 9a is the distribution of CVs at the current epoch in the entire Galaxy for our standard model (see Table 1). The stellar birthrate function and IMF in our Standard Model \[eq. (8)\] are normalized in such a way that there are $`0.6`$ stars born in the Galaxy per year with a mass $`>0.8M_{}`$, just above the threshold for producing a remnant white dwarf by the current epoch. Thus, the “absolute values” of the numbers plotted in Fig. 9 can be appropriately scaled up or down for either lower or higher assumed birthrates.
If the numbers of CV systems are scaled by the types of “observability factors” discussed above, before the histogram is produced, the results are the dashed and dotted histograms superposed in Figure 9a. As discussed above, in conjunction with Fig. 8, this qualitatively helps to explain the relative numbers of CVs observed above the period gap compared with the number observed below (see especially the dotted histogram). Inspection of the compilation of CVs with known orbital periods given in Warner (1995) reveals that our histogram shown in Fig. 9a with the $`\dot{M}`$ scaling provides qualitative agreement with current observational results, especially considering the many observational selection effects that exist (e.g., magnitude-limited color surveys, large-amplitude but infrequent outbursts compared with semi-periodic lower amplitude outbursts, discovery in X-ray surveys, etc.). The distributions of orbital period shown in Figure 9b are for systems that have not yet evolved to the minimum orbital period (solid curve), and systems that have evolved beyond the period minimum (dashed curve) - no scaling in $`\dot{M}`$ has been applied here.
The distributions of white dwarf masses and donor masses at the current epoch are shown in Figure 10 (left and right panels respectively) for four different ranges of orbital period. The He and CO white dwarfs are easy to distinguish by mass. Note that for systems with $`P_{orb}>4`$ hr, which typically have donor stars with masses $`>0.4M_{}`$, there are few He white dwarfs, since the mass transfer would tend to be unstable. The distributions of donor star masses show a steady trend toward higher masses at the longer periods, as expected. This results qualitatively from the fact that the larger orbital periods require less dense, and therefore usually more massive, stars. Note that the distributions shown in this figure are not produced with sufficient resolution in $`P_{orb}`$ to allow one to make quantitative predictions as to what mass donors are needed to validate the basic paradigm for the period gap. Such information may be found, however, in Figs. 5 and 12, and Table 2.
The distribution of mass ratios $`q`$ ($`M_2`$/$`M_{WD}`$) is shown in Figure 11 for two different orbital period ranges. Attempts to determine $`q`$ observationally can be made from, for example, superhump period analysis or spectroscopic analysis. The observational distribution for $`q`$ in short period ($`P_{orb}<2`$ hr) CVs has recently been compiled (Mennickent et al. 1999), and is seen to show an approximate Gaussian distribution with $`<q>`$ = 0.14. However, observational selection effects allow few CVs with small $`q`$ to be discovered due to their intrinsic faintness. Our results (Fig. 11 - top panel) show that the actual distribution should not drop off at $`q`$ values lower then 0.14, but rather should peak at values of $`q`$ = 0.05-0.1, with an overall distribution that is clearly non-Gaussian. Discovery and observation of additional faint (short period) CVs are needed in order to confirm this theoretical prediction.
## 6 Test of the Basic Paradigm
The rapid rate of mass loss for donor stars in CVs just above the period gap should lead to significant thermal bloating of the donor. Thus, in the conventional paradigm for the formation of the period gap, this mass loss rate is abruptly decreased at orbital periods near 3 hr and the donor star shrinks inside its Roche lobe (see §2), leading to the cessation of mass transfer. Specific choices of the parameters utilized in any such evolutionary model change the bloating factor quantitatively, but do not change the overall evolution qualitatively. To demonstrate this, we show in Figure 12 CVs at the current epoch in the $`M_2P_{orb}`$ plane for four different sets of model parameters (see Table 1). Panel (A) is for our Standard Model, while the other panels are for models where (B) the proportionality constant in the magnetic braking formula was reduced by a factor of 2 ($`C_{MB}=1/2`$), (C) the specific angular momentum carried away by mass lost from the system in nova explosions is twice that of the white dwarf ($`\alpha =2\alpha _{WD}`$), and (D) all mass transferred to the white dwarf is ultimately retained by the white dwarf (i.e., $`\beta =1`$; in this somewhat artificial model, white dwarfs are allowed to exceed the Chandrasekhar Limit).
We see from a study of Figure 12 that the effects of thermal bloating on the mass of the donor stars in CVs for orbital periods just above the gap are qualitatively similar for all four models. The actual factors by which the masses are lower than would be inferred by making the assumption that the donor has a main-sequence mass-radius relation range from $`2550\%`$; the exact range depends on which model parameters are chosen and whether one includes the CVs with He white dwarfs where the mass transfer can be only marginally stable. To quantify the effect of thermal bloating on mass determinations, we have carried out weighted least squares fits of polynomials to each of the “upper branches” shown in Fig. 12. The results are given in Table 2 which also includes the evaluation of the polynomial fit at $`P_{orb}=3`$ hr. As we can see from Table 2, the effect of thermal bloating on the inferred donor mass of CVs is quite significant, and potentially testable for any of these models.
We note that, in general, the spread in values of $`M_2`$, at a given $`P_{orb}`$, around the best-fit $`M_2(P_{orb})`$ curve is substantially smaller than the mean deviation from an $`M_2(P_{orb})`$ curve based on the assumption of a main-sequence mass-radius relation, especially in the crucial period range of 3-5 hr. A large part of this spread is due to the different values of the mass of the accretor, $`M_1`$, with lower values of $`M_2`$ corresponding to the lower values of $`M_1`$ (see also the discussion in §5). However, we do not attempt here to produce fits of the more general form $`M_2(P_{orb},M_1)`$. Such fits are not straightforward to construct since, among other things, there is the added complication of the existence of a minimum value for $`M_1`$ at any $`P_{orb}`$ (due to issues of mass-transfer stability; see §4.2). In any case, the main effect to be confirmed observationally concerns the substantially reduced values of $`M_2`$ just above the period gap (3-5 hr), compared to what would be expected if the donor stars followed a main-sequence mass-radius relation. If sufficient numbers of high quality mass determinations of the secondary stars can be made, and this basic effect is confirmed, then a secondary goal would be to look for a weak, but positive correlation between $`M_2`$ and $`M_1`$.
In Figure 13 we plot the polynomial fits that we made to the upper branches in the $`M_2P_{orb}`$ plane for the four different models. For comparison we show the $`M_2P_{orb}`$ relation that would be obtained if the donor star followed a main-sequence radius-mass relation (the one derived from our bipolytrope code). This set of curves shows quantitatively how mass determinations based solely on $`P_{orb}`$ are affected by the thermal bloating effect. Note how the effect should go from a maximum at $`3`$ hr to quite small at $`P_{orb}5.5`$ hr.
Finally, we point out that if, in fact, the period gap is in any way related to a relaxation from thermal bloating, then the inferred effect on mass determinations based on the orbital period must be approximately in the range of $`2550\%`$. To demonstrate this, we note that in the basic paradigm for producing the period gap, the system masses do not change from the upper boundary of the gap ($`P_{upper}`$) to the lower boundary ($`P_{lower}`$), while the radius shrinks from its bloated state, characterized by a bloating factor, $`f`$, to nearly its main-sequence radius at the lower edge of the gap. A simple application of Kepler’s 3rd law for the case of a Roche-lobe filling star (which is true at both the upper and lower edges of the gap) shows that
$`f=\left({\displaystyle \frac{P_{upper}}{P_{lower}}}\right)^{2/3},`$ (13)
where $`f`$ must range from $`1.21.3`$, depending on whether the period gap is taken to be 3/4 of an hour in width or 1 hour, respectively (we have assumed that the gap is centered at 2.5 hr). From equation (6) we see that this value of $`f`$ should reduce the inferred mass, at the top edge of the period gap, by amounts ranging from $`30\%40\%`$, in basic agreement with our more detailed population synthesis study. (For a related discussion see Beuermann et al. 1998.)
This type of discrepancy between the mass inferred for a secondary star, based on the CV orbital period and the assumption that its radius is that of a main-sequence star has probably already led to a number of incorrect mass determinations reported in the literature, particularly for systems with $`P_{orb}`$ between 3-5.5 hr. For example, at an orbital period of 3.2 hours, the mass assigned to a CV secondary would be 0.35 M while our calculations show that it would actually be only 0.26 $`\pm 0.02M_{}`$, although bloated in size. For a known or inferred mass ratio of say $`q=0.4`$, we would then calculate a white dwarf mass of 0.89 M, when the true white dwarf mass is only 0.65 M. Thus, ignoring the bloating effect in the secondary stars in CVs with orbital periods of 3-5 hr, can lead to a significant overestimation of both component masses.
It is interesting to note here that the secondary stars which are farthest from thermal equilibrium are those in systems with orbital periods just above the period gap (see Figs 4, 5 & 13). Observationally, this orbital period region (3-4 hr) essentially contains only high mass transfer rate, novalike (NL) types of CVs. The inferred high mass transfer rates for these systems would then be expected, on theoretical grounds, to lead to a large bloating of the secondary stars, and hence lower masses than might otherwise be anticipated. Precise observational determinations of the secondary star masses in NLs would allow a confirmation of this basic effect which is, in fact, required if the period gap is to be explained by the interrupted magnetic braking scenario.
Figure 13 provides our theoretical predictions for the most likely mass of the secondary star at any given orbital period (see Table 2). Observational determinations accurate to a few percent would be needed in order to differentiate between the four models presented; but, accuracies of only $`10\%`$ will allow a test of the bloating model in general, and the predicted deviation of the donor star from the main sequence. This is a challenging observational project, however, since the systems with orbital periods in the 3-5 hr range are ones in which the secondary star is rarely directly observed. IR spectral studies (e.g., Howell, et al. 2000; Mason, et al. 2000; Dhillon, et al. 1997) have looked in detail for the secondary star in a number of CVs with only marginally successful results. In these CVs, spectral identification of absorption features due to the secondary star is difficult since the lines are rotationally broadened and filled in by radiation from the accretion disk. For the critical 3-5 hr period range, a signal-to-noise ratio of $`>`$100 in the continuum will be needed to allow the atomic and molecular features of the secondary to be observed against the high background accretion-disk dominated continuum. We therefore advocate high signal-to-noise, orbital phase-resolved, near- and mid-IR spectroscopic observations with large ground-based telescopes (e.g., Gemini, Keck), and eventually with SIRTF, of sources such as the brightest NLs and other CVs which have $`P_{orb}`$ = 3-5 hr.
## 7 Summary and Conclusions
In this paper we briefly reviewed our current understanding of the secular evolution of CVs through their mass transfer phase, including the currently accepted model for the 2-3 hr “period gap” in the orbital period distribution. The results of evolution calculations for a representative sample of individual systems are presented, both as functions of time and of orbital period. A population synthesis code, that starts with some $`3\times 10^6`$ primordial binaries, was then used to generate $`2\times 10^4`$ systems which evolve successfully through the CV phase of mass transfer. This allows for a more complete exploration of parameter space. The results are displayed as probability densities in the $`\dot{M}P_{orb}`$, $`M_2P_{orb}`$, $`R_2P_{orb}`$ $`qP_{orb}`$, and $`T_{eff}/L_2P_{orb}`$ planes, for CVs at the current epoch. This method of displaying the results can lead to considerable insight into the relationships among the various system parameters. We find that for CVs with orbital periods above 5.5 hr and below the period gap (but above the period minimum) the secondary stars closely follow the main-sequence R-M relation (cf. Beuermann et al. 1998). However, for those with $`P_{orb}`$ between 3-5.5 hr, the effect of bloating causes them to deviate substantially from this same relation.
Among our more interesting results, we have shown that the donor star masses in CVs with orbital periods just above the period gap should be as much as $`3050\%`$ lower than would be inferred on the assumption that the donor stars obey a main-sequence radius-mass relation. This conclusion is only valid if the basic underlying cause of the period gap is thermal bloating of the donor star for systems above the period gap (see §§1-6). On the basis of our results, we have proposed a direct observational test of, in particular, the basic paradigm of the period gap and, more generally, our overall understanding of the evolution of CVs. This test involves the challenging, but realistic, task of making relatively accurate (e.g., $`10\%`$) determinations of the secondary masses in about a half dozen CVs in the period range of 3-4 hr. If the masses are consistent with the assumption of a main-sequence radius-mass relation for the donor stars, then the currently accepted explanation of the period gap cannot be correct, and the very existence of the gap would pose a major conundrum. If, on the other hand, the masses are mostly consistent with the lower values predicted in this work, then a substantial part of our basic understanding of the secular evolution of CVs will be validated.
Previously, much observational attention in CV studies has been focused on determinations of the white dwarf masses. While this is clearly of great interest, we hope with this work to stimulate more interest in the important issue of determining the secondary masses.
This research was supported in part by NASA under ATP grants GSFC-070 and NAG5-8500 (to S.B.H.), and NAG5-7479 and NAG5-4057 (to S.A.R.). L.A.N. acknowledges the financial support of NSERC (Canada) and thanks CITA and the University of Toronto for a Reinhardt Fellowship and for their hospitality. We thank M. Politano for a number of useful discussions relating to this work. We are grateful to an anonymous referee who made numerous helpful and insightful comments that led to significant improvements in the paper. We also thank D. MacCannell and G. Esquerdo for their technical assistance.
Fig. 1 – Evolution with time of the mass transfer rate, $`\dot{M}`$, and orbital period, $`P_{orb}`$, for several model cataclysmic variable systems. Left panel - the evolution of a single CV with initial masses ($`M_2=0.9M_{}`$; $`M_{WD}=1.1M_{}`$). This system first comes into Roche-lobe contact at $`P_{orb}=6`$ hr, evolves through the period gap, to the minimum in $`P_{orb}`$, and back up to longer periods by $`10^{10}`$ yr. Right panel - the evolutions of a selection of four other illustrative initial binary constituent masses, $`M_2`$,$`M_{WD}`$ = 0.2,0.4 (solid), 0.35,0.35 (dotted), 0.3,0.6 (dashed), and 0.65,0.7 (long dashed), all in units of $`M_{}`$.
Fig. 2 – Evolution with orbital period, $`P_{orb}`$, of the mass transfer rate, $`\dot{M}`$, secondary mass, $`M_2`$, and secondary radius, $`R_2`$, for several illustrative model cataclysmic variable systems. The initial masses for the systems whose evolutions are displayed in the left and right sets of panels are the same as described in Fig. 1.
Fig. 3 – Computed population of cataclysmic variables at the current epoch in the $`\dot{M}P_{orb}`$ plane for our Standard Model (see Table 1). Here $`\dot{M}`$ is the mass transfer rate, and $`P_{orb}`$ is the orbital period. The color represents the logarithm of the number of systems in a particular $`\dot{M}P_{orb}`$ cell, of which there are 100 per hour interval in $`P_{orb}`$ and 100 per decade in $`\dot{M}`$. The color scale is given on the right side of the figure. We note that the scattered, isolated (red) points in the image below the main tracks are minor numerical artifacts of the evolution code that occasionally appear when the Roche lobe makes initial contact with the atmosphere of the donor star. One of these dots corresponds to only $`0.1`$ CVs in the entire Galaxy at the current epoch, and so is of no significance.
Figs. 4 – Computed population of cataclysmic variables at the current epoch in the $`R_2P_{orb}`$ plane for our Standard Model (see Table 1). Here $`R_2`$ is the radius of the donor star. The color represents the logarithm of the number of systems in a particular $`R_2P_{orb}`$ cell, of which there are 100 per 0.1 $`R_{}`$ and 100 per hour interval in $`P_{orb}`$. The color scale is given on the right side of the figure.
Figs. 5 – Computed population of cataclysmic variables at the current epoch in the $`M_2P_{orb}`$ plane for our Standard Model (see Table 1). Here $`M_2`$ is the mass of the donor star. The color represents the logarithm of the number of systems in a particular $`M_2P_{orb}`$ cell, of which there are 100 per 0.1 $`M_{}`$ and 100 per hour interval in $`P_{orb}`$. The color scale is given on the right side of the figure.
Fig. 6 – Computed population of cataclysmic variables at the current epoch in the $`qP_{orb}`$ plane for our Standard Model (see Table 1); $`qM_2/M_{WD}`$. The color represents the logarithm of the number of systems in a particular $`qP_{orb}`$ cell, of which there are 100 per $`\mathrm{\Delta }q=0.1`$ and 100 per hour interval in $`P_{orb}`$. The color scale is given on the right side of the figure.
Fig. 7 – Computed population of cataclysmic variables at the current epoch in the $`T_{eff}`$$`P_{orb}`$ plane (left panel), and the $`LuminosityP_{orb}`$ plane (right panel) for our Standard Model (see Table 1). We show both the stellar luminosity (top curve) and the core nuclear luminosity (lower distributions). The color represents the logarithm of the number of systems in a particular $`LP_{orb}`$ or $`T_{eff}P_{orb}`$ cell of which there are 100 per decade in $`L`$, 100 per 500K in $`T_{eff}`$, and 100 per hour interval in $`P_{orb}`$. The color scale for both plots is given on the right.
Fig. 8 – Same as Figure 3, except that the population has been scaled by $`\dot{M}^1`$ to crudely take into account observational selection effects.
Fig. 9 – Computed orbital period distributions for cataclysmic variables at the current epoch. Left panel - solid curve is the distribution for all systems that appear in Figure 3; the dashed curve was produced by scaling the contributions of each system evolved by $`\dot{M}^{3/2}`$ while the dotted curve is for an $`\dot{M}^1`$ scaling (see text). The $`\dot{M}^{3/2}`$\- and $`\dot{M}^1`$-scaled curves have been shifted vertically by arbitrary amounts for ease in comparison. Right panel - solid curve is for all systems in Fig. 3 which have not yet reached orbital period minimum; dashed curve is for systems that have evolved past the orbital period minimum.
Fig. 10 – Computed distributions of the secondary (right panels) and white dwarf masses (left panels) in cataclysmic variables at the current epoch. The mass distributions are ordered according to the range of orbital period. The dotted histogram (upper right) is for post-period minimum CVs and has been arbitrarily divided by 1.5 for presentation purposes.
Fig. 11 – Computed distribution of mass ratios in cataclysmic variables at the current epoch. The top panel is for systems with orbital periods in the range of 1-3 hr (which includes all post period-gap systems), while the bottom panel is for systems above the period gap.
Fig. 12 – Same as Figure (5), except that in addition to the Standard Model (A), the results for three other models are shown (see Tables 1 and 2): (B) reduced magnetic braking constant; (C) specific angular momentum lost with the ejected matter is twice that of the white dwarf; and (D) conservative mass transfer and retention by the white dwarf.
Fig. 13 – Secondary (donor) mass, $`M_2`$ as a function of orbital period. The solid curve is based on the assumption that the donor star fills its Roche lobe and has a radius-mass relation appropriate to stars on the main sequence (i.e., eq. ) The main-sequence models were generated with the same bipolytrope code that was used to carry out the binary stellar evolution calculations and are discussed in the text. The dashed curves are polynomial fits to the $`M_2P_{orb}`$ relations derived from the population synthesis study shown in Figure 12. The labels, A through D, correspond to the four different panels in Figure 12.
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# 1 Introduction
## 1 Introduction
During the last two years, it has become clear that the proper way to formulate chiral fermions on the lattice is to impose the exact chiral symmetry on the lattice, namely, the Ginsparg-Wilson relation
$`D\gamma _5+\gamma _5D=2aDR\gamma _5D,`$ (1)
where $`D`$ is the lattice Dirac operator, $`a`$ is the lattice spacing, and $`R`$ is a positive definite Hermitian operator which commutes with $`\gamma _5`$. Equation (1) should be regarded as a generalized chiral symmetry which contains the usual chiral symmetry in the continuum limit ( $`a0`$ ). However, it should be noted that this exact chiral symmetry ( $`R=1`$ ) had been existing in the Overlap formalism , even before the GW relation was rediscovered. Therefore, unless one can explicitly construct a GW Dirac operator $`D`$ without using the Overlap, and such a $`D`$ satisfies all physical requirements; otherwise it is unlikely that the GW relation would turn out to be more fundamental than the Overlap.
Recently, Fujikawa proposed a generalization of the Ginsparg-Wilson relation as
$`\gamma _5(\gamma _5D)+(\gamma _5D)\gamma _5=2a^{2k+1}(\gamma _5D)^{2k+2},k=0,1,2,\mathrm{}`$ (2)
Multiplying both sides of (2) by $`\gamma _5`$, we obtain
$`\gamma _5D+D\gamma _5=2aD(a\gamma _5D)^{2k}\gamma _5D,k=0,1,2,\mathrm{}`$ (3)
which is equivalent to the original GW relation (1) with
$`R=(a\gamma _5D)^{2k},k=0,1,2,\mathrm{}`$ (4)
It can be shown that $`R`$ is Hermitian and commutes with $`\gamma _5`$ for any $`D`$ which is $`\gamma _5`$-Hermitian and satisfies (3) \[ the proof is below Eq. (16) \]. The motivation of considering $`k>0`$ in (4) is to improve the chiral symmetry at small lattice spacings. Further, Fujikawa has constructed a sequence ( $`k>0`$ ) of these GW Dirac operators based on the Neuberger-Dirac operator ( $`k=0`$ ) . However, the price one has to pay for the improved chiral symmetry is a less localized $`D`$, since in the limit $`k\mathrm{}`$, $`D`$ must tend to a chirally symmetric and nonlocal $`D_c`$, as a consequence of the Nielson-Ninomiya no-go theorem . Therefore, it is not clear whether one may have any advantages in practice by considering $`k>0`$ in (4). Nevertheless, from a theoretical viewpoint, it is interesting to see how one can construct a sequence ( $`k=1,2,\mathrm{}`$ ) of topologically proper $`D`$ satisfying the GW relation (3), in addition to the Neuberger-Dirac operator ( $`k=0`$ ).
In this paper, we examine several aspects of Fujikawa’s proposal. In section 2, we derive the analytical properties of the GW Dirac operator satisfying (1) with $`R=(a\gamma _5D)^{2k}`$. Then, in section 3, we analyze the construction of higher-order Overlap Dirac operators, and derive their general properties. In section 4, we compare the chiral properties of the higher-order ( $`k=1,2`$ ) Overlap Dirac operators to those of the Neuberger-Dirac operator, by computing the fermion propagator, the axial anomaly and the fermion determinant in two-dimensional background $`U(1)`$ gauge fields. Finally, we discuss and conclude in section 5.
## 2 General analytical properties
In this section, we begin with general considerations of the GW relation, and then derive the analytical properties of the GW Dirac operators satisfying (1) with $`R=(a\gamma _5D)^{2k}`$.
In general, one can assume that the lattice Dirac operator $`D`$ satisfies the Ginsparg-Wilson relation in the form
$`D\gamma _5f(D)+g(D)\gamma _5D=0,`$ (5)
where $`f`$ and $`g`$ are any analytic functions. Then the fermionic action $`𝒜_f=\overline{\psi }D\psi `$ is invariant under the global chiral transformation
$`\psi \mathrm{exp}[\theta \gamma _5f(D)]\psi `$ (6)
$`\overline{\psi }\overline{\psi }\mathrm{exp}[\theta g(D)\gamma _5]`$ (7)
where $`\theta `$ is a global parameter.
If we set $`f(D)=\text{1I}aRD`$ and $`g(D)=\text{1I}aDR`$, then (5) becomes
$`D\gamma _5+\gamma _5D=aD(R\gamma _5+\gamma _5R)D`$ (8)
where $`R`$ is any operator.
Since the massless Dirac operator in continuum is chirally symmetric ( $`𝒟\gamma _5+\gamma _5𝒟=0`$ ) and antihermitian ( $`𝒟^{}=𝒟`$ ), so it is $`\gamma _5`$-Hermitian ( $`𝒟^{}=\gamma _5𝒟\gamma _5`$ ). Thus, we require that the lattice Dirac operator $`D`$ also preserves this symmetry at any lattice spacing, i.e.,
$`D^{}=\gamma _5D\gamma _5.`$ (9)
Then multiplying (8) by $`\gamma _5`$ and using (9), we obtain
$`D^{}+D=aD^{}(R+\gamma _5R\gamma _5)D=aD(R+\gamma _5R\gamma _5)D^{}.`$ (10)
Evidently, only the part of $`R`$ which commutes with $`\gamma _5`$ can enter (10). Recall that any operator $`R`$ can be decomposed into two parts as
$`R={\displaystyle \frac{1}{2}}(R+\gamma _5R\gamma _5)+{\displaystyle \frac{1}{2}}(R\gamma _5R\gamma _5)`$
where the first ( second ) term on r.h.s. commutes ( anticommutes ) with $`\gamma _5`$. Therefore, without loss, one can assume that $`R`$ commutes with $`\gamma _5`$. Thus, (10) becomes
$`D^{}+D=2aD^{}RD`$ (11)
Taking the adjoint of (11), we immediately obtain
$`R=R^{}.`$ (12)
Then (8) becomes the usual GW relation
$`D\gamma _5+\gamma _5D=2aDR\gamma _5D,`$ (13)
where $`R`$ is any Hermitian operator which commutes with $`\gamma _5`$.
Fujikawa’s proposal is equivalent to setting
$`R=(a\gamma _5D)^{2k},k=0,1,2,\mathrm{}`$ (14)
in the usual GW relation (13). Then (13) becomes
$`D\gamma _5+\gamma _5D=2aD(a\gamma _5D)^{2k}\gamma _5D.`$ (15)
It is obvious that $`R`$ is Hermitian since $`D`$ is $`\gamma _5`$-Hermitian. Note that
$`\gamma _5(a\gamma _5D)^2=(a\gamma _5D)^2\gamma _5`$ (16)
since
$`\gamma _5(a\gamma _5D)^2`$
$`=`$ $`a^2\gamma _5(\gamma _5D+D\gamma _5)(\gamma _5D)a^2(\gamma _5D)\gamma _5(\gamma _5D+D\gamma _5)+a^2(\gamma _5D)\gamma _5(D\gamma _5)`$
$`=`$ $`2(a\gamma _5D)^{2k+2}2(a\gamma _5D)^{2k+2}+(a\gamma _5D)^2\gamma _5`$
$`=`$ $`(a\gamma _5D)^2\gamma _5`$
where (15) has been used in the second equality. Then it follows that $`R=(a\gamma _5D)^{2k}`$ commutes with $`\gamma _5`$.
Further, Eq. (16) gives
$`\gamma _5(\gamma _5D)^2\gamma _5=(\gamma _5D)^2`$ (17)
which yields
$`DD^{}=D^{}D,\text{ }D\text{ is normal}`$ (18)
where (9) has been used. Since $`D`$ is normal, $`D`$ and $`D^{}`$ have common eigenfunctions and their eigenvalues are either real or come in complex conjugate pairs,
$`D\varphi _s=\lambda _s\varphi _s,`$ (19)
$`D^{}\varphi _s=\lambda ^{}\varphi _s,`$ (20)
and the eigenfunctions $`\{\varphi _s\}`$ form a complete orthonormal set. Then the general analytical properties of the eigenmodes as derived in Ref. ( Eqs. (37), (38) and (41) in Ref. ) for $`R=1/2`$ also hold for the $`R`$ in (14). For the sake of completeness, we outline the derivations as follows. Writing $`R`$ as
$`R=(a\gamma _5D)^{2k}=a^{2k}(\gamma _5D\gamma _5D)^k=a^{2k}(D^{}D)^k,`$ (21)
and applying Eq. (11) to $`\varphi _s`$, we obtain the eigenvalue equation
$`\lambda _s+\lambda _s^{}=2a^{2k+1}(\lambda _s^{})^{k+1}(\lambda _s)^{k+1}.`$ (22)
Multiplying both sides of (22) by $`(\lambda _s)^k(\lambda _s^{})^k`$, we get
$`(\lambda _s)^{k+1}(\lambda _s^{})^k+(\lambda _s)^k(\lambda _s^{})^{k+1}=2a^{2k+1}(\lambda _s^{})^{2k+1}(\lambda _s)^{2k+1}`$ (23)
which can be rewritten as
$`\left|(a\lambda _s)^{k+1}(a\lambda _s^{})^k{\displaystyle \frac{1}{2}}\right|={\displaystyle \frac{1}{2}}`$ (24)
Thus the eigenvalues in the form $`z=(a\lambda _s)^{k+1}(a\lambda _s^{})^k`$ fall on the circle centered at $`1/2`$ with radius $`1/2`$. The real eigenvalues ( if any ) of $`D`$ are at $`\lambda =0`$ ( $`z=0`$ ) and $`\lambda =a^1`$ ( $`z=1`$ ). Writing $`a\lambda _s=r\mathrm{exp}(i\theta )`$, we have $`z=r^{2k+1}\mathrm{exp}(i\theta )`$. Thus, for $`k=0`$, the eigenvalues of $`D`$ fall on the circle centered at $`1/2`$ with radius $`1/2`$, while for $`k>0`$, on the deformed circle stretched symmetrically in $`\pm \widehat{y}`$ directions with fixed points at zero and $`a^1`$, bounded inside the region $`0xa^1`$.
From (9) and (20), we obtain
$`D\gamma _5\varphi _s=\lambda _s^{}\gamma _5\varphi _s.`$ (25)
Multiplying both sides of (25) by $`\varphi _s^{}`$ and using the adjoint of (20), we get
$`\lambda _s\varphi _s^{}\gamma _5\varphi _s=\lambda _s^{}\varphi _s^{}\gamma _5\varphi _s`$ (26)
This implies that the chirality of any complex eigenmode is zero,
$`\chi _s\varphi _s^{}\gamma _5\varphi _s=0\text{ if }\lambda _s\lambda _s^{}.`$ (27)
If $`\lambda _s`$ is real ( zero or $`a^1`$ ), then Eqs. (25) and (19) imply that $`\varphi _s`$ has definite chirality $`+1`$ or $`1`$ :
$`\gamma _5\varphi _s=\pm \varphi _s,\text{ if }\lambda _s=\lambda _s^{}.`$ (28)
A useful property of chirality is that the total chirality of all eigenmodes must vanish,
$`{\displaystyle \underset{s}{}}\chi _s`$ $`=`$ $`{\displaystyle \underset{s}{}}\varphi _s^{}\gamma _5\varphi _s`$ (29)
$`=`$ $`{\displaystyle \underset{s}{}}{\displaystyle \underset{x}{}}{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\beta }{}}[\varphi _s^\alpha (x)]^{}\gamma _5^{\alpha \beta }\varphi _s^\beta (x)`$
$`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{\beta }{}}\gamma _5^{\alpha \beta }\delta _{\alpha \beta }=0`$
where the completeness relation
$$\underset{x}{}\underset{s}{}[\varphi _s^\alpha (x)]^{}\varphi _s^\beta (x)=\delta ^{\alpha \beta }$$
(30)
has been used. Since the chirality of any complex eigenmode is zero, then Eq. (29) gives the chirality sum rule for real eigenmodes,
$`N_++n_+=N_{}+n_{}`$ (31)
where $`n_+(n_{})`$ denotes the number of zero modes of positive ( negative ) chirality, and $`N_+(N_{})`$ the number of nonzero real ( $`a^1`$ ) eigenmodes of positive ( negative ) chirality. Then we immediately see that any zero mode must be accompanied by a real ( $`a^1`$ ) eigenmode with opposite chirality, and the index of $`D`$ is
$`\text{index}(D)n_{}n_+=(N_{}N_+).`$ (32)
It should be emphasized that the chiral properties (28), (27) and the chirality sum rule (31) hold for any normal $`D`$ satisfying the $`\gamma _5`$-Hermiticity, as shown in Ref. . However, in nontrivial gauge backgrounds, whether $`D`$ possesses any zero modes or not relies on the topological characteristics of $`D`$, which cannot be guaranteed by the conditions such as the locality, free of species doublings, correct continuum behavior, $`\gamma _5`$-Hermiticity and the GW relation.
## 3 Higher-order realization of Overlap Dirac operator
Now the task is to construct a topologically proper $`D`$ which is local, free of species doubling, $`\gamma _5`$-Hermitian, having correct continuum behavior, and satisfies the GW relation (3). So far, the only viable way to construct a topologically proper $`D`$ is the Overlap,
$`D={\displaystyle \frac{1}{2a}}(\text{1I}+\gamma _5ϵ),ϵ^2=\text{1I}`$ (33)
which satisfies the GW relation (1) with $`R=1`$. There are many different ways to implement the Hermitian $`ϵ`$ in (33). However, it is required to be able to capture the topology of the gauge background. That means, one-half of the difference of the numbers ( $`h_\pm `$ ) of positive ( $`+1`$ ) and negative ( $`1`$ ) eigenvalues of $`ϵ`$ is equal to the background topological charge $`Q`$,
$`{\displaystyle \underset{x}{}}\text{tr}(a\gamma _5D(x,x))={\displaystyle \frac{1}{2}}{\displaystyle \underset{x}{}}\text{tr}(ϵ)={\displaystyle \frac{1}{2}}(h_+h_{})=Q,`$ (34)
where tr denotes the trace over the Dirac and color space. Otherwise, the axial anomaly of $`D`$ cannot agree with the topological charge density in a nontrivial gauge background. Henceforth, we shall regard any lattice Dirac operator which is constructed through the general form of Overlap (33) as a realization of the Overlap Dirac operator.
An explicit realization of $`ϵ`$ in (33) is the Neuberger-Dirac operator with
$`ϵ={\displaystyle \frac{H_w}{\sqrt{H_w^2}}}`$ (35)
where
$`H_w`$ $`=`$ $`\gamma _5(D_Wm_0a^1),0<m_0<2r_w,`$ (36)
$`D_W`$ $`=`$ $`\gamma _\mu t_\mu +W,\text{ }D_W\text{ : massless Wilson-Dirac operator },`$ (37)
$`t_\mu (x,y)`$ $`=`$ $`{\displaystyle \frac{1}{2a}}[U_\mu (x)\delta _{x+\widehat{\mu },y}U_\mu ^{}(y)\delta _{x\widehat{\mu },y}],`$ (38)
$`W(x,y)`$ $`=`$ $`{\displaystyle \frac{r_w}{2a}}{\displaystyle \underset{\mu }{}}\left[2\delta _{x,y}U_\mu (x)\delta _{x+\widehat{\mu },y}U_\mu ^{}(y)\delta _{x\widehat{\mu },y}\right].`$ (39)
( The Wilson parameter $`r_w`$ is usually set to one. )
$`\gamma _\mu =\left(\begin{array}{cc}0& \sigma _\mu \\ \sigma _\mu ^{}& 0\end{array}\right),`$
and
$`\sigma _\mu \sigma _\nu ^{}+\sigma _\nu \sigma _\mu ^{}=2\delta _{\mu \nu }.`$
Note that the parameter $`m_0`$ plays a crucial role in detecting the topology of the gauge background.
Now the problem is how to generalize this construction to $`D`$ satisfying the GW relation with $`R=(a\gamma _5D)^{2k}`$ for $`k>0`$. We can multiply both sides of the GW relation (1) by $`R`$ and redefine $`D^{}=RD`$, then we have
$`D^{}\gamma _5+\gamma _5D^{}=2D^{}\gamma _5D^{}`$ (41)
where $`D^{}=RD`$ is $`\gamma _5`$-hermitian since
$`D^{}=D^{}R=\gamma _5D\gamma _5R=\gamma _5DR\gamma _5=R\gamma _5D\gamma _5=\gamma _5RD\gamma _5=\gamma _5D^{}\gamma _5.`$
Now (41) is in the same form of the GW relation with $`R=1`$. Thus, one can construct $`D^{}`$ in the same way as the Overlap
$`D^{}=RD={\displaystyle \frac{1}{2a}}(\text{1I}+\gamma _5ϵ),ϵ^2=\text{1I},`$ (42)
provided that a proper realization of $`ϵ`$ can be obtained. An explicit construction based on the Neuberger-Dirac operator ( $`k=0`$ ) has been generalized to higher-orders ( $`k>0`$ ) by Fujikawa .
In the following, we formulate Fujikawa’s construction in a more transparent way. Using $`R\gamma _5=\gamma _5R`$, we can write
$`D^{}=RD=(a\gamma _5D)^{2k}D=(a\gamma _5D)^{2k}\gamma _5\gamma _5D=a^1\gamma _5(a\gamma _5D)^{2k+1},`$ (43)
which yields
$`D=a^1\gamma _5(a\gamma _5D^{})^{1/(2k+1)},`$ (44)
where the $`(2k+1)`$-th real root of the Hermitian operator $`a\gamma _5D^{}`$ is assumed.
Then (44) suggests that if the $`ϵ`$ in (42) is expressed in terms of a Hermitian operator $`H`$,
$`ϵ={\displaystyle \frac{H}{\sqrt{H^2}}},`$ (45)
then $`H`$ is required to be proportional to $`(\gamma _\mu 𝒟_\mu )^{2k+1}`$ plus higher-order terms in the continuum limit such that $`D`$ behaves as $`\gamma _\mu 𝒟_\mu `$ after taking the $`(2k+1)`$-th root in (44). Thus, $`H`$ must contain the term $`(\gamma _\mu t_\mu )^{2k+1}`$, where $`\gamma _\mu t_\mu `$ is the naive lattice fermion operator defined in (38). Then additional terms must be required in order to remove the species doublers in the term $`(\gamma _\mu t_\mu )^{2k+1}`$. So, we add the Wilson term to the $`(2k+1)`$-th power, i.e., $`W^{2k+1}`$. Finally, a negative mass term $`(m_0a^1)^{2k+1}`$ is inserted such that $`ϵ`$ is able to detect the topological charge $`Q`$ of the gauge background, i.e.,
$`{\displaystyle \frac{1}{2}}{\displaystyle \underset{x}{}}\text{tr}\left({\displaystyle \frac{H}{\sqrt{H^2}}}\right)=Q.`$ (46)
Putting all these terms together, we have
$`H=\gamma _5[(\gamma _\mu t_\mu )^{2k+1}+W^{2k+1}(m_0a^1)^{2k+1}],`$ (47)
which, at $`k=0`$, reduces to the $`H_w`$ in Eq. (36). Then (44) can be rewritten as
$`D=a^1\left({\displaystyle \frac{1}{2}}\right)^{1/(2k+1)}\gamma _5\left(\gamma _5+{\displaystyle \frac{H}{\sqrt{H^2}}}\right)^{1/(2k+1)},`$ (48)
where $`H`$ is defined in (47). This is the higher-order realization of Overlap Dirac operator, as constructed by Fujikawa .
Next we derive some general properties of the Overlap Dirac operator $`D`$ defined in (48).
The fermion propagator $`S_F(x,y)`$ is defined by
$`S_F(x,y)={\displaystyle \frac{1}{Z}}{\displaystyle \underset{z}{}d\overline{\psi }(z)d\psi (z)\text{e}^{\overline{\psi }D\psi }\psi (x)\overline{\psi }(y)}`$
where
$`Z={\displaystyle \underset{z}{}d\overline{\psi }(z)d\psi (z)\text{e}^{\overline{\psi }D\psi }}`$
In a background gauge field of zero topological charge ( $`Q=0`$ ), the fermion propagator is
$`S_F(x,y)=D^1(x,y).`$ (49)
In the naive continuum limit with $`r_wa^1`$ and $`m_0a^1`$ kept finite, the ( free ) fermion propagator in momentum space can be obtained after some straightforward algebras,
$$\stackrel{~}{S_F}(p)=a\gamma _5\left[\gamma _5\left(\text{1I}+\frac{1}{(\gamma _\mu t_\mu )^{2k+1}}T(p)\right)\right]^{1/(2k+1)}$$
(50)
where
$`t_\mu `$ $`=`$ $`ia^1\mathrm{sin}(p_\mu a),`$ (51)
$`t^2`$ $`=`$ $`a^2{\displaystyle \underset{\mu }{}}\mathrm{sin}^2(p_\mu a),`$ (52)
$`W(p)`$ $`=`$ $`r_wa^1{\displaystyle \underset{\mu }{}}[1\mathrm{cos}(p_\mu a)],`$ (53)
$`u(p)`$ $`=`$ $`[W(p)]^{2k+1}(m_0a^1)^{2k+1},`$ (54)
$`N(p)`$ $`=`$ $`\sqrt{(t^2)^{2k+1}+u^2(p)},`$ (55)
$`T(p)`$ $`=`$ $`N(p)u(p).`$ (56)
Around $`p0`$, $`t_\mu ip_\mu `$, $`T(p)2(m_0a^1)^{2k+1}`$, and the fermion propagator becomes
$`\stackrel{~}{S_F}(p)`$ $``$ $`a\gamma _5\left[\gamma _5\left(\text{1I}+{\displaystyle \frac{1}{(i\gamma _\mu p_\mu )^{2k+1}}}2(m_0a^1)^{2k+1}\right)\right]^{1/(2k+1)}`$ (57)
$``$ $`2^{1/(2k+1)}m_0{\displaystyle \frac{1}{i\gamma _\mu p_\mu }}+a+a(1\delta _{k,0})\mathrm{\Sigma }_k(ap)`$ (58)
where except for $`k=0`$, a momentum-dependent term denoted by $`a\mathrm{\Sigma }_k(ap)`$ is present in the scalar part, which may lead to additive mass renormalization. Evidently, we have to fix the value of $`m_0`$ to
$`m_0=\left({\displaystyle \frac{1}{2}}\right)^{1/(2k+1)}`$ (59)
such that in the limit ( $`a0`$ ) the fermion propagator (58) agrees with the continuum propagator.
For $`m_0(0,2r_w)`$, on a $`d`$-dimensional lattice ( $`d`$ = even ), at any one the $`2^d1`$ corners of the Brillouin zone \[ i.e., $`ap=(\pi ,0,\mathrm{},0)`$, $`(0,\pi ,\mathrm{},0)`$, $`\mathrm{}`$, $`(\pi ,\pi ,\mathrm{},\pi )`$ \], we have $`N(p)=u(p)>0`$, thus $`T(p)=u(p)u(p)=0`$, and all doubled modes are decoupled from the fermion propagator (50).
In general, we consider an arbitrary value of $`m_0`$. At the origin ( $`p=0`$ ) and the $`2^d1`$ corners of the Brillouin zone, we have $`\mathrm{sin}(p_\mu a)=0`$, so $`T(p)`$ becomes
$`T(p)=|u(p)|u(p),`$ (60)
where the possible values of $`u(p)`$ are :
$`u(p)`$ $`=`$ $`(m_0a^1)^{2k+1},`$
$`(2r_wa^1)^{2k+1}(m_0a^1)^{2k+1},`$
$`(4r_wa^1)^{2k+1}(m_0a^1)^{2k+1},`$
$`\mathrm{}`$
$`\mathrm{}`$
$`(2dr_wa^1)^{2k+1}(m_0a^1)^{2k+1}.`$
Here the first value of $`u(p)`$ corresponds to all components of $`p`$ equal to zero, the second value to one of components equal to $`\pi /a`$, and so on, and the last value to all components equal to $`\pi /a`$. Note that
$`(2nr_wa^1)^{2k+1}(m_0a^1)^{2k+1}`$
$`=`$ $`(2nr_wa^1m_0a^1)[(2nr_wa^1)^{2k}+\mathrm{}+(m_0a^1)^{2k}],n=0,1,\mathrm{},d.`$
Therefore the sign of $`u(p)`$ is independent of the order $`k`$. From Eq. (60), we see that if $`u(p)0`$, then $`T(p)=0`$, and this doubled mode is decoupled from the fermion propagator (50) for any order $`k`$. Since the chiral charge of a doubled mode is equal to $`(1)^n`$, where $`n`$ is the number of momentum components equal to $`\pi /a`$, then the total chiral charge $`Q_5`$ of all massless ( primary and doubled ) fermion modes contributing to the fermion propagator (50) can be determined. Then $`\text{index}(D)=Q_5Q`$ for a gauge background with topological charge $`Q`$. Thus, in the naive continuum limit, the index of $`D`$ ( as a function of $`m_0`$ ) is independent of the order $`k`$, same as the index of the Neuberger-Dirac operator ( $`k=0`$ ), which has been determined in Ref. . Explicitly,
$$\text{index}[D(m_0)]=\{\begin{array}{cc}\frac{(1)^{n+1}(d1)!}{(dn)!(n1)!}Q,\hfill & \text{ }2(n1)r_w<m_0<2nr_w\hfill \\ \text{ }\hfill & \text{ for }n=1,\mathrm{},d\text{ ; }\hfill \\ 0,\hfill & \text{ otherwise. }\hfill \end{array}$$
(61)
In particular, for $`d=4`$,
$$\text{index}[D(m_0)]=\{\begin{array}{cc}Q,\hfill & 0<m_0<2r_w,\hfill \\ 3Q,\hfill & 2r_w<m_0<4r_w,\hfill \\ 3Q,\hfill & 4r_w<m_0<6r_w,\hfill \\ Q,\hfill & 6r_w<m_0<8r_w,\hfill \\ 0,\hfill & \text{ otherwise. }\hfill \end{array}$$
(62)
For the Neuberger-Dirac operator ( $`k=0`$ ), there exists an exact discrete symmetry of the index on any finite lattice with even number of sites in each dimension
$`\text{index}[D(m_0)]=\text{index}[D(2dr_wm_0)],\text{ for }k=0,`$ (63)
which holds for any background gauge configuration. However, for higher-order Overlap Dirac operators ( $`k>0`$ ), this discrete symmetry is not exact on a finite lattice; only in the naive continuum limit, this discrete symmetry can be realized as in (61). Nevertheless, at $`m_0=dr_w`$, it can be shown that the index is exactly zero for all $`k`$,
$`\text{index}[D(m_0=dr_w)]=0,\text{ for all }k0,`$ (64)
on any finite lattice with even number of sites in each dimension, and for any background gauge configuration.
## 4 Tests
In this section, we compare the chiral properties of the higher-order Overlap Dirac operators ( $`k=1,2`$ ) to those of the Neuberger-Dirac operator ( $`m_0=1`$ and $`R=1/2`$ ), by computing the fermion propagator, the axial anomaly and the fermion determinant in two-dimensional background $`U(1)`$ gauge fields. Our notations for the two-dimensional background gauge field are the same as those in Ref. ( Eqs. (7)-(11) in Ref. ).
Note that the Neuberger-Dirac operator is conventionally written as
$`D=a^1(\text{1I}+V)=a^1(\text{1I}+\gamma _5{\displaystyle \frac{H_w}{\sqrt{H_w^2}}}),m_0=1;`$ (65)
which satisfies the GW relation with $`R=1/2`$. However, the zeroth order ( $`k=0`$ ) Overlap Dirac operator is
$`D={\displaystyle \frac{1}{2a}}(\text{1I}+\gamma _5{\displaystyle \frac{H_w}{\sqrt{H_w^2}}}),m_0=1/2;`$ (66)
which satisfies the GW relation with $`R=1`$. In the following, we shall use the Neuberger-Dirac operator (65) in place of the zeroth order Overlap (66). All numerical results for the $`k=0`$ case are obtained using the Neuberger-Dirac operator (65) rather than (66).
First of all, we checked that the eigenvalues of a higher-order ( $`k=1,2`$ ) Overlap Dirac operator fall on the deformed circle which is stretched symmetrically in $`\pm \widehat{y}`$ directions. In a nontrivial gauge background, the real eigenmodes ( zero and $`a^1`$ ) have definite chirality and satisfy the chirality sum rule (31), and each complex eigenmode has zero chirality (27). The Atiyah-Singer index theorem ( $`n_{}n_+=Q`$ ) is satisfied in all cases ( $`k=0,1,2`$ ) for gauge configurations fulfiling the topological bound
$`a^2|\overline{\rho }(x)|<ϵ_10.28x`$ (67)
where $`a^2\overline{\rho }(x)`$ is the topological charge inside the unit square of area $`a^2`$ centered at $`x`$,
$`\overline{\rho }(x)={\displaystyle \frac{1}{a^2}}{\displaystyle _{x_1a/2}^{x_1+a/2}}𝑑y_1{\displaystyle _{x_2a/2}^{x_2+a/2}}𝑑y_2{\displaystyle \frac{1}{2\pi }}F_{12}(y).`$ (68)
The value of $`m_0`$ in the higher-order ( $`k>0`$ ) Overlap Dirac operator is fixed according to Eq. (59), $`m_0=2^{1/(2k+1)}`$, while $`m_0=1`$ for the Neuberger-Dirac operator.
### 4.1 Fermion propagator
In the following, we first compute the free fermion propagator on a two dimensional lattice, for the Neuberger-Dirac operator and higher-order ( $`k=1,2`$ ) Overlap Dirac operators respectively. We compare them to the exact solution of the massless fermion propagator on the torus. Then we turn on a background gauge field to examine the behaviors of the scalar part $`S_0`$ and the pseudoscalar part $`S_5`$ in the higher order ( $`k=1,2`$ ) fermion propagators.
In general, the free fermion propagator can be written as
$`S_F(x)=S_0(x)+\gamma _\mu S_\mu (x),`$ (69)
where $`S_0(x)=0`$ for the massless fermion in continuum; and $`S_0(x)=(a/2)\delta _{x,0}`$ for the Neuberger-Dirac operator.
First, we examine the $`S_\mu (x)`$ components in (69). In Table 1, we list the component $`S_1(x)`$ along the diagonal ( $`x_1=x_2`$ ) of a $`16\times 16`$ lattice with antiperiodic boundary conditions. One of the end points of the propagator is fixed at the origin, while the other end point is located at a site along the diagonal. Note that, by symmetry, $`S_2(x)=S_1(x)`$ along the diagonal. From the data in Table 1, we immediately see that in all cases ( Neuberger, $`k=1,2`$ ), $`S_1(x)`$ agrees very well with the exact solution on the torus. In general, all $`S_\mu `$ components of the free fermion propagators are in good agreement with the exact solution for any $`x=(x_1,x_2)`$.
Next, we examine the scalar part $`S_0(x)`$ in the free fermion propagator of the higher-order ( $`k=1,2`$ ) Overlap Dirac operators (48). In Table 2, we list $`S_0(x)`$ along the diagonal of the $`16\times 16`$ lattice with antiperiodic boundary conditions. From the data in Table 2, we see that $`S_0(x)`$ is local for both $`k=1`$ and $`k=2`$. However, we note that $`S_0(x)`$ in the second order ( $`k=2`$ ) case is less localized than that of the first order ( $`k=1`$ ), as expected. It seems that $`|S_0(x)|`$ in both ( $`k=1,2`$ ) orders can be fitted by an exponentially decay function for $`0<|x|<8\sqrt{2}`$.
In a background gauge field, the fermion propagator can be written as
$`S_F(x,y)=\left(\begin{array}{cc}S_0(x,y)+S_5(x,y)& S_R(x,y)\\ S_L(x,y)& S_0(x,y)S_5(x,y)\end{array}\right),`$ (72)
where $`S_0(x,y)=S_5(x,y)=0`$ for the massless fermion in continuum; and $`S_0(x,y)=(a/2)\delta _{x,y}`$ and $`S_5(x,y)=0`$ for the Neuberger-Dirac operator. However, for higher-order ( $`k>0`$ ) Overlap Dirac operators, both $`S_0(x,y)`$ and $`S_5(x,y)`$ are not proportional to $`\delta _{x,y}`$. If $`S_0(x,y)`$ ( $`S_5(x,y)`$ ) turns out to be nonlocal, then it would cause additive mass renormalization and the poles in the fermion propagator will be shifted accordingly.
Now we turn on a background $`U(1)`$ gauge field with parameters ( $`h_1=0.1`$, $`h_2=0.2`$, $`A_1^{(0)}=0.3`$, $`A_2^{(0)}=0.4`$ and $`n_1=n_2=1`$, as defined in Eqs. (7) and (8) in Ref. ). Then we examine the behaviors of $`S_0(x,y)`$ and $`S_5(x,y)`$ for the higher-order ( $`k=1,2`$ ) Overlap Dirac operators. We find that both $`S_0(x,y)`$ and $`S_5(x,y)`$ are local in the higher-order ( $`k=1,2`$ ) fermion propagators. In Table 3, we list the real parts and imaginary parts of $`S_0(x,0)`$ and $`S_5(x,0)`$ for the second order ( $`k=2`$ ) fermion propagator, along the diagonal ( $`x_1=x_2`$ ) of the $`16\times 16`$ lattice.
In general, the scalar part $`S_0(x,y)`$ and the pseudoscalar part $`S_5(x,y)`$ in the higher order ( $`k=1,2`$ ) fermion propagators seem to be local, especially for near continuum gauge configurations. However, one cannot exclude the possibility that they may cause the perturbative instability of the pole of the fermion propagator . On the other hand, for the Neuberger-Dirac operator, we are sure that $`S_0(x,y)=a/2\delta _{x,y}`$ and $`S_5(x,y)=0`$ for any gauge configuration, as well as the perturbative stability of the pole of the fermion propagator . So, from this viewpoint, the chiral properties of Neuberger-Dirac operator are better than those of higher-order Overlap Dirac operators.
### 4.2 Axial anomaly
The axial anomaly of GW Dirac operator $`D`$ satisfying (1) is
$`𝒜_L(x)=a\text{tr}[\gamma _5(RD)(x,x)]`$ (73)
where the trace runs over the Dirac and color space. Substituting $`R=(a\gamma _5D)^{2k}`$ into (73), we obtain
$`𝒜_L(x)=\text{tr}[(a\gamma _5D)^{2k+1}(x,x)].`$ (74)
The sum of the axial anomaly over all sites is equal to the index of $`D`$,
$`{\displaystyle \underset{x}{}}𝒜_L(x)=n_{}n_+.`$ (75)
If the index of $`D`$ is equal to the topological charge $`Q`$ of the gauge background, then the sum of the axial anomaly is equal to $`Q`$. However, it does not necessarily imply that $`𝒜_L(x)`$ would agree with the topological charge density at each site. This happens only when $`D`$ is local.
Since the higher-order ( $`k>0`$ ) Overlap Dirac operator (48) is also topologically proper ( i.e., its index agrees with the background topological charge for any gauge background satisfying the topological bound ), then it follows that its axial anomaly would agree with the topological charge density at each site if $`D`$ is local. ( i.e., the gauge configuration satisfies the locality bound which is more restrictive than the topological bound ).
In the following, we compute the axial anomaly $`𝒜_L(x)`$ in a two-dimensional background $`U(1)`$ gauge field, for the Neuberger-Dirac operator and higher-order ( $`k=1,2`$ ) Overlap Dirac operators respectively. We compare them to the topological charge density $`\overline{\rho }(x)`$ (68) of the gauge background on the torus.
The deviation of the axial anomaly of a lattice Dirac operator in a gauge background can be measured in terms of
$`\delta ={\displaystyle \frac{1}{N_s}}{\displaystyle \underset{x}{}}{\displaystyle \frac{|𝒜_L(x)a^2\overline{\rho }(x)|}{a^2|\overline{\rho }(x)|}}`$ (76)
where $`N_s`$ is the total number of sites of the lattice, and $`\overline{\rho }(x)`$ is the topological charge density inside the square of area $`a^2`$ centered at $`x`$.
In a nontrivial $`U(1)`$ gauge background with parameters $`Q=1`$, $`h_1=0.1`$, $`h_2=0.2`$, $`A_1^{(0)}=0.3`$, $`A_2^{(0)}=0.4`$ and $`n_1=n_2=1`$ ( as defined in Eqs. (7) and (8) in Ref. ) on the $`12\times 12`$ lattice with $`a=1`$, the axial anomaly $`𝒜_L(x)`$ and its deviation $`\delta `$ are computed for the Neuberger-Dirac operator and higher-order ( $`k=1,2`$ ) Overlap Dirac operators respectively. The results are :
$$\delta =\{\begin{array}{cc}0.110,\hfill & \text{Neuberger},\hfill \\ 0.193,\hfill & k=1,\hfill \\ 0.351,\hfill & k=2.\hfill \end{array}$$
(77)
The relatively large deviations of axial anomaly in higher-order ( $`k=1,2`$ ) Overlap Dirac operators indicate that they are less localized than the Neuberger-Dirac operator. And the locality of $`D`$ gets worse as the order $`k`$ goes higher ( i.e., $`D_{k=2}`$ is less localized than $`D_{k=1}`$ ). We have confirmed this by examining $`|D(x,y)|`$ versus $`|xy|`$ explicitly. This implies that the $`ϵ`$ in the locality bound
$`\text{1I}U(p)<ϵ,\text{for all plaquettes}`$ (78)
for a higher order ( $`k>0`$ ) Overlap Dirac operator is more restrictive ( smaller ) than that of the Neuberger-Dirac operator, and it gets smaller as the order goes higher.
For all background gauge configurations we have tested, the Neuberger-Dirac operator always gives anomaly deviation ( $`\delta `$ ) smaller than those of the higher-order ( $`k=1,2`$ ) Overlap Dirac operators.
### 4.3 Fermion determinant
The fermion determinant $`\text{det}(D)`$ is proportional to the exponentiation of the one-loop effective action which is the summation of any number of external sources interacting with one internal fermion loop. It is one of the most crucial quantities to be examined in any lattice fermion formulations. The determinant of $`D`$ is the product of all its eigenvalues
$`\text{det}(D)=(1+e^{\text{i}\pi })^{(n_++n_{})}\text{det}^{}(D)`$ (79)
where $`\text{det}^{}(D)`$ is equal to the product of all non-zero eigenvalues. Since the eigenvalues of $`D`$ are either real of come in complex conjugate pairs, $`\text{det}^{}(D)`$ must be real and positive. For $`Q=0`$, then $`n_++n_{}=0`$ and $`\text{det}(D)=\text{det}^{}(D)`$. For $`Q0`$, then $`n_++n_{}0`$ and $`\text{det}(D)=0`$, but $`\text{det}^{}(D)`$ still provides important information about the spectrum. In continuum, exact solutions of fermion determinants in the general background $`U(1)`$ gauge fields on a torus ( $`L_1\times L_2`$ ) was obtained in Ref. . In the following we compute $`\text{det}^{}(D)`$ for the Neuberger-Dirac operator and higher-order ( $`k=1,2`$ ) Overlap Dirac operators respectively, and then compare them with the exact solutions in continuum. For simplicity, we turn off the harmonic part ( $`h_1=h_2=0`$ ) and the local sinusoidal fluctuations ( $`A_1^{(0)}=A_2^{(0)}=0`$ ) and examine the change of $`\text{det}^{}(D)`$ with respect to the topological charge $`Q`$. For such gauge configurations, the exact solution is
$`\text{det}^{}[D(Q)]=N\sqrt{\left({\displaystyle \frac{L_1L_2}{2|Q|}}\right)^{|Q|}}`$ (80)
where the normalization constant $`N`$ is fixed by
$$N=\sqrt{\frac{2}{L_1L_2}}$$
such that $`\text{det}^{}[D(1)]=1`$.
In Table 4 and Table 5, the fermion determinants $`\text{det}^{}(D)`$ are listed for $`8\times 8`$ and $`16\times 16`$ lattice respectively. It is clear that the Neuberge-Dirac operator always produces results better than those of the higher-order ( $`k=1,2`$ ) Overlap Dirac operators. The first order ( $`k=1`$ ) Overlap Dirac operator performs better than the second order ( $`k=2`$ ) one. This is essentially due to the fact that $`D`$ becomes less localized as the order ( $`k`$ ) goes higher.
## 5 Discussions and Conclusions
We can understand the emergence of Fujikawa’s proposal by the following considerations.
If one requires that $`D`$ is $`\gamma _5`$-hermitian,
$`D^{}=\gamma _5D\gamma _5,`$
and normal,
$`D^{}D=DD^{},`$
( Note that these two conditions are sufficient to guarantee that the real eigenmodes of $`D`$ have definite chirality (28) and satisfy the chirality sum rule (31), and each complex eigenmode has zero chirality (27). ), then one immediately obtains
$`\gamma _5D\gamma _5D=D\gamma _5D\gamma _5.`$
Multiplying above equation by $`\gamma _5`$, we obtain
$`\gamma _5\gamma _5D\gamma _5D=\gamma _5D\gamma _5D\gamma _5,`$
which can be rewritten as
$`\gamma _5(a\gamma _5D)^2=(a\gamma _5D)^2\gamma _5.`$
Since $`(a\gamma _5D)^2`$ is Hermitian and commutes with $`\gamma _5`$, one finds an example of $`R=(a\gamma _5D)^2`$ which depends on $`D`$. Then it is straightforward to generalize this $`R`$ to any powers,
$`R=(a\gamma _5D)^{2k},k=0,1,2,\mathrm{}`$
Substituting this $`R`$ into the GW relation (1), we obtain (3) which is equivalent to Fujikawa’s proposal (2).
It seems to us that Fujikawa’s higher-order realization of the Overlap Dirac operator may not be feasible for practical computations in lattice QCD, in view of its locality, chiral properties and computational accessibility in comparison with those of the Neuberger-Dirac operator. Nevertheless, from a theoretical viewpoint, it has widened our scope and deepened our understanding of the Overlap which does capture one of the fundamental aspects of the nature.
Acknowledgement
This work was supported by the National Science Council, R.O.C. under the grant number NSC89-2112-M002-017. The motivation of this work emerged after a visit to the Center for Subatomic Structure of Matter ( CSSM ) at University of Adelaide, where I had stimulating discussions with Kazuo Fujikawa. I also thank David Adams, Kazuo Fujikawa, Urs. Heller and Tony Williams for interesting discussions at CSSM, and I am grateful to David Adams and Tony Williams for their kind hospitality.
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# Acoustic Signatures in the Primary Microwave Background Bispectrum
## I Introduction
Why measure the bispectrum of the cosmic microwave background (CMB) radiation anisotropy? Simple inflationary models predict that the CMB anisotropy field is nearly random Gaussian, and that two-point statistics completely specify statistical properties of CMB. However, our universe may not be so simple. Higher order statistics, such as the three-point correlation function, or its harmonic transform, the angular bispectrum, are potential probes of the physics of generating the primordial fluctuations. Since gravitationally induced non-linearities are small at $`z1300`$, CMB is expected to be the best probe of the primordial non-Gaussianity.
In the inflationary scenario, the quantum fluctuations of the scalar (inflaton) field generate the observed matter and radiation fluctuations in the universe. In the stochastic inflationary scenario of Starobinsky, the quantum fluctuations decohere to generate the classical fluctuations. There are two potential sources of non-Gaussianity in this inflationary model: (a) the non-linear coupling between the classical inflaton field and the observed fluctuation field, and (b) the non-linear coupling between the quantum noise field and the classical fluctuation field. The former has been investigated by Salopek and Bond, while the latter has been explored by Gangui et al.. Calzetta and Hu and Matacz present an alternative treatment of the decoherence process that leads to different results for the primordial density perturbation from those obtained by Starobinsky. Matacz’s treatment makes similar predictions for the level of non-Gaussianity to the Starobinsky’s treatment. These studies conclude that in the slow roll regime, the fluctuations are Gaussian. However, features in the inflaton potential can produce significant non-Gaussianity.
There have been claims for both the non-detection and the detection of the non-Gaussianity in the COBE map. Banday, Zaroubi and Górski argued the non-cosmological origin of the COBE non-Gaussianity. MAP and Planck will measure the fluctuation field down to angular scales $`0.^{}2`$ and $`0.^{}1`$, and test these claims.
Previous work on the primary non-Gaussianity has focused on very large angular scale, where the temperature fluctuations trace the primordial fluctuations. This is valid on the COBE scale. For MAP and Planck; however, we need the full effect of the radiation transfer function. In this paper, we develop a formalism for doing this, and then present numerical results. Both the formalism and the numerical results are main results of this paper. We also discuss how well we can separate the primary bispectrum from various secondary bispectra.
This paper is organized as follows. Sec. II defines the bispectrum, the Gaunt integral, and particularly the new quantity called the “reduced” bispectrum, which plays a fundamental role in estimating the physical property of the bispectrum. Sec. III formulates the primary bispectrum that uses the full radiation transfer function, and presents the numerical results of the primary bispectrum and the skewness. Sec. IV estimates the secondary bispectra from the coupling between the Sunyaev–Zel’dovich and the weak lensing effects, and from the extragalactic radio and infrared sources. Sec. V studies how well we can measure each bispectrum, and how well we can discriminate among various bispectra. Sec. VI is devoted to further discussion and our conclusion.
## II Defining the “reduced” bispectrum
The observed CMB temperature fluctuation field $`\mathrm{\Delta }T(\widehat{𝐧})/T`$ is expanded into the spherical harmonics:
$$a_{lm}d^2\widehat{𝐧}\frac{\mathrm{\Delta }T(\widehat{𝐧})}{T}Y_{lm}^{}(\widehat{𝐧}),$$
(1)
where hats denote unit vectors. The CMB angular bispectrum is given by
$$B_{l_1l_2l_3}^{m_1m_2m_3}a_{l_1m_1}a_{l_2m_2}a_{l_3m_3},$$
(2)
and the angle-averaged bispectrum is defined by
$$B_{l_1l_2l_3}\underset{m_1m_2m_3}{}\left(\begin{array}{ccc}l_1& l_2& l_3\\ m_1& m_2& m_3\end{array}\right)B_{l_1l_2l_3}^{m_1m_2m_3},$$
(3)
where the matrix is the Wigner-$`3j`$ symbol. The bispectrum $`B_{l_1l_2l_3}^{m_1m_2m_3}`$ must satisfy the triangle conditions and selection rules: $`m_1+m_2+m_3=0`$, $`l_1+l_2+l_3=\mathrm{even}`$, and $`\left|l_il_j\right|l_kl_i+l_j`$ for all permutations of indices. Thus, $`B_{l_1l_2l_3}^{m_1m_2m_3}`$ consists of the Gaunt integral, $`𝒢_{l_1l_2l_3}^{m_1m_2m_3}`$, defined by
$`𝒢_{l_1l_2l_3}^{m_1m_2m_3}`$ $``$ $`{\displaystyle d^2\widehat{𝐧}Y_{l_1m_1}(\widehat{𝐧})Y_{l_2m_2}(\widehat{𝐧})Y_{l_3m_3}(\widehat{𝐧})}`$ (4)
$`=`$ $`\sqrt{{\displaystyle \frac{\left(2l_1+1\right)\left(2l_2+1\right)\left(2l_3+1\right)}{4\pi }}}\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)\left(\begin{array}{ccc}l_1& l_2& l_3\\ m_1& m_2& m_3\end{array}\right).`$ (9)
$`𝒢_{l_1l_2l_3}^{m_1m_2m_3}`$ is real, and satisfies all the conditions mentioned above.
Given the rotational invariance of the universe, $`B_{l_1l_2l_3}`$ is written as
$$B_{l_1l_2l_3}^{m_1m_2m_3}=𝒢_{l_1l_2l_3}^{m_1m_2m_3}b_{l_1l_2l_3},$$
(10)
where $`b_{l_1l_2l_3}`$ is an arbitrary real symmetric function of $`l_1`$, $`l_2`$, and $`l_3`$. This form of equation (10) is necessary and sufficient to construct generic $`B_{l_1l_2l_3}^{m_1m_2m_3}`$ under the rotational invariance. Thus, we shall frequently use $`b_{l_1l_2l_3}`$ instead of $`B_{l_1l_2l_3}^{m_1m_2m_3}`$ in this paper, and call this function the “reduced” bispectrum, as $`b_{l_1l_2l_3}`$ contains all physical information in $`B_{l_1l_2l_3}^{m_1m_2m_3}`$. Since the reduced bispectrum does not contain the Wigner-$`3j`$ symbol that merely ensures the triangle conditions and selection rules, it is easier to calculate and useful to quantify the physical properties of the bispectrum.
The observable quantity, the angle-averaged bispectrum $`B_{l_1l_2l_3}`$, is obtained by substituting equation (10) into (3),
$$B_{l_1l_2l_3}=\sqrt{\frac{(2l_1+1)(2l_2+1)(2l_3+1)}{4\pi }}\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)b_{l_1l_2l_3},$$
(11)
where we have used the identity:
$$\underset{m_1m_2m_3}{}\left(\begin{array}{ccc}l_1& l_2& l_3\\ m_1& m_2& m_3\end{array}\right)𝒢_{l_1l_2l_3}^{m_1m_2m_3}=\sqrt{\frac{(2l_1+1)(2l_2+1)(2l_3+1)}{4\pi }}\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right).$$
(12)
Alternatively, one can define the bispectrum in the flat-sky approximation,
$$a(𝐥_1)a(𝐥_1)a(𝐥_3)=(2\pi )^2\delta ^{(2)}\left(𝐥_1+𝐥_2+𝐥_3\right)B(𝐥_1,𝐥_2,𝐥_3),$$
(13)
where $`𝐥`$ is the two dimensional wave-vector on the sky. This definition of $`B(𝐥_1,𝐥_2,𝐥_3)`$ corresponds to equation (10), given the correspondence of $`𝒢_{l_1l_2l_3}^{m_1m_2m_3}\delta ^{(2)}\left(𝐥_1+𝐥_2+𝐥_3\right)`$ in the flat-sky limit. Thus,
$$b_{l_1l_2l_3}B(𝐥_1,𝐥_2,𝐥_3)\text{(flat-sky approximation)},$$
(14)
is satisfied. This fact also would motivate us to use the reduced bispectrum $`b_{l_1l_2l_3}`$ rather than the angular averaged bispectrum $`B_{l_1l_2l_3}`$. Note that $`b_{l_1l_2l_3}`$ is similar to $`\widehat{B}_{l_1l_2l_3}`$ defined by Magueijo. The relation is $`b_{l_1l_2l_3}=\sqrt{4\pi }\widehat{B}_{l_1l_2l_3}`$.
## III Primary Bispectrum and Skewness
### A Model of the primordial non-Gaussianity
If the primordial fluctuations are adiabatic scalar fluctuations, then
$$a_{lm}=4\pi (i)^l\frac{d^3𝐤}{(2\pi )^3}\mathrm{\Phi }(𝐤)g_{Tl}(k)Y_{lm}^{}(\widehat{𝐤}),$$
(15)
where $`\mathrm{\Phi }(𝐤)`$ is the primordial curvature perturbation in the Fourier space, and $`g_{Tl}(k)`$ is the radiation transfer function. $`a_{lm}`$ thus takes over the non-Gaussianity, if any, from $`\mathrm{\Phi }(𝐤)`$. Although equation (15) is valid only if the universe is flat, it is straightforward to extend this to an arbitrary geometry. The isocurvature fluctuations can be similarly calculated by using the entropy perturbation and the proper transfer function.
In this paper, we explore the simplest weak non-linear coupling case:
$$\mathrm{\Phi }(𝐱)=\mathrm{\Phi }_L(𝐱)+f_{NL}\left(\mathrm{\Phi }_L^2(𝐱)\mathrm{\Phi }_L^2(𝐱)\right),$$
(16)
in real space, where $`\mathrm{\Phi }_L(𝐱)`$ denotes the linear gaussian part of the perturbation. $`\mathrm{\Phi }(𝐱)=0`$ is guaranteed. Henceforth, we shall call $`f_{NL}`$ the non-linear coupling constant. This model is based upon the slow-roll inflationary scenario. Salopek and Bond and Gangui et al. found that $`f_{NL}`$ is given by a certain combination of the slope and the curvature of the inflaton potential. In the notation of Gangui et al., $`\mathrm{\Phi }_3=2f_{NL}`$. Gangui et al. found that $`\mathrm{\Phi }_310^2`$ in the quadratic and the quartic inflaton potential models.
In the Fourier space, $`\mathrm{\Phi }(𝐤)`$ is decomposed into two parts:
$$\mathrm{\Phi }(𝐤)=\mathrm{\Phi }_L(𝐤)+\mathrm{\Phi }_{NL}(𝐤),$$
(17)
and accordingly,
$$a_{lm}=a_{lm}^L+a_{lm}^{NL},$$
(18)
where $`\mathrm{\Phi }_{NL}(𝐤)`$ is the non-linear part defined by
$$\mathrm{\Phi }_{NL}(𝐤)f_{NL}\left[\frac{d^3𝐩}{(2\pi )^3}\mathrm{\Phi }_L(𝐤+𝐩)\mathrm{\Phi }_L^{}(𝐩)(2\pi )^3\delta ^{(3)}(𝐤)\mathrm{\Phi }_L^2(𝐱)\right].$$
(19)
One can confirm that $`\mathrm{\Phi }(𝐤)=0`$ is satisfied. In this model, a non-vanishing component of the $`\mathrm{\Phi }(𝐤)`$-field bispectrum is
$$\mathrm{\Phi }_L(𝐤_1)\mathrm{\Phi }_L(𝐤_2)\mathrm{\Phi }_{NL}(𝐤_3)=2(2\pi )^3\delta ^{(3)}(𝐤_1+𝐤_2+𝐤_3)f_{NL}P_\mathrm{\Phi }(k_1)P_\mathrm{\Phi }(k_2),$$
(20)
where $`P_\mathrm{\Phi }(k)`$ is the linear power spectrum given by $`\mathrm{\Phi }_L(𝐤_1)\mathrm{\Phi }_L(𝐤_2)=(2\pi )^3P_\mathrm{\Phi }(k_1)\delta ^{(3)}(𝐤_1+𝐤_2)`$. We have also used $`\mathrm{\Phi }_L(𝐤+𝐩)\mathrm{\Phi }_L^{}(𝐩)=(2\pi )^3P_\mathrm{\Phi }(p)\delta ^{(3)}(𝐤)`$, and $`\mathrm{\Phi }_L^2(𝐱)=(2\pi )^3d^3𝐤P_\mathrm{\Phi }(k)`$.
Substituting equation (15) into (2), using equation (20) for the $`\mathrm{\Phi }(𝐤)`$-field bispectrum, and then integrating over angles $`\widehat{𝐤}_1`$, $`\widehat{𝐤}_3`$, and $`\widehat{𝐤}_3`$, we obtain the primary CMB angular bispectrum,
$`B_{l_1l_2l_3}^{m_1m_2m_3}`$ $`=`$ $`a_{l_1m_1}^La_{l_2m_2}^La_{l_3m_3}^{NL}+a_{l_1m_1}^La_{l_2m_2}^{NL}a_{l_3m_3}^L+a_{l_1m_1}^{NL}a_{l_2m_2}^La_{l_3m_3}^L`$ (21)
$`=`$ $`2𝒢_{l_1l_2l_3}^{m_1m_2m_3}{\displaystyle _0^{\mathrm{}}}r^2𝑑r\left[b_{l_1}^L(r)b_{l_2}^L(r)b_{l_3}^{NL}(r)+b_{l_1}^L(r)b_{l_2}^{NL}(r)b_{l_3}^L(r)+b_{l_1}^{NL}(r)b_{l_2}^L(r)b_{l_3}^L(r)\right],`$ (22)
where
$`b_l^L(r)`$ $``$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{}}}k^2𝑑kP_\mathrm{\Phi }(k)g_{Tl}(k)j_l(kr),`$ (23)
$`b_l^{NL}(r)`$ $``$ $`{\displaystyle \frac{2}{\pi }}{\displaystyle _0^{\mathrm{}}}k^2𝑑kf_{NL}g_{Tl}(k)j_l(kr).`$ (24)
Note that $`b_l^L(r)`$ is a dimensionless quantity, while $`b_l^{NL}(r)`$ has a dimension of $`L^3`$.
One confirms that the form of equation (10) holds. Thus, the reduced bispectrum, $`b_{l_1l_2l_3}=B_{l_1l_2l_3}^{m_1m_2m_3}\left(𝒢_{l_1l_2l_3}^{m_1m_2m_3}\right)^1`$ (Eq.(10)), for the primordial non-Gaussianity is
$`b_{l_1l_2l_3}^{primary}=2{\displaystyle _0^{\mathrm{}}}r^2𝑑r\left[b_{l_1}^L(r)b_{l_2}^L(r)b_{l_3}^{NL}(r)+b_{l_1}^L(r)b_{l_2}^{NL}(r)b_{l_3}^L(r)+b_{l_1}^{NL}(r)b_{l_2}^L(r)b_{l_3}^L(r)\right].`$ (25)
$`b_{l_1l_2l_3}^{primary}`$ is fully specified by a single constant parameter $`f_{NL}`$, as the cosmological parameters will be precisely determined by measuring the CMB angular power spectrum $`C_l`$ (e.g., ). It should be stressed again that this is the special case in the slow-roll limit. If the slow-roll condition is not satisfied, then $`f_{NL}=f_{NL}(k_1,k_2,k_3)`$ at equation (20). Wang and Kamionkowski have developed the formula to compute $`B_{l_1l_2l_3}`$ from the generic form of $`\mathrm{\Phi }(𝐤)`$-field bispectrum. Our formula (Eq.(22)) agrees with theirs, given our form of the $`\mathrm{\Phi }(𝐤)`$-field bispectrum (Eq.(20)).
Even if the inflation produced Gaussian fluctuations, Pyne and Carroll pointed out that the general relativistic second-order perturbation theory would produce terms of $`f_{NL}𝒪(1)`$. For generic slow-roll models, these terms dominate the primary non-Gaussianity.
### B Numerical results of the primary bispectrum
We evaluate the primary CMB bispectrum (Eqs.(22)–(25)) numerically. We compute the full radiation transfer function $`g_{Tl}(k)`$ with the CMBFAST code, and assume the single power law spectrum, $`P_\mathrm{\Phi }(k)k^{n4}`$, for the primordial curvature fluctuations. The integration over $`k`$ (Eqs.(23) and (24)) is done by the algorithm used in CMBFAST. The cosmological model is the scale-invariant standard cold dark matter model with $`\mathrm{\Omega }_m=1`$, $`\mathrm{\Omega }_\mathrm{\Lambda }=0`$, $`\mathrm{\Omega }_b=0.05`$, $`h=0.5`$, and $`n=1`$, and with the power spectrum $`P_\mathrm{\Phi }(k)`$ normalized to COBE. Although this model is almost excluded by current observations, it is still useful to depict the basic effects of the transfer function on the bispectrum.
Figure 1 shows $`b_l^L(r)`$ (Eq.(23)) and $`b_l^{NL}(r)`$ (Eq.(24)) for several different values of $`r`$. $`r=c\left(\tau _0\tau \right)`$, where $`\tau `$ is the conformal time, and $`\tau _0`$ is at the present. In our model, $`c\tau _0=11.8\mathrm{Gpc}`$, and the decoupling epoch occurs at $`c\tau _{}=235\mathrm{Mpc}`$ at which the differential visibility has a maximum. Our $`c\tau _0`$ includes the radiation effect on the expansion of universe, otherwise $`c\tau _0=12.0\mathrm{Gpc}`$. $`\tau _{}`$ is the epoch when the most of the primary signal is generated. $`b_l^L(r)`$ and $`C_l`$ look very similar one another in the shape and the amplitude at $`l100`$, although the amplitude in the Sachs–Wolfe regime is different by a factor of $`3`$. This is because $`C_l`$ is proportional to $`P_\mathrm{\Phi }(k)g_{Tl}^2(k)`$, while $`b_l^L(r)P_\mathrm{\Phi }(k)g_{Tl}(k)`$, where $`g_{Tl}=1/3`$. $`b_l^L(r)`$ has a good phase coherence over wide range of $`r`$, while the phase of $`b_l^{NL}(r)`$ in high-$`l`$ regime oscillates rapidly as a function of $`r`$. This strongly damps the integrated result of the bispectrum (Eq.(22)) in high-$`l`$ regime. The main difference between $`C_l`$ and $`b_l(r)`$ is that $`b_l(r)`$ changes a sign, while $`C_l`$ does not.
Looking at figure 1, we find $`l^2b_l^L2\times 10^9`$ and $`b_l^{NL}f_{NL}^110^{10}\mathrm{Mpc}^3`$. The most signal coming from the decoupling, the volume element at $`\tau _{}`$ is $`r_{}^2\mathrm{\Delta }r_{}(10^4)^2\times 10^2\mathrm{Mpc}^3`$, and thus we estimate an order of magnitude of the primary reduced bispectrum (Eq.(25)) as
$$b_{lll}^{primary}l^4\left[2r_{}^2\mathrm{\Delta }r_{}\left(l^2b_l^L\right)^2b_l^{NL}\times 3\right]l^4\times 2\times 10^{17}f_{NL}.$$
(26)
Since $`b_l^{NL}f_{NL}^1r_{}^2\delta (rr_{})`$ (see Eq.(29)), $`r_{}^2\mathrm{\Delta }r_{}b_l^{NL}f_{NL}^11`$. This rough estimate agrees with the numerical result below (figure 2).
Figure 2 shows the integrated bispectrum (Eq.(22)) divided by the Gaunt integral $`𝒢_{l_1l_2l_3}^{m_1m_2m_3}`$, which is basically $`b_{l_1l_2l_3}^{primary}`$. Since the signal comes primarily from the decoupling epoch $`\tau _{}`$ as mentioned above, the integration boundary is chosen as $`c\left(\tau _02\tau _{}\right)rc\left(\tau _00.1\tau _{}\right)`$. We use a step-size of $`0.1c\tau _{}`$, as we found that a step size of $`0.01c\tau _{}`$ gives very similar results. While the bispectrum is a 3-d function, we show different 1-d slices of the bispectrum in this figure. $`ł_2(l_2+1)l_3(l_3+1)a_{l_1m_1}^{NL}a_{l_2m_2}^La_{l_3m_3}^L\left(𝒢_{l_1l_2l_3}^{m_1m_2m_3}\right)^1/(2\pi )^2`$ is plotted as a function of $`l_3`$ in the upper panel, while $`l_1(l_1+1)l_2(l_2+1)a_{l_1m_1}^La_{l_2m_2}^La_{l_3m_3}^{NL}\left(𝒢_{l_1l_2l_3}^{m_1m_2m_3}\right)^1/(2\pi )^2`$ is plotted in the lower panel. $`l(l+1)/(2\pi )`$ is multiplied for each $`b_l^L(r)`$ which contains $`P_\mathrm{\Phi }(k)`$ so as the Sachs–Wolfe plateau at $`l_310`$ is easily seen in figure 2. $`l_1`$ and $`l_2`$ are chosen so as $`(l_1,l_2)=(9,11),(99,101),(199,201)`$, and $`(499,501)`$. We find that the $`(l_1,l_2)=(199,201)`$ mode, the first acoustic peak mode, has the largest signal in this family of parameters. The upper panel has a prominent first acoustic peak, and strongly damped oscillations in high-$`l`$ regime. The lower panel also has a first peak, but damps more slowly. The typical amplitude of the reduced bispectrum is $`l^4b_{lll}^{primary}f_{NL}^110^{17}`$, which agrees with an order of magnitude estimate (Eq.(26)).
Our formula (Eq.(25)) and numerical results agree with Gangui et al. calculation in the Sachs–Wolfe regime, where $`g_{Tl}(k)j_l(kr_{})/3`$, and thus
$$b_{l_1l_2l_3}^{primary}6f_{NL}\left(C_{l_1}^{SW}C_{l_2}^{SW}+C_{l_1}^{SW}C_{l_3}^{SW}+C_{l_2}^{SW}C_{l_3}^{SW}\right)\text{(Sachs–Wolfe approximation)}.$$
(27)
Each term is in the same order as equation (25). $`C_l^{SW}`$ is the CMB angular power spectrum in the Sachs–Wolfe approximation,
$$C_l^{SW}\frac{2}{9\pi }_0^{\mathrm{}}k^2𝑑kP_\mathrm{\Phi }(k)j_l^2(kr_{}).$$
(28)
In deriving equation (27) from (25), we approximated $`b_l^{NL}(r)`$ (Eq.(24)) to
$$b_l^{NL}(r)\left(\frac{f_{NL}}{3}\right)\frac{2}{\pi }_0^{\mathrm{}}k^2𝑑kj_l(kr_{})j_l(kr)=\frac{f_{NL}}{3}r_{}^2\delta (rr_{}).$$
(29)
The Sachs–Wolfe approximation (Eq.(27)) is valid only when $`l_1`$, $`l_2`$, and $`l_3`$ are all less than $`10`$, where Gangui et al.’s formula gives $`6\times 10^{20}`$ in figure 2. It should be stressed again that the Sachs–Wolfe approximation gives the qualitatively different result from our full calculation (Eq.(25)) at $`l_i10`$. The full bispectrum does change a sign, while the approximation never changes a sign because of the use of $`C_l^{SW}`$. The acoustic oscillation and the sign change are actually great advantages, when we try to separate the primary bispectrum from various secondary bispectra. We shall study this point later.
### C Primary skewness
The skewness $`S_3`$,
$$S_3\left(\frac{\mathrm{\Delta }T(\widehat{𝐧})}{T}\right)^3$$
(30)
is the simplest statistic characterizing the non-Gaussianity. $`S_3`$ is expanded in terms of $`B_{l_1l_2l_3}`$ (Eq.(3)) or $`b_{l_1l_2l_3}`$ (Eq.(10)) as
$`S_3`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{l_1l_2l_3}{}}\sqrt{{\displaystyle \frac{\left(2l_1+1\right)\left(2l_2+1\right)\left(2l_3+1\right)}{4\pi }}}\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)B_{l_1l_2l_3}W_{l_1}W_{l_2}W_{l_3}`$ (33)
$`=`$ $`{\displaystyle \frac{1}{2\pi ^2}}{\displaystyle \underset{2l_1l_2l_3}{}}\left(l_1+{\displaystyle \frac{1}{2}}\right)\left(l_2+{\displaystyle \frac{1}{2}}\right)\left(l_3+{\displaystyle \frac{1}{2}}\right)\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)^2b_{l_1l_2l_3}W_{l_1}W_{l_2}W_{l_3},`$ (36)
where $`W_l`$ is the experimental window function. We have used equation (11) to replace $`B_{l_1l_2l_3}`$ by the reduced bispectrum $`b_{l_1l_2l_3}`$ in the last equality. Since $`l=0`$ and $`1`$ modes are not observable, we have excluded them from the summation. Throughout this paper, we consider the single-beam window function, $`W_l=e^{l(l+1)/(2\sigma _b^2)}`$, where $`\sigma _b=\mathrm{FWHM}/\sqrt{8\mathrm{ln}2}`$. Since $`\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)^2b_{l_1l_2l_3}`$ is symmetric under permutation of indices, it is useful to change the way of summation as
$$\underset{2l_1l_2l_3}{}6\underset{2l_1l_2l_3}{}.$$
(37)
Since this reduces the number of summations by a factor of $`6`$, we shall use this convention henceforth.
The upper panel of figure 3 plots $`S_3(<l_3)`$, which is $`S_3`$ summed up to a certain $`l_3`$, for FWHM beam-sizes of $`7^{}`$, $`13^{}`$, and $`5^{}.5`$. These values correspond to beam-sizes of COBE, MAP, and Planck experiments, respectively. Figure 3 also plots the infinitesimal thin-beam case. MAP, Planck, and the ideal experiments measure very similar $`S_3`$ one another, despite the fact that Planck and the ideal experiments can use much more number of modes than MAP. The reason is as follows. Looking at equation (36), one finds that $`S_3`$ is the linear integration of $`b_{l_1l_2l_3}`$ over $`l_i`$. Thus, integrating oscillations in $`b_{l_1l_2l_3}^{primary}`$ around zero (see figure 2) damps the non-Gaussian signal in small angular scales, $`l300`$. Since the COBE scale is basically dominated by the Sachs–Wolfe effect, no oscillation, the cancellation affects $`S_3`$ less significantly than in MAP and Planck scales, while Planck suffers from severe cancellation in small angular scales. Even Planck and the ideal experiments measure the same amount of $`S_3`$ as MAP does. As a result, measured $`S_3`$ almost saturates at the MAP resolution scale, $`l500`$.
We conclude this section by noting that when we can calculate the expected form of the bispectrum, then it is a “matched filter” for detecting the non-Gaussianity in the data, and thus much more powerful tool than the skewness in which the information is lost through the coarse-graining.
## IV Secondary sources of the CMB bispectrum
Even if the CMB bispectrum were significantly detected in the CMB map, the origin would not necessarily be primordial, but rather would be various foregrounds such as the Sunyaev–Zel’dovich effect (hereafter SZ), the weak lensing effect, extragalactic radio sources, and so on. In order to isolate the primordial origin from others, we have to know the accurate form of bispectra produced by the foregrounds.
### A Coupling between the weak lensing and the Sunyaev–Zel’dovich effects
The coupling between the SZ and the weak lensing effects would produce an observable effect in the bispectrum. The CMB temperature field including the SZ and the lensing effects is expanded as
$$\frac{\mathrm{\Delta }T(\widehat{𝐧})}{T}=\frac{\mathrm{\Delta }T^P\left(\widehat{𝐧}+\mathrm{\Theta }(\widehat{𝐧})\right)}{T}+\frac{\mathrm{\Delta }T^{SZ}(\widehat{𝐧})}{T}\frac{\mathrm{\Delta }T^P(\widehat{𝐧})}{T}+\left(\frac{\mathrm{\Delta }T^P(\widehat{𝐧})}{T}\right)\mathrm{\Theta }(\widehat{𝐧})+\frac{\mathrm{\Delta }T^{SZ}(\widehat{𝐧})}{T},$$
(38)
where $`P`$ denotes the primary anisotropy, $`\mathrm{\Theta }(\widehat{𝐧})`$ is the lensing potential:
$$\mathrm{\Theta }(\widehat{𝐧})2_0^r_{}𝑑r\frac{r_{}r}{rr_{}}\mathrm{\Phi }(r,\widehat{𝐧}r),$$
(39)
and $`SZ`$ denotes the SZ effect:
$$\frac{\mathrm{\Delta }T^{SZ}(\widehat{𝐧})}{T}=y(\widehat{𝐧})j_\nu ,$$
(40)
where $`j_\nu `$ is the spectral function of the SZ effect. $`y(\widehat{𝐧})`$ is the Compton $`y`$-parameter given by
$$y(\widehat{𝐧})y_0\frac{dr}{r_{}}\frac{T_\rho (r,\widehat{𝐧}r)}{\overline{T}_{\rho 0}}a^2(r),$$
(41)
where
$$y_0\frac{\sigma _T\overline{\rho }_{gas0}k_B\overline{T}_{\rho 0}r_{}}{\mu _em_pm_ec^2}=4.3\times 10^4\mu _e^1\left(\mathrm{\Omega }_bh^2\right)\left(\frac{k_B\overline{T}_{\rho 0}}{1\mathrm{keV}}\right)\left(\frac{r_{}}{10\mathrm{Gpc}}\right).$$
(42)
$`T_\rho \rho _{gas}T_e/\overline{\rho }_{gas}`$ is the electron temperature weighted by the gas mass density, the overline denotes the volume average, and the subscript 0 means the present epoch. We adopt $`\mu _e^1=0.88`$, where $`\mu _e^1n_e/(\rho _{gas}/m_p)`$ is the number of electrons per proton mass in fully ionized medium. Other quantities have their usual meanings.
Transforming equation (38) into harmonic space,
$`a_{lm}`$ $`=`$ $`a_{lm}^P+{\displaystyle \underset{l^{}m^{}}{}}{\displaystyle \underset{l^{\prime \prime }m^{\prime \prime }}{}}(1)^m𝒢_{ll^{}l^{\prime \prime }}^{mm^{}m^{\prime \prime }}{\displaystyle \frac{l^{}(l^{}+1)l(l+1)+l^{\prime \prime }(l^{\prime \prime }+1)}{2}}a_{l^{}m^{}}^P\mathrm{\Theta }_{l^{\prime \prime }m^{\prime \prime }}+a_{lm}^{SZ}`$ (43)
$`=`$ $`a_{lm}^P+{\displaystyle \underset{l^{}m^{}}{}}{\displaystyle \underset{l^{\prime \prime }m^{\prime \prime }}{}}(1)^{m+m^{}+m^{\prime \prime }}𝒢_{ll^{}l^{\prime \prime }}^{mm^{}m^{\prime \prime }}{\displaystyle \frac{l^{}(l^{}+1)l(l+1)+l^{\prime \prime }(l^{\prime \prime }+1)}{2}}a_{l^{}m^{}}^P\mathrm{\Theta }_{l^{\prime \prime }m^{\prime \prime }}^{}+a_{lm}^{SZ},`$ (44)
where $`𝒢_{l_1l_2l_3}^{m_1m_2m_3}`$ is the Gaunt integral (Eq.(9)). Substituting equation (44) into (2), and using the identity $`𝒢_{l_1l_2l_3}^{m_1m_2m_3}=𝒢_{l_1l_2l_3}^{m_1m_2m_3}`$, we obtain the bispectrum,
$$B_{l_1l_2l_3}^{m_1m_2m_3}=𝒢_{l_1l_2l_3}^{m_1m_2m_3}\left[\frac{l_1(l_1+1)l_2(l_2+1)+l_3(l_3+1)}{2}C_{l_1}^P\mathrm{\Theta }_{l_3m_3}^{}a_{l_3m_3}^{SZ}+\text{5 permutations}\right].$$
(45)
The form of equation (10) is confirmed, and then the reduced bispectrum $`b_{l_1l_2l_3}^{szlens}`$ includes terms in the square bracket.
The cross-correlation power spectrum of the lensing and the SZ effects, $`\mathrm{\Theta }_{lm}^{}a_{lm}^{SZ}`$, appearing in equation (45) was first derived by Goldberg and Spergel. They assumed the linear pressure bias model proposed by Persi et al.: $`T_\rho =\overline{T}_\rho b_{gas}\delta `$, and the mean temperature evolution of $`\overline{T}_\rho \overline{T}_{\rho 0}(1+z)^1`$ for $`z<2`$ as roughly suggested by recent hydrodynamic simulations. Then they derived
$$\mathrm{\Theta }_{lm}^{}a_{lm}^{SZ}j_\nu \frac{4y_0b_{gas}l^2}{3\mathrm{\Omega }_mH_0^2}_0^z_{}𝑑z\frac{dr}{dz}D^2(z)(1+z)^2\frac{r_{}r(z)}{r_{}^2r^5(z)}P_\mathrm{\Phi }\left(k=\frac{l}{r(z)}\right),$$
(46)
where $`D(z)`$ is the linear growth factor. Simulations without non-gravitational heating suggest that $`\overline{T}_{\rho 0}0.20.4\mathrm{keV}`$ and $`b_{gas}510`$, and similar numbers are obtained by analytic estimations. In this pressure bias model, free parameters except cosmological parameters are $`\overline{T}_{\rho 0}`$ and $`b_{gas}`$. However, both actually depend on cosmological models. Since $`l^3\mathrm{\Theta }_{lm}^{}a_{lm}^{SZ}2\times 10^{10}j_\nu \overline{T}_{\rho 0}b_{gas}`$ and $`l^2C_l^P6\times 10^{10}`$,
$$b_{lll}^{szlens}l^3\left[\left(l^2C_l^P\right)\left(l^3\mathrm{\Theta }_{lm}^{}a_{lm}^{SZ}\right)\times 5/2\right]l^3\times 3\times 10^{19}j_\nu \overline{T}_{\rho 0}b_{gas},$$
(47)
where $`\overline{T}_{\rho 0}`$ is in units of 1 keV, and $`b_{l_1l_2l_3}=B_{l_1l_2l_3}^{m_1m_2m_3}\left(𝒢_{l_1l_2l_3}^{m_1m_2m_3}\right)^1`$ (Eq.(10)) is the reduced bispectrum. Thus, comparing this to equation (26), we obtain
$$\frac{b_{lll}^{primary}}{b_{lll}^{szlens}}l^1\times 10\left(\frac{f_{NL}}{j_\nu \overline{T}_{\rho 0}b_{gas}}\right).$$
(48)
This estimate suggests that the primary bispectrum is overwhelmed by the SZ–lensing bispectrum in small angular scales. This is why we have to separate the primary from the SZ–lensing effect.
### B Extragalactic radio and infrared sources
The bispectrum from extragalactic radio and infrared sources whose fluxes $`F`$ are less than a certain detection threshold $`F_d`$ is relatively simple to estimate, when they are assumed to be Poisson distributed. Since the Poisson distribution has the white noise spectrum, the reduced bispectrum (Eq.(10)) is constant, $`b_{l_1l_2l_3}^{ps}=b^{ps}=\mathrm{constant}`$, then we obtain
$$B_{l_1l_2l_3}^{m_1m_2m_3}=𝒢_{l_l1_2l_3}^{m_1m_2m_3}b^{ps},$$
(49)
where
$$b^{ps}(<F_d)g^3(x)_0^{F_d}𝑑FF^3\frac{dn}{dF}=g^3(x)\frac{n(>F_d)}{3\beta }F_d^3.$$
(50)
The assumption of the Poisson distribution is fairly good approximation as found by Toffolatti et al.. $`dn/dF`$ is the differential source count per unit solid angle, and $`n(>F_d)_{F_d}^{\mathrm{}}𝑑F(dn/dF)`$. The power law count, $`dn/dFF^{\beta 1}`$ with $`\beta <2`$, has been assumed. $`xh\nu /k_BT(\nu /56.80\mathrm{GHz})(T/2.726\mathrm{K})^1`$, and
$$g(x)2\frac{(hc)^2}{(k_BT)^3}\left(\frac{\mathrm{sinh}x/2}{x^2}\right)^2\frac{1}{67.55\mathrm{MJy}\mathrm{sr}^1}\left(\frac{T}{2.726\mathrm{K}}\right)^3\left(\frac{\mathrm{sinh}x/2}{x^2}\right)^2.$$
(51)
$`b^{ps}`$ is otherwise written in terms of the Poisson angular power spectrum $`C^{ps}`$:
$$C^{ps}(<F_d)g^2(x)_0^{F_d}𝑑FF^2\frac{dn}{dF}=g^2(x)\frac{n(>F_d)}{2\beta }F_d^2,$$
(52)
as
$$b^{ps}(<F_d)=\frac{(2\beta )^{3/2}}{3\beta }\left[n(>F_d)\right]^{1/2}\left[C^{ps}(<F_d)\right]^{3/2}.$$
(53)
Toffolatti et al. estimated $`n(>F_d)300\mathrm{sr}^1`$ for $`F_d0.2\mathrm{Jy}`$ at 217 GHz. This $`F_d`$ corresponds to $`5\sigma `$ detection threshold for Planck experiment at 217 GHz. Refregier, Spergel and Herbig extrapolated Toffolatti et al.’s estimation to 94 GHz, and obtained $`n(>F_d)7\mathrm{sr}^1`$ for $`F_d2\mathrm{Jy}`$, which corresponds to MAP $`5\sigma `$ threshold. These values yield
$`C^{ps}(90\mathrm{GHz},<2\mathrm{Jy})`$ $``$ $`2\times 10^{16},`$ (54)
$`C^{ps}(217\mathrm{GHz},<0.2\mathrm{Jy})`$ $``$ $`1\times 10^{17}.`$ (55)
Thus, rough estimates for $`b^{ps}`$ are
$`b^{ps}(90\mathrm{GHz},<2\mathrm{Jy})`$ $``$ $`2\times 10^{25},`$ (56)
$`b^{ps}(217\mathrm{GHz},<0.2\mathrm{Jy})`$ $``$ $`5\times 10^{28}.`$ (57)
While we assumed the Euclidean source count ($`\beta =3/2`$) here for definiteness, this does not affect an order of magnitude estimates above. Since the primary reduced bispectrum $`l^4`$ (Eq.(26)) and the SZ–lensing reduced bispectrum $`l^3`$ (Eq.(47)), the Poisson bispectrum rapidly becomes to dominate the total bispectrum in small angular scales,
$`{\displaystyle \frac{b_{lll}^{primary}}{b^{ps}}}`$ $``$ $`l^4\times 10^7\left({\displaystyle \frac{f_{NL}}{b^{ps}/10^{25}}}\right),`$ (58)
$`{\displaystyle \frac{b_{lll}^{szlens}}{b^{ps}}}`$ $``$ $`l^3\times 10^6\left({\displaystyle \frac{j_\nu \overline{T}_{\rho 0}b_{gas}}{b^{ps}/10^{25}}}\right).`$ (59)
For example, the SZ–lensing bispectrum measured by MAP experiment is overwhelmed by point sources at $`l100`$.
## V Measuring Bispectra
### A Fisher matrix
We shall discuss the detectability of CMB experiments to the primary non-Gaussianity in the bispectrum. We also need to separate it from secondary bispectra. Suppose that we try to fit the observed bispectrum $`B_{l_1l_2l_3}^{obs}`$ by theoretically calculated bispectra which include both primary and secondary sources. Then we minimize $`\chi ^2`$ defined by
$$\chi ^2\underset{2l_1l_2l_3}{}\frac{\left(B_{l_1l_2l_3}^{obs}_iA_iB_{l_1l_2l_3}^{(i)}\right)^2}{\sigma _{l_1l_2l_3}^2},$$
(60)
where $`i`$ denotes a component such as the primary, the SZ and lensing effects, extragalactic sources, and so on. Unobservable modes $`l=0`$ and $`1`$ are removed. In case that the non-Gaussianity is small, the cosmic variance of the bispectrum is given by the six-point function of $`a_{lm}`$. The variance of $`B_{l_1l_2l_3}`$ is then calculated as
$$\sigma _{l_1l_2l_3}^2B_{l_1l_2l_3}^2B_{l_1l_2l_3}^2𝒞_{l_1}𝒞_{l_2}𝒞_{l_3}\mathrm{\Delta }_{l_1l_2l_3},$$
(61)
where $`\mathrm{\Delta }_{l_1l_2l_3}`$ takes values 1, 2, and 6 for cases of that all $`l`$’s are different, two of them are same, and all are same, respectively. $`𝒞_lC_l+C_l^N`$ is the total CMB angular power spectrum, which includes the power spectrum of the detector noise $`C_l^N`$. $`C_l^N`$ is calculated analytically using the formula derived by Knox with the noise characteristics of the relevant experiments. We do not include $`C_l`$ from secondary sources, as they are totally subdominant compared with the primary $`C_l`$ and $`C_l^N`$ for relevant experiments. For example, inclusion of $`C_l`$ from extragalactic sources (Eqs.(54) or (55)) changes our results less than 10%.
Taking $`\chi ^2/A_i=0`$, we obtain the normal equation,
$$\underset{j}{}\left[\underset{2l_1l_2l_3}{}\frac{B_{l_1l_2l_3}^{(i)}B_{l_1l_2l_3}^{(j)}}{\sigma _{l_1l_2l_3}^2}\right]A_j=\underset{2l_1l_2l_3}{}\frac{B_{l_1l_2l_3}^{obs}B_{l_1l_2l_3}^{(i)}}{\sigma _{l_1l_2l_3}^2}.$$
(62)
Thus, we define the Fisher matrix $`F_{ij}`$ as
$$F_{ij}\underset{2l_1l_2l_3}{}\frac{B_{l_1l_2l_3}^{(i)}B_{l_1l_2l_3}^{(j)}}{\sigma _{l_1l_2l_3}^2}=\frac{2}{\pi }\underset{2l_1l_2l_3}{}\left(l_1+\frac{1}{2}\right)\left(l_2+\frac{1}{2}\right)\left(l_3+\frac{1}{2}\right)\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)^2\frac{b_{l_1l_2l_3}^{(i)}b_{l_1l_2l_3}^{(j)}}{\sigma _{l_1l_2l_3}^2},$$
(63)
where we have used equation (11) to replace $`B_{l_1l_2l_3}`$ by the reduced bispectrum $`b_{l_1l_2l_3}`$ (see Eq.(10) for definition). Since the covariance matrix of $`A_i`$ is $`F_{ij}^1`$, we define the signal-to-noise ratio $`(S/N)_i`$ for a component $`i`$, the correlation coefficient $`r_{ij}`$ between different components $`i`$ and $`j`$, and the degradation parameter $`d_i`$ of $`(S/N)_i`$ due to $`r_{ij}`$ as
$`\left({\displaystyle \frac{S}{N}}\right)_i`$ $``$ $`{\displaystyle \frac{1}{\sqrt{F_{ii}^1}}},`$ (64)
$`r_{ij}`$ $``$ $`{\displaystyle \frac{F_{ij}^1}{\sqrt{F_{ii}^1F_{jj}^1}}},`$ (65)
$`d_i`$ $``$ $`F_{ii}F_{ii}^1.`$ (66)
Note that $`r_{ij}`$ does not depend on amplitudes of bispectra, but shapes. $`d_i`$ is defined so as $`d_i=1`$ for zero degradation, while $`d_i>1`$ for degraded $`(S/N)_i`$. Spergel and Goldberg and Cooray and Hu considered the diagonal component of $`F_{ij}^1`$, while we study all components in order to discuss the separatability between various bispectra.
An order of magnitude estimation of $`S/N`$ as a function of a certain angular resolution $`l`$ is possible as follows. Since the number of modes contributing to $`S/N`$ increases as $`l^{3/2}`$ and $`l^3\left(\begin{array}{ccc}l& l& l\\ 0& 0& 0\end{array}\right)^20.36\times l`$, we estimate $`(S/N)_i(F_{ii})^{1/2}`$ as
$$\left(\frac{S}{N}\right)_i\frac{1}{3\pi }l^{3/2}\times l^{3/2}\left|\left(\begin{array}{ccc}l& l& l\\ 0& 0& 0\end{array}\right)\right|\times \frac{l^3b_{lll}^{(i)}}{(l^2C_l)^{3/2}}l^5b_{lll}^{(i)}\times 4\times 10^{12},$$
(67)
where we have used $`l^2C_l6\times 10^{10}`$.
Table I and II tabulate all components of $`F_{ij}`$ and $`F_{ij}^1`$, respectively. Table III summarizes $`(S/N)_i`$, while table IV tabulates $`d_i`$ in the diagonal, and $`r_{ij}`$ in the off-diagonal parts.
### B Measuring primary bispectrum
Figure 4 shows the numerical results of differential $`S/N`$ for the primary bispectrum at $`\mathrm{ln}l_3`$ interval, $`\left[d(S/N)^2/d\mathrm{ln}l_3\right]^{1/2}f_{NL}^1`$, in the upper panel, and $`(S/N)(<l_3)f_{NL}^1`$, which is $`S/N`$ summed up to a certain $`l_3`$, in the lower panel. The detector noises $`C_l^N`$ have been computed for COBE 4-yr map, for MAP 90 GHz channel, and for Planck 217 GHz channel, but the effect of limited sky coverage is neglected. Figure 4 also shows results for the ideal experiment with no noise: $`C_l^N=0`$. Both $`\left[d(S/N)^2/d\mathrm{ln}l_3\right]^{1/2}`$ and $`(S/N)(<l_3)`$ are monotonically increasing function with $`l_3`$ as roughly $`l_3`$ up to $`l_32000`$ for the ideal experiment.
Beyond $`l_32000`$, an enhancement of the damping tail in $`C_l`$ because of the weak lensing effect stops $`\left[d(S/N)^2/d\mathrm{ln}l_3\right]^{1/2}`$ and then $`(S/N)(<l_3)`$ increasing. This leads to an important constraint on the observation; even for the ideal noise-free and the infinitesimal thin-beam experiment, there is an upper limit on the value of $`S/N0.3f_{NL}`$. For a given realistic experiment, $`\left[d(S/N)^2/d\mathrm{ln}l_3\right]^{1/2}`$ has a maximum at a scale near the beam-size.
The total $`(S/N)f_{NL}^1`$ are $`1.7\times 10^3`$, $`5.8\times 10^2`$, and $`0.19`$ for COBE, MAP and Planck experiments, respectively (see table III). In order to obtain $`S/N>1`$, therefore, we need $`f_{NL}>600,20`$, and $`5`$ for each corresponding experiment, while the ideal experiment requires $`f_{NL}>3`$ (see table V). These values are also roughly obtained by substituting equation (26) into (67),
$$\left(\frac{S}{N}\right)_{primary}l\times 10^4f_{NL}.$$
(68)
The degradation parameters $`d_{primary}`$ are 1.46, 1.01, and 1.00 for COBE, MAP, and Planck experiments, respectively (see table IV). This means that MAP and Planck experiments will separate the primary bispectrum from others at 1% or better accuracies. Since the primary and other secondary sources change monotonically in the COBE angular scales, COBE cannot discriminate between them very well. In the MAP and Planck scales, however, the primary bispectrum starts oscillating around zero, and then is well separated in shape from other secondaries, as the secondaries do not oscillate. This is good news for the forthcoming high angular resolution CMB experiments.
### C Measuring secondary bispectra
The signal-to-noises for measuring the SZ–lensing bispectrum $`(S/N)_{szlens}`$ in units of $`\left|j_\nu \right|\overline{T}_{\rho 0}b_{gas}`$ are $`1.8\times 10^4`$, $`0.34`$, and $`6.2`$ for COBE, MAP, and Planck experiments, respectively (see table III). $`\overline{T}_{\rho 0}`$ is in units of 1 keV. Using equations (67) and (47), we roughly estimate $`(S/N)_{szlens}`$ as
$$\left(\frac{S}{N}\right)_{szlens}l^2\times 10^6\left|j_\nu \right|\overline{T}_{\rho 0}b_{gas}.$$
(69)
Thus, $`(S/N)_{szlens}`$ increases with the angular resolution more rapidly than the primary bispectrum (see Eq.(68)). Since $`\left|j_\nu \right|\overline{T}_{\rho 0}b_{gas}`$ should be an order of unity, COBE and MAP would not be expected to detect the SZ–lensing bispectrum; however, Planck would be sensitive enough to detect, depending on the frequency, i.e., a value of $`j_\nu `$. For example, 217 GHz is totally insensitive to the SZ effect as $`j_\nu 0`$, while $`j_\nu =2`$ in the Rayleigh–Jeans regime.
The degradation parameters $`d_{szlens}`$ are 3.89, 1.16, and 1.00 for COBE, MAP, and Planck experiments, respectively (see table IV). Thus, Planck will separate the SZ–lensing bispectrum from other effects. Note that $`(S/N)_{szlens}`$ values must be an order of magnitude estimation, as our cosmological model is the COBE normalized SCDM yielding $`\sigma _8=1.2`$. Since this $`\sigma _8`$ is about a factor of 2 greater than the cluster normalization with $`\mathrm{\Omega }_m=1`$, and $`20\%`$ greater than the normalization with $`\mathrm{\Omega }_m=0.3`$. Thus, this factor tends to overestimate $`\mathrm{\Theta }_{lm}^{}a_{lm}^{SZ}`$ (Eq.(46)) by a factor of several. On the other hand, using the linear power spectrum for $`P_\mathrm{\Phi }(k)`$ rather than the non-linear power spectrum tends to underestimate the effect by a factor of several at $`l3000`$. However, our main goal is to discriminate between shapes of various bispectra, not amplitudes, so that this factor does not affect our conclusion on the degradation parameters $`d_i`$.
For the extragalactic radio and infrared sources, we estimated the signal-to-noises as $`5.7\times 10^7(b^{ps}/10^{25})`$, $`2.2(b^{ps}/10^{25})`$, and $`52(b^{ps}/10^{27})`$ for COBE, MAP, and Planck experiments, respectively (see table III), and the degradation parameters $`d_{ps}`$ are 3.45, 1.14, and 1.00 (see table IV). Since
$$\left(\frac{S}{N}\right)_{ps}l^5\times 10^{13}\left(\frac{b^{ps}}{10^{25}}\right),$$
(70)
from equation (67), $`S/N`$ of the bispectrum from point sources increases very rapidly with the angular resolution. Our estimate that MAP will detect the bispectrum from point sources is consistent with the results found by Refregier, Spergel and Herbig. Although MAP cannot separate the Poisson bispectrum from the SZ–lensing bispectrum very well (see $`r_{ij}`$ in table IV), it would not matter as the SZ–lensing bispectrum would be too small to be measured by MAP. Planck will do an excellent job on separating all kinds of bispectra, at least including the primary signal, SZ–lensing coupling, and extragalactic point sources, on the basis of the shape difference.
### D Measuring primary skewness
For the skewness, we define $`S/N`$ as
$$\left(\frac{S}{N}\right)^2\frac{S_3^2}{\sigma _{S_3}^2},$$
(71)
where the variance is
$`\sigma _{S_3}^2`$ $``$ $`\left(S_3\right)^2=6{\displaystyle _1^1}{\displaystyle \frac{d\mathrm{cos}\theta }{2}}\left[𝒞(\theta )\right]^3`$ (72)
$`=`$ $`6{\displaystyle \underset{l_1l_2l_3}{}}{\displaystyle \frac{\left(2l_1+1\right)\left(2l_2+1\right)\left(2l_3+1\right)}{(4\pi )^3}}\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)^2𝒞_{l_1}𝒞_{l_2}𝒞_{l_3}W_{l_1}^2W_{l_2}^2W_{l_3}^2`$ (75)
$`=`$ $`{\displaystyle \frac{9}{2\pi ^3}}{\displaystyle \underset{2l_1l_2l_3}{}}\left(l_1+{\displaystyle \frac{1}{2}}\right)\left(l_2+{\displaystyle \frac{1}{2}}\right)\left(l_3+{\displaystyle \frac{1}{2}}\right)\left(\begin{array}{ccc}l_1& l_2& l_3\\ 0& 0& 0\end{array}\right)^2𝒞_{l_1}𝒞_{l_2}𝒞_{l_3}W_{l_1}^2W_{l_2}^2W_{l_3}^2.`$ (78)
In the last equality, we have used the symmetry of summed quantity with respect to indices (Eq.(37)), and removed unobservable modes $`l=0`$ and $`1`$. Typically $`\sigma _{S_3}10^{15}`$, as $`\sigma _{S_3}\left[𝒞(0)\right]^{3/2}10^{15}`$, where $`𝒞(\theta )`$ is the temperature auto correlation function including noise. The lower panel of figure 3 shows $`\sigma _{S_3}(<l_3)`$, which is $`\sigma _{S_3}(<l_3)`$ summed up to a certain $`l_3`$, for COBE, MAP, and Planck experiments as well as the ideal experiment. Since $`𝒞_lW_l^2=C_le^{l(l+1)\sigma _b^2}+w^1`$, where $`w^1`$ determines the white noise power spectrum of the detector noise according to the Knox’s formula, the dominance of second term beyond the experimental angular resolution scale, $`l\sigma _b^1`$, keeps $`\sigma _{S_3}(<l_3)`$ slightly increasing with $`l_3`$, while $`S_3(<l_3)`$ becomes constant beyond that (see the upper panel of figure 3). As a result, $`S/N`$ starts somewhat decreasing beyond the resolution. We use the maximum $`S/N`$ for estimating the minimum value of $`f_{NL}`$ needed to detect the primary $`S_3`$. We find that $`f_{NL}>800`$, 80, 70, and 60 for COBE, MAP, Planck, and the ideal experiments, respectively, with all-sky coverage.
These $`f_{NL}`$ values are systematically larger than those needed to detect $`B_{l_1l_2l_3}`$ by a factor of 1.3, 4, 14, and 20, respectively (see table V). Higher the angular resolution, less sensitive measuring the primary $`S_3`$ than $`B_{l_1l_2l_3}`$. This is because the cancellation effect in smaller angular scales due to the oscillation of $`B_{l_1l_2l_3}`$ damps $`S_3`$.
## VI Discussion and Conclusion
Using the full radiation transfer function , we have computed numerically the primary cosmic microwave background bispectrum (Eq.(22)) and skewness (Eq.(36)) down to arcminutes angular scales. The primary bispectrum oscillates around zero (figure 2), thus the primary skewness saturates at the MAP angular resolution scale, $`l500`$ (figure 3). We have introduced the “reduced” bispectrum defined by equation (10), and confirmed that this quantity is more useful to describe the physical property of the bispectrum than the full bispectrum (Eq.(2)).
Figure 5 compares the expected signal-to-noise ratio for detecting the primary non-Gaussianity based on the bispectrum (Eq.(64)) to that based on the skewness (Eq.(71)). It shows that the bispectrum is almost an order of magnitude more sensitive to the non-Gaussianity than the skewness. We conclude that when we can compute the predicted form of the bispectrum, it becomes a “matched filter” for detecting the non-Gaussianity in the data, and thus much more powerful tool than the skewness. Table V summarizes $`f_{NL}`$ required for detecting the primary non-Gaussianity using the bispectrum or the skewness with COBE, MAP, Planck, and the ideal experiments. This shows that even the ideal experiment needs $`f_{NL}>3`$ in order to detect the primary bispectrum.
We estimated the secondary bispectra from the coupling between the Sunyaev–Zel’dovich (SZ) and the weak lensing effects, and from the extragalactic radio and infrared sources. Only Planck will detect the SZ–lensing bispectrum, while both MAP and Planck will detect the bispectrum from extragalactic point sources (table III).
We also studied how well we can discriminate among the primary, the SZ–lensing coupling, and the extragalactic point sources bispectra. We found that MAP and Planck will separate the primary from other secondary sources at 1% or better accuracies. This conclusion is due to the presence of acoustic oscillation in the primary bispectrum that does not appear in the secondary bispectra. The SZ–lensing coupling and the extragalactic sources are well separately measured by Planck experiment, although COBE and MAP cannot discriminate between them (table IV).
Our arguments about the ability to discriminate among various bispectra were fully based upon the shape difference, and thus did not take into account the spectral difference in the frequency space. As pointed out by , the multi-band observation is so efficient to discriminate among the primary signal and the other foreground contaminants for measuring the CMB anisotropy power spectrum. Their scheme should be effective on the bispectrum as well, and the accuracy of the foreground removal will be improved further. Thus, we expect that MAP and Planck will measure the primary bispectrum separately from the foregrounds.
The simplest inflationary scenario usually predicts small $`f_{NL}`$ $`(10^2)`$, and the second order perturbation theory yields $`f_{NL}1`$. Thus, the significant detection of the primary bispectrum or the skewness with any experiments means that the simplest inflationary scenario needs to be modified. According to our results, if the reported detections of the bispectrum in the COBE map were the cosmological origin, then MAP and Planck would detect the primary bispectrum much more significantly. Although Banday, Zaroubi and Górski pointed out the one of those detections could be accounted for by the experimental systematic effects of COBE, the other is claimed to be significant even after removing such the systematics.
Although we have not discussed so far, the spatial distribution of emissions from interstellar dust is a potential source of the microwave non-Gaussianity. Since it is very hard to estimate the bispectrum analytically, the dust map compiled by Schlegel, Finkbeiner and Davis could be used to estimate the dust bispectrum. For example, we found that the dimensionless skewness parameter defined by $`(\mathrm{\Delta }T)^3/(\mathrm{\Delta }T)^2^{3/2}`$ is as large as 51. We used the publicly available HEALPix-formatted $`100\mu \mathrm{m}`$ map which contains 12,582,912 pixels without sky cut. The mean intensity in the map was $`14.8\mathrm{MJy}\mathrm{sr}^1`$. Of course, this skewness is largely an overestimate for the CMB measurement in reality; we need to cut a fraction of sky which contains the Galactic plane, and then this will greatly reduce the non-Gaussianity. Nevertheless, residual non-Gaussianity is still a source of the microwave bispectrum, and has to be taken into account. Moreover, the form of the bispectrum measured in the dust map would reflect the physics of interstellar dust, which is highly uncertain at present, and thus studying the interstellar dust bispectrum would be challenging field.
## Acknowledgments
We would like to thank Naoshi Sugiyama and Licia Verde for useful comments, and Uro$`\stackrel{ˇ}{\mathrm{s}}`$ Seljak and Matias Zaldarriaga for making their CMBFAST code publicly available. E. K. acknowledges a fellowship from the Japan Society for the Promotion of Science. D. N. S. is partially supported by the MAP/MIDEX program.
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# 1 Introduction
## 1 Introduction
Classical solutions of nonabelian gauge theories have played an important role in a variety of contexts. Classical solutions in Higgs theories may play an important role in cosmology. They may also be relevant in models of confinement. Different classical objects may affect cosmology, symmetry breaking, etc. in different ways. Therefore, it is of considerable importance to find all classical solutions and investigate their properties.
Vortex solutions are solitons in 2+1 dimensions and are stringlike extended objects in 3+1 dimenions. In 3+1 dimensions they have infinite energy (the energy per unit length is finite) but condensed vortices contribute a finite amount to the free energy per unit volume. Nonabelian vortex configurations were discussed in -; explict vortex solutions were first found in ref. . The existence of nonabelian vortices is the consequence of nontrivial topological classes in the mapping $`S_1SU(N)/Z_N.`$ The homotopy group of this mapping is $`Z_N`$, implying the existence of $`N1`$ distinct stable vortices. As the symmetry, classifying vortices, is the center of the gauge group $`SU(N)`$, one needs to introduce Higgs fields that break $`SU(N)`$ symmetry, but not the center $`Z_N`$. The smallest representation for the Higgs fields, such that they commute with the center, is the adjoint representation. Therefore, one needs to use one or more adjoint Higgs bosons to break the symmetry. Symmetry breaking induced by a single adjoint Higgs boson is not complete. The adjoint Higgs, when diagonalized, commutes with the ‘diagonal’ generators, the elements of the Cartan subgroup, $`[U(1)]^{N1}.`$ The relevant classical objects in such a theory are ’t Hooft-Polyakov monopoles. Thus, at least two adjoint Higgs bosons are needed to break the symmetry down to its center.
Vortex solutions found in correspond to $`SU(N)`$ adjoint Higgs theories with $`N`$ Higgs bosons. In fact, one would think that a ‘minimal’ solution could be found with only two Higgs bosons. The first Higgs boson breaks the symmetry down to the maximal abelian subgroup and then another Higgs, that is kept non-parallel with the first one, can break all the remaining continuous symmetries. The purpose of this paper is to show that vortex solutions in $`SU(3)`$ gauge theory with two adjoint Higgs bosons exist and to study the properties of these solutions.
The equations of motion in Abelian - and nonabelian vortex model were shown to reduce to linear, Bogomol’nyi equations at critical values of the coupling constant. This phenomenon was shown to be related to the increase of an underlying supersymmetry of the model. - The equations of motions we obtain for the $`SU(3)`$ Higgs theory also linearize and decouple at critical couplings. The relationship with increased supersymmetry can also be ascertained as the fields decouple into a couple of abelian vortices at the critical coupling.
In the next section we will briefly review the solutions of field equations for $`SU(3)`$ theory offered in Ref. . In section 3 we will present our two Higgs model and the ansatz for solving the equations of motion. In section 4 we will discuss the critical coupling, the Bogomol’nyi eqautions and their relationship to supersymmetry, followed by a concluding section.
## 2 Vortex solutions in $`SU(N)`$ gauge theory with $`N`$ Higgs
As usual in discussing time-independent classical solutions we will consider the Hamiltonian, the negative of the Lagrangian in the absence of time derivative. The Hamiltonian for a cylindrically symmetric solution is of the form
$$H=d^2x\left[\frac{1}{4}G_{\mu \nu }^2+\frac{1}{2}\underset{A=1}{\overset{N}{}}(D_\mu \mathrm{\Phi }^{(A)})^2+V(\mathrm{\Phi }^{(A)})\right],$$
(1)
Here
$`A_\mu `$ $`=`$ $`A_\mu ^at^a,a=1,2,\mathrm{},N^21`$
$`D_\mu `$ $`=`$ $`_\mu +ie[A_\mu ,]`$
$`G_{\mu \nu }`$ $`=`$ $`_\mu A_\nu _\nu A_\mu +ie[A_\mu ,A_\nu ]`$ (2)
where $`t^a`$ are the $`SU(N)`$ generators. We are considering $`N`$ Higgs scalars $`\mathrm{\Phi }^{(A)}`$ in the adjoint representation and the potential $`V[\mathrm{\Phi }]`$ chosen so as to ensure complete symmetry breaking.
Vortex solutions to the equations of motion associated with Hamiltonian (1) have been found in by making an ansatz that ensures non-trivial topology and maximum symmetry breaking. Since the scalars are in the adjoint representation, the center $`Z_N`$ of $`SU(N)`$ is the surviving symmetry subgroup. Then, the relevant homotopy group for classifying topologically inequivalent configurations is non-trivial, $`\pi _1(SU(N)/Z_N)=Z_N`$. One then has $`N1`$ topologically non-trivial inequivalent possible solutions which can be associated with gauge group elements $`U_n`$ ($`n=1,2,\mathrm{},N1`$ labeling the homotopy classes). If we call $`\varphi `$ the azimuthal angle in a plane perpendicular to the vortex, then $`U_n(\varphi )`$ should satisfy, when a turn around a closed contour is made,
$$U_n(\varphi +2\pi )=\mathrm{exp}\left(i\frac{2\pi (n+Nk)}{N}\right)U_n(\varphi ),n=1,2,\mathrm{},N1,kZ$$
(3)
Condition (3) can be realized by writing
$$U_n(\varphi )=\mathrm{diag}(\mathrm{exp}(i(n+Nk)\frac{\varphi }{N}),\mathrm{},\mathrm{exp}(i(n+Nk)\frac{\varphi }{N}),\mathrm{exp}(i\varphi \frac{N1}{N}(n+Nk)))$$
(4)
Then, one can construct a gauge field configuration $`A_\mu ^n`$ belonging to the $`n`$ class so that it satisfies, at infinity,
$$\underset{\rho \mathrm{}}{lim}A_\mu ^n=\frac{i}{e}U_n^{}(\varphi )_\varphi U_n(\varphi )_\mu \varphi =\frac{1}{e}M_n_\mu \varphi $$
(5)
One has explicitly
$$M_n=(n+Nk)\mathrm{diag}(\frac{1}{N},\frac{1}{N},\mathrm{},\frac{1}{N},\frac{1N}{N})$$
(6)
and hence $`M_n`$ can be written in terms of the $`(N1)`$ $`SU(N)`$ generators $`H_i`$ spanning the Cartan subalgebra of $`SU(N)`$,
$$M_n=(n+Nk)\underset{i=1}{\overset{N1}{}}m^iH_i$$
(7)
where $`m^i`$ are the magnetic weights, as defined in .
In view of (5), the natural ansatz for a vortex solution with topological charge $`n`$ is
$$A_\mu ^n=\frac{1}{e}_\mu \varphi M_na(\rho )$$
(8)
with $`a(\rho )`$ such that $`G_{\mu \nu }0`$ as $`\rho \mathrm{}`$, fast enough to ensure the finiteness of the energy.
The finiteness of energy also requires that, at infinity, the Higgs scalars $`\mathrm{\Phi }^{(A)},(A=1,2,\mathrm{},N)`$ take their vacuum value, minimizing the symmetry breaking potential. Moreover,
$$\underset{\rho \mathrm{}}{lim}D_\mu [\mathrm{\Phi }^{(A)}]=0$$
(9)
Condition (9) can be achieved either by taking the scalars in the Cartan algebra of $`SU(N)`$ or in its complement. Let us write the $`SU(N)`$ generators in the Cartan-Weyl basis, with $`H_i`$ the $`N1`$ generators spanning the Cartan algebra and $`E_{\pm \alpha }`$ those in its complement,
$`[H_i,E_{\pm \alpha }]`$ $`=`$ $`\pm \alpha _iE_{\pm \alpha }`$
$`[E_\alpha ,E_\alpha ]`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N1}{}}}\alpha _iH_i`$ (10)
where $`\alpha _i=\alpha ^i`$ are the roots in an orthonormal basis. Then, one can choose the symmetry breaking potential so that the first $`S`$ scalars $`\overline{\mathrm{\Phi }}^{(1)},\overline{\mathrm{\Phi }}^{(2)},\mathrm{}\overline{\mathrm{\Phi }}^{(S)}`$ are in the Cartan algebra and the rest, $`\mathrm{\Phi }^{(1)},\mathrm{\Phi }^{(2)},\mathrm{}\mathrm{\Phi }^{(T)}`$ in its complement, $`S+T=N`$. Now, in order to satisfy (9), one necessarily has
$`\underset{\rho \mathrm{}}{lim}\overline{\mathrm{\Phi }}^{(A)}(\rho ,\varphi )`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N1}{}}}C_j^{(A)}H_j`$
$`\underset{\rho \mathrm{}}{lim}\mathrm{\Phi }^{(A)}(\rho ,\varphi )`$ $`=`$ $`U_n^{}(\varphi )\left({\displaystyle \underset{\pm \alpha }{}}\eta _\alpha ^{(A)}E_\alpha \right)U_n(\varphi )=U_n^{}(\varphi )\eta ^{(A)}U_n(\varphi )`$ (11)
with $`C_j^{(A)}`$ and $`\eta _\alpha ^{(A)}`$ constants. The constants $`\eta ^{(A)}`$ should be adjusted to so that they would mimimize $`V(\mathrm{\Phi }^{(A)})`$.
In view of the conditions described above, a consistent ansatz for a $`Z_N`$ vortex configuration can be proposed in the form
$`\overline{\mathrm{\Phi }}^{(A)}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N1}{}}}C_j^{(A)}H_j`$
$`\mathrm{\Phi }^{(A)}`$ $`=`$ $`U_n^{}(\varphi )\left({\displaystyle \underset{\pm \alpha }{}}\eta _\alpha ^{(A)}\psi _\alpha ^{(A)}(\rho )E_\alpha \right)U_n(\varphi )`$
$`A_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{e}}a(\rho )M_n`$
$`A_\rho `$ $`=`$ $`A_0=A_z=0`$ (12)
Here we have taken the $`\overline{\mathrm{\Phi }}^{(A)}`$ scalars to be constant everywhere. $`F_\alpha ^{(A)}(\rho )`$ and $`a(\rho )`$ should satisfy the boundary conditions
$`\underset{\rho \mathrm{}}{lim}\psi _\alpha ^{(A)}(\rho )=1,`$ $`\underset{\rho \mathrm{}}{lim}a(\rho )=n.`$ (13)
Ansatz (12) implies that
$$D_\varphi \mathrm{\Phi }^{(A)}=(na(\rho ))_\varphi \mathrm{\Phi }^{(A)}$$
(14)
Given the ansatz for the $`\overline{\mathrm{\Phi }}`$-type scalars, the equations of motion derived from (1) take the form
$`D_\mu G^{\mu \nu }`$ $`=`$ $`ie{\displaystyle \underset{A=1}{\overset{N1}{}}}[D_\nu \mathrm{\Phi }^{(A)},\mathrm{\Phi }^{(A)}]`$
$`D_\mu D^\mu \mathrm{\Phi }^{(A)}`$ $`=`$ $`{\displaystyle \frac{\delta V}{\delta \mathrm{\Phi }^{(A)}}}`$ (15)
That is, appart from the potential, the $`\overline{\mathrm{\Phi }}^{(A)}`$ fields play no role in the dynamics. Concerning the other scalars $`\mathrm{\Phi }^{(A)}`$, separability of the equations of motion into radial and angular parts imposes
$`[M_n,[M_n,\mathrm{\Phi }^{(A)}]]`$ $`=`$ $`R_n^A(\rho )\mathrm{\Phi }^{(A)}`$
$`{\displaystyle \underset{i=1}{\overset{N1}{}}}[\mathrm{\Phi }^{(A)},[\mathrm{\Phi }^{(A)},M_n]]`$ $`=`$ $`S_n^A(\rho )M_n`$ (16)
One can see that these conditions simplify the ansatz (12) to
$`\overline{\mathrm{\Phi }}^{(A)}`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N1}{}}}C_j^{(A)}H_j`$
$`\mathrm{\Phi }^{(A)}`$ $`=`$ $`\eta ^{(A)}\psi ^{(A)}(\rho )U_n^{}(\varphi )\left(E_{\alpha _A}+E_{\alpha _A}\right)U_n(\varphi )`$
$`A_\varphi `$ $`=`$ $`{\displaystyle \frac{1}{e}}a(\rho )M_n`$
$`A_\rho `$ $`=`$ $`A_0=A_z=0`$ (17)
In order to characterize the vortex solutions from the topological point of view one can introduce an “electromagnetic tensor” $`𝒢_{\mu \nu }`$ analogous to that proposed by Polyakov for the $`SO(3)`$ monopole solution . In view of the ansatz for the gauge field, it is natural to take
$$𝒢_{\mu \nu }=\frac{\mathrm{tr}\left(M_nG_{\mu \nu }\right)}{\left(\mathrm{tr}\left(M_n^2\right)\right)^{1/2}}$$
(18)
Then, the flux $`\mathrm{\Phi }`$ associated to the magnetic field $`𝒢_{12}`$ reads, for the $`n`$-vortex solution
$$\mathrm{\Phi }=(n+Nk)\mathrm{\Phi }_0$$
(19)
with $`\mathrm{\Phi }_0=2\pi /e`$. Let us recall that $`n=1,2,\mathrm{},N1`$ indicates the topological class to which the configuration belongs while $`kZ`$ is related to gauge transformations that, although leading to the same behavior at infinity (and hence are topologically trivial), cannot be well defined everywhere and then are not gauge equivalent everywhere, thus giving, for a fixed $`n`$, different values for the magnetic flux .
Although the analysis of the radial equations of motion and their solution can be performed for arbitrary $`N`$, let us concentrate in the $`SU(3)`$ vortex solution, for which two topologically inequivalent classes exist. The associated $`U_n(\varphi )`$ are (we take for simplicity $`k=0`$)
$$U_n(\varphi )=e^{in\varphi \lambda _8/\sqrt{3}},n=1,2$$
(20)
One then has
$$M_n=\frac{n\lambda _8}{\sqrt{3}}$$
(21)
An explicit realization of the Cartan Algebra is
$$H_1=\frac{\lambda _3}{2},H_2=\frac{\lambda _8}{2}$$
(22)
where $`\lambda _3`$ and $`\lambda _8`$ are the diagonal Gell-Mann matrices. One then has, for the two-component magnetic weight (7)
$$\stackrel{}{m}=(0,2/\sqrt{3})$$
(23)
Concerning the step generators $`E_\alpha `$, they can be written in terms of the Gell-Mann matrices $`\lambda _i`$ in the form
$`E_{\alpha _1}+E_{\alpha _1}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\lambda _4`$
$`E_{\alpha _2}+E_{\alpha _2}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\lambda _6`$
$`E_{\alpha _3}+E_{\alpha _3}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\lambda _1`$ (24)
The solution found in corresponds to just one $`\overline{\mathrm{\Phi }}`$-type scalar,
$$\overline{\mathrm{\Phi }}=B\lambda _3+C\lambda _8$$
(25)
and two $`\mathrm{\Phi }`$-type ones
$`\mathrm{\Phi }^{(1)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}\eta ^{(1)}\psi ^{(1)}(\rho )U_n^{}(\varphi )\lambda _4U_n(\varphi )`$
$`\mathrm{\Phi }^{(2)}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}\eta ^{(2)}\psi ^{(2)}(\rho )U_n^{}(\varphi )\lambda _6U_n(\varphi )`$ (26)
With this choice, the radial equations of motion read
$`{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{d}{d\rho }}\left(\rho {\displaystyle \frac{d\psi ^{(A)}}{d\rho }}\right)\left({\displaystyle \frac{na(\rho )}{\rho }}\right)^2\psi ^{(A)}v^{(A)}(\rho )\psi ^{(A)}(\rho )`$ $`=`$ $`0`$
$`\rho {\displaystyle \frac{d}{d\rho }}\left({\displaystyle \frac{1}{\rho }}{\displaystyle \frac{da}{d\rho }}\right){\displaystyle \frac{e}{2}}\left(\left(\eta ^{(1)}\psi ^{(1)}\right)^2+\left(\eta ^{(2)}\psi ^{(2)}\right)^2\right)(na(\rho ))`$ $`=`$ $`0`$ (27)
where $`v^{(A)}(\rho )`$ stands for the derivative of the potential with respect to $`\mathrm{\Phi }^{(A)}`$.
The symmetry breaking potential proposed in can be written in the form
$`V(\mathrm{\Phi }^{(A)},\overline{\mathrm{\Phi }})`$ $`=`$ $`{\displaystyle \frac{g_1\eta ^4}{8}}{\displaystyle \underset{A=1}{\overset{2}{}}}\left({\displaystyle \frac{1}{2}}\mathrm{Tr}[\mathrm{\Phi }^{(A)}]^21\right)^2+{\displaystyle \frac{\overline{g}_1\eta ^4}{8}}\left({\displaystyle \frac{1}{2}}\mathrm{Tr}[\overline{\mathrm{\Phi }}\overline{\mathrm{\Phi }}]1\right)^2`$
$`+{\displaystyle \frac{g_2\eta ^4}{4}}\left(\mathrm{Tr}[\mathrm{\Phi }^{(1)}\mathrm{\Phi }^{(2)}]\right)^2+d_{abc}\overline{\mathrm{\Phi }}^a\left({\displaystyle \underset{A=1}{\overset{2}{}}}f_A\mathrm{\Phi }^{(A)b}\mathrm{\Phi }^{(A)c}+h\mathrm{\Phi }^{(1)b}\mathrm{\Phi }^{(2)c}\right)`$
where $`\mathrm{\Phi }^{(A)}=\mathrm{\Phi }^{(A)b}\lambda ^b`$ and $`d_{abc}`$ is the completely symmetric $`SU(3)`$ tensor. In (2) and in our subseqent discussion we will use Higgs fields that become normalized in the limit $`\rho \mathrm{}`$. One can see that the choice of the same coupling constant $`g_1`$ for the quartic coupling of the $`\mathrm{\Phi }^{(A)}`$ fields implies that $`f^{(1)}=f^{(2)}`$ and reduces system (27) to that arising in the $`U(1)`$ case, which is solved, at critical coupling, by the solutions of the original Bogomol’nyi equations.
## 3 A 2-Higgs vortex in $`SU(3)`$ gauge theory
In this section we shall present a ‘minimal’ $`SU(3)`$ solution with only two Higgs fields $`\mathrm{\Phi }^{(A)}`$ ($`A=1,2`$) in the adjoint. The Hamiltonian of the model is defined uniquely up to the Higgs potential. There is a considerable freedom in the Higgs potential. In a way we consider a Higgs potential simpler than that of the previous section, but in an other way we generalize it such that it will disallow solutions of the form discussed in Ref. . Vortex solutions for a similar generalization of the $`SU(2)`$ Higgs potential were shown to exist in Ref.
The Higgs potential we propose is identical in form to that of Ref. for $`SU(2)`$:
$$V(\mathrm{\Phi }^{(1)},\mathrm{\Phi }^{(2)})=\frac{g_1\eta ^4}{8}\underset{A=1}{\overset{2}{}}\left(\frac{1}{2}\mathrm{Tr}[\mathrm{\Phi }^{(A)}]^21\right)^2+\frac{g_2\eta ^4}{4}\left(\frac{1}{2}\mathrm{Tr}[\mathrm{\Phi }^{(1)}\mathrm{\Phi }^{(2)}]c\right)^2.$$
(29)
The generalization compared to Ref. appears in the nonzero value of the constant $`c`$, that is the cosine of the ‘angle’ between the two Higgs fields at infinity. The brackets of (29) must vanish at infinity to keep the Hamiltonian finite. Thus, unlike in previous models the Higgs fields are required not to be orthogonal at infinity. Admittedly, the model we study here is less general in the sense that the self coupling of the two Higgs bosons is assumed to be identical.
The field equations derived from the Lagrangian, analogous to (15) and (14), are
$$D_\mu G_{\mu \nu }ie\eta ^2\underset{A=1}{\overset{2}{}}[\mathrm{\Phi }^{(A)},D_\nu \mathrm{\Phi }^{(A)}]=0$$
(30)
and
$$D_\mu D_\mu \mathrm{\Phi }^{(A)}g_1\eta ^2\mathrm{\Phi }^{(A)}\left(\frac{1}{2}\mathrm{Tr}[\mathrm{\Phi }^{(A)}]^21\right)g_2\eta ^2\mathrm{\Phi }^{(B)}\left(\frac{1}{2}\mathrm{Tr}[\mathrm{\Phi }^{(1)}\mathrm{\Phi }^{(2)}]c\right)=0,$$
(31)
where $`A=1,2`$ and then $`B=2,1`$.
As in the previous section, the ansatz we use for finding vortex solutions is based on the philosophy that the vortex solution is associated with a singular gauge transformation that maps circles linked with the vortex to a smooth transformation connecting two elements of the center. Choosing $`U_n(\varphi )`$, as in (20), the Higgs fields are defined as
$$\mathrm{\Phi }^{(i)}(x)=U_n(\varphi )\psi ^{(i)}(\rho )U_n^{}(\varphi ),$$
(32)
where $`i=1,2`$ for the two Higgs bosons.
The ansatz for the gauge field,
$$A_\mu (x)=_\mu \varphi [a_8(\rho )\lambda _8+a_3(\rho )\lambda _3],$$
(33)
is diagonal in gauge space. We will later show that unlike for vortices of the previous section the component $`a_3(\rho )`$ must be different from zero, despite the fact that this component does not contribute to the vortex at $`\rho \mathrm{}`$. The gauge field of (33) satisfies the gauge fixing condition $`_\mu A_\mu =0`$. Taking the derivative of the Higgs field generates a vortex contribution in the $`\lambda _8`$ gauge direction. The form of the gauge field was chosen to be able to cancel this vortex at infinity in the covariant derivative. Without such a cancellation the term of the Hamiltonian containing the covariant derivative of the Higgs fields would diverge.
We still need to show that the forms chosen for the fields are consistent with field equations (30) and (31). Before doing so we will further restrict the form of our solution. We will assume that the Higgs fields have only components
$$\psi ^{(A)}=\psi _4^{(A)}\lambda _4+\psi _6^{(A)}\lambda _6,$$
(34)
where $`\lambda _4`$ and $`\lambda _6`$ are off diagonal Gell-Mann matrices matrices. Two is the minimal number of components needed to satisfy the all the constraints on the normalization of the Higgs fields at $`\rho \mathrm{}`$ simultaneously. The two Higgs fields, provided their coefficients are not identical, break $`SU(3)`$ symmetry completely, down to its center, $`Z_3`$.
Let us now show that the gauge structure we propose is consistent with the field equations. First of all consider (31). The two equations, for the choices $`A=1`$ and 2, are consistent with the solution $`\psi _I^{(1)}=\pm \psi _I^{(2)}`$. We will show that the choice
$$\psi _4^{(1)}=\psi _4^{(2)}\psi _4,\psi _6^{(1)}=\psi _6^{(2)}\psi _6$$
(35)
is also consistent with (30). Under the assumptions (32)-(35) (30) can be calculated as
$`D_\mu G_{\mu \nu }ie\eta ^2{\displaystyle \underset{A=1}{\overset{2}{}}}[\mathrm{\Phi }^{(A)},D_\nu \mathrm{\Phi }^{(A)}]=_\mu \varphi \left[\lambda _8\rho {\displaystyle \frac{d}{d\rho }}\left({\displaystyle \frac{1}{\rho }}{\displaystyle \frac{da_8}{d\rho }}\right)+\lambda _3\rho {\displaystyle \frac{d}{d\rho }}\left({\displaystyle \frac{1}{\rho }}{\displaystyle \frac{da_3}{d\rho }}\right)\right]`$ (36)
$``$ $`2e\eta ^2_\mu \varphi [(\psi _4)^2ea_+(\sqrt{3}\lambda _8+\lambda _3)+(\psi _6)^2ea_{}(\sqrt{3}\lambda _8\lambda _3)]=0,`$
where
$$a_\pm =\sqrt{3}a_8+\frac{n}{e}\pm a_3.$$
(37)
Clearly the space and isospace structures are consistent and (36) leads to two scalar equations for the two unknown functions, $`a_+`$ and $`a_{}`$. These equations are
$$\rho \frac{d}{d\rho }\left(\frac{1}{\rho }\frac{da_+}{d\rho }\right)4e^2\eta ^2[2(\psi _4)^2a_++(\psi _6)^2a_{}]=0,$$
(38)
and
$$\rho \frac{d}{d\rho }\left(\frac{1}{\rho }\frac{da_{}}{d\rho }\right)4e^2\eta ^2[(\psi _4)^2a_++2(\psi _6)^2a_{}]=0,$$
(39)
In a similar way, the scalar equations reduce to two equations for the two components, $`\psi _4`$ and $`\psi _6`$
$$\frac{1}{\rho }\frac{d}{d\rho }\left(\rho \frac{d\psi _4}{d\rho }\right)\frac{a_+^2}{\rho ^2}\psi _4g_1\eta ^2\psi _4(\psi _4^2+\psi _6^21)g_2\eta ^2\psi _4(\psi _4^2\psi _6^2c)=0,$$
(40)
and
$$\frac{1}{\rho }\frac{d}{d\rho }\left(\rho \frac{d\psi _6}{d\rho }\right)\frac{a_{}^2}{\rho ^2}\psi _6g_1\eta ^2\psi _6(\psi _4^2+\psi _6^21)+g_2\eta ^2\psi _6(\psi _4^2\psi _6^2c)=0.$$
(41)
The boundary conditions for the four fields are the following:
$`a_\pm (0)={\displaystyle \frac{n}{e}}`$
$`\underset{\rho \mathrm{}}{lim}a_\pm (\rho )=0`$
$`\psi _4(0)=\psi _6(0)=0`$
$`\underset{\rho \mathrm{}}{lim}\psi _4(\rho )=\sqrt{{\displaystyle \frac{1+c}{2}}},\underset{\rho \mathrm{}}{lim}\psi _6(\rho )=\sqrt{{\displaystyle \frac{1c}{2}}}.`$ (42)
Now at this point it should be obvious that $`a_3=0`$, equivalent to $`a_+=a_{}`$ is not an admissible solution. If $`a_+=a_{}`$ then from (38) and (39) it follows that $`\psi _4=\psi _6`$. Such a solution would not satisfy the boundary condition (42), unless $`c=0`$.
Note that at $`c=0`$ $`\psi _4=\psi _6`$ and $`a_+=a_{}`$. In other words the $`a_3`$ component of the gauge field vanishes. Then, after appropriate rescaling, the vortex defined by (38)-(41) coincides with that defined by (27), provided we choose $`g_1=\overline{g}_1`$ and $`\eta ^{(1)}=\eta ^{(2)}.`$
We have not been able to prove analytically the existence of solutions of these equations. In a future publication we will study the solutions numerically. At special values of the couplings, however, the second order equations become first order. The system of equations also decouples and can be rescaled to a form identical to a combination of two critical abelian vortices. Abelian vortices have been well studied and the existense of solutions has been shown.
The form of the solutions for two-adjoint-Higgs model is unique up to gauge transformations. A gauge transformation can always bring $`U_n(\varphi )`$ to the form used above. Then the gauge field, commuting with $`U_n(\varphi )`$ should only have components $`a_8`$ or $`a_3`$. Furthermore, the combination of constraint
$$\underset{A}{}[\mathrm{\Phi }^{(A)},_\mu \mathrm{\Phi }_{(A)}]=0$$
and of the field equations for the two Higgs fields can only be satisfied with at most two nonvanishing components of $`\mathrm{\Phi }^{(A)}`$. Choosing these as $`\mathrm{\Phi }_4`$ and $`\mathrm{\Phi }_6`$ we arrive at the choice of this section.<sup>1</sup><sup>1</sup>1Components that can be transformed into each other by a global $`U(1)\times U(1)`$ transformations and therefore satisfy the same field equations are not counted as different. For our choice of the components $`\mathrm{\Phi }_5`$ and $`\mathrm{\Phi }_7`$ can be eliminated by global $`U(1)\times U(1)`$ transformations.
## 4 Critical coupling
At a critical coupling the second order differential equations for the gauge and Higgs field of abelian vortex solutions can be transformed to linear equations ,. We will show below that the solution found in the previous section also satisfies linear equations at a critical coupling. Furthermore, we will also observe that the first order equations decouple into equations coupling the gauge field $`a_+`$ with $`\psi _4`$ and the gauge field $`a_{}`$ with the Higgs field $`\psi _6`$ only.
First of all, it will be advantageus to express Hamiltonian (1) in terms of the Higgs component $`a_8`$, $`a_3`$ (or $`a_+`$ and $`a_{}`$), $`\psi _4`$, and $`\psi _6`$. One obtains
$`H`$ $`=`$ $`2\pi {\displaystyle _0^{\mathrm{}}}\rho d\rho [{\displaystyle \frac{1}{2\rho ^4}}(\rho a_8^{}a_8)^2+{\displaystyle \frac{1}{2\rho ^4}}(\rho a_3^{}a_3)^2+{\displaystyle \frac{\eta ^2}{\rho ^2}}(\psi _4^2a_+^2+\psi _6^2a_{}^2)`$ (43)
$`+`$ $`\eta ^2(\psi _4^2+\psi _6^2)+{\displaystyle \frac{g_1\eta ^4}{4}}(\psi _4^2+\psi _6^21)^2+{\displaystyle \frac{g_2\eta ^4}{4}}(\psi _4^2\psi _6^2c)^2]`$
The variation of (43) results in field equations (38)-(41).
Inspired by Ref. we write rearrange the Hamiltonian into an alternative form
$`H`$ $`=`$ $`2\pi {\displaystyle _0^{\mathrm{}}}\rho d\rho \{\eta ^2[\psi _{4}^{}{}_{}{}^{}+\gamma \rho \psi _4w_+]^2+\eta ^2[\psi _6^{}+\delta \rho \psi _6w_{}]^2`$ (44)
$`+`$ $`{\displaystyle \frac{1}{2\rho ^4}}[\rho a_8^{}a_8\alpha \rho ^2\eta ^2(\psi _4^2+\psi _6^21)]^2+{\displaystyle \frac{1}{2\rho ^4}}[\rho a_3^{}a_3\beta \rho ^2\eta ^2(\psi _4^2\psi _6^2c)]^2`$
$`+`$ $`{\displaystyle \frac{f_1\eta ^4}{4}}(\psi _4^2+\psi _6^21)^2+{\displaystyle \frac{f_2\eta ^4}{4}}(\psi _4^2\psi _6^2c)^2+{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{dX}{d\rho }}\},`$
where $`\alpha `$, $`\beta `$, $`\gamma `$, and $`\delta `$ are yet undetermined constants and $`X`$ is an undetermined form. Comparing (44) with (43) provides the following values for the constants:
$$\gamma =\delta =\frac{n}{|n|}$$
(45)
$$\alpha =\sqrt{3}\beta =2\sqrt{3}e\frac{n}{|n|}.$$
(46)
Furthermore, one obtains
$$X=|n|(\psi _4^2+\psi _6^2),$$
(47)
$$f_1=g_124e^2,$$
(48)
and
$$f_2=g_28e^2,$$
(49)
Substituting these values back into the Hamiltonian we can see that the Hamiltonian is minimized with a minimum value of $`2\pi |n|`$ at the critical couplings
$$g_1=24e^2,$$
(50)
and
$$g_2=8e^2,$$
(51)
if the fields satisfy the following Bogomol’nyi type equations:
$$\psi _4^{}=\frac{e}{\rho }\psi _4a_+,$$
(52)
$$\psi _6^{}=\frac{e}{\rho }\psi _6a_{},$$
(53)
$$a_8^{}\frac{1}{\rho }a_8=\frac{\eta ^2\rho }{2\sqrt{3}e}(\psi _4^2+\psi _6^21),$$
(54)
and
$$a_3^{}\frac{1}{\rho }a_3=\frac{\eta ^2\rho }{2e}(\psi _4^2\psi _6^2c).$$
(55)
In fact, appropriate linear combinations of (54) and (55) can be taken to arrive at the equations
$$a_+^{}\frac{1}{\rho }a_+=\frac{\eta ^2\rho }{e}\left(\psi _4^2\frac{1+c}{2}\right),$$
(56)
and
$$a_{}^{}\frac{1}{\rho }a_{}=\frac{\eta ^2\rho }{e}\left(\psi _6^2\frac{1c}{2}\right),$$
(57)
Now field equation (52) and (56) decouple from (53) and (57). Both of these systems are identical to the systems one obtains for the gauge and Higgs fields at critical coupling for $`SU(2)`$ vortices. This is not an accident, as a gauge rotation can transform $`\psi _4`$ (or $`\psi _6`$) into $`\psi _1`$ and $`a_+`$ (or $`a_{}`$) into $`a_3`$ simultaneously. In the next section we will connect the $`SU(2)`$ vortices with abelian vortices at critical coupling.
## 5 Relation to Abelian vortices and supersymmetry
One can rewrite eqs.(52) and (56) so that they coincide with the Bogomol’nyi equations of an Abelian Higgs model where $`\psi _4`$ is identified with the modulus of the Higgs scalar and $`a_+`$ identified with the $`A_\varphi `$ component of the Abelian gauge field. Indeed, if one calls
$`f`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}F}{e}}\psi _4`$
$`a_+{\displaystyle \frac{n}{e}}`$ $``$ $`\sqrt{3}a_8+a_3=A`$ (58)
then, eqs.(52) and (56) become
$`f^{}(\rho )`$ $`=`$ $`(n+eA){\displaystyle \frac{f}{\rho }}`$
$`{\displaystyle \frac{1}{\rho }}{\displaystyle \frac{dA_\phi }{d\rho }}`$ $`=`$ $`{\displaystyle \frac{e}{2}}\left(f^2{\displaystyle \frac{F^2}{e^2}}(1+c)\right)`$ (59)
These are nothing but the first order (BPS) equations, as originally written in (see eqs. (3.5) and (3.6) in that paper), once an axially symmetric ansatz is imposed in the form
$`\mathrm{\Phi }`$ $`=`$ $`f(\rho )\mathrm{exp}(in\varphi )`$
$`A_\phi `$ $`=`$ $`{\displaystyle \frac{1}{\rho }}A(\rho )`$ (60)
Here we have called $`\mathrm{\Phi }`$ the complex scalar field and $`A_\mu `$ the $`U(1)`$ gauge field. The condition for the Higgs field at infinity corresponds to
$$\underset{\rho \mathrm{}}{lim}f(\rho )=\frac{F}{e}\sqrt{1+c}$$
(61)
The solution to eqs. (59) corresponds to a vortex with magnetic flux $`\mathrm{\Phi }=(2n\pi )/e`$. First order equations (59) solve the second order Euler-Lagrange equations for the Abelian Higgs model with symmetry breaking potential
$$V_{U(1)}=\frac{\lambda }{8}\left(|\varphi |^2|\varphi _0|\right)^2$$
(62)
provided the $`\varphi ^4`$ coupling constant $`\lambda `$ is chosen as ,
$$\lambda =e^2$$
(63)
We then see that eqs.(59) correspond to the Bogomol’nyi equations for a vortex with topological charge $`1`$. Of course, the $`+1`$ topological charge equation is also obtainable just by changing eq.(58) to
$`f`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}F}{e}}\psi _4`$
$`a_+`$ $`=`$ $`{\displaystyle \frac{n}{e}}+A`$ (64)
The same can be done for magnetic flux $`n=\pm 2`$ , etc. An analogous identification can be done concerning $`\psi _6`$ and $`w_{}`$. One gets the same equations as in (59) with $`cc`$.
The connection between BPS relations and supersymmetry has been thoroughfully analysed for the Abelian Higgs model, including the case in which a Chern-Simons term is added to the Maxwell term -. The outcome is that in order to achieve the $`N=2`$ supersymmetric extension of the purely bosonic model, one is forced to impose the condition (63), exactly as it happens when trying to find a BPS bound proceeding à la Bogomol’nyi . In the supersymmetric framework, the bound for the energy coincides with the central charge of the $`N=2`$ SUSY algebra, which can be seen to coincide with the magnetic flux, related to the topological charge.
Once the connection between the non-Abelian $`SU(3)`$ model presented in Section 3 and the $`U(1)`$ model is established, the supersymmetric analysis can be done in a very simple way. Indeed, since the first order system (52)-(57) decouples into two systems, one for $`(w_+,\psi _4)`$ and the other for $`(w_{},\psi _6)`$, one can analyse them separately. Let us consider for example the former system. Identification (58) implies that the $`U(1)`$ Lagrangian from which eqs.(52) and (56) can be derived takes the form
$$L_{U(1)}=\frac{1}{4}F_{\mu \nu }F^{\mu \nu }+\frac{1}{2}(_\mu \mathrm{\Phi }^{}ieA_\mu \mathrm{\Phi }^{})(_\mu \mathrm{\Phi }+ieA_\mu \mathrm{\Phi })\frac{e^2}{8}\left(|\mathrm{\Phi }|^2|\mathrm{\Phi }_0|\right)^2$$
(65)
In view of the axial symmetry of the problem (no $`x^3`$ dependence), one should consider $`\mu =0,1,2`$; that is, one is effectively working in $`2+1`$ dimensional space-time. In (65) $`A_\mu `$ and $`\mathrm{\Phi }`$ are connected with $`w_+`$ and $`\psi _4`$ according to eqs.(58),(60).
The $`N=2`$ supersymmetric extension of the model defined by Lagrangian (65) can be written in the form
$`L_{N=2}`$ $`=`$ $`\{{\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }+{\displaystyle \frac{1}{2}}(_\mu M)(^\mu M)+{\displaystyle \frac{1}{2}}(D_\mu \mathrm{\Phi })^{}(D^\mu \mathrm{\Phi }){\displaystyle \frac{e^2}{4}}M^2|\mathrm{\Phi }|^2`$ (66)
$``$ $`{\displaystyle \frac{e^2}{8}}(|\mathrm{\Phi }|^2\mathrm{\Phi }_{0}^{}{}_{}{}^{2})^2+{\displaystyle \frac{i}{2}}\overline{\mathrm{\Sigma }}\overline{)}\mathrm{\Sigma }+{\displaystyle \frac{i}{2}}\overline{\psi }\overline{)}D\psi {\displaystyle \frac{e}{2}}M\overline{\psi }\psi `$
$``$ $`{\displaystyle \frac{e}{2}}(\overline{\psi }\mathrm{\Sigma }\mathrm{\Phi }+h.c.)`$
Here $`M`$ is a real scalar, $`\psi `$ and $`\mathrm{\Sigma }`$ ($`\mathrm{\Sigma }=\chi +i\xi `$) being Dirac fermions.
Lagrangian (66) is invariant under the supersymmetric transformations
$`\widehat{\delta }\mathrm{\Sigma }`$ $`=`$ $`\left({\displaystyle \frac{1}{2}}ϵ^{\mu \nu \lambda }F_{\mu \nu }\gamma _\lambda +{\displaystyle \frac{e}{2}}(|\mathrm{\Phi }|^2\mathrm{\Phi }_{0}^{}{}_{}{}^{2})+i\overline{)}M\right)\eta _c,\widehat{\delta }A_\mu =i\overline{\eta }_c\gamma _\mu \xi `$
$`\widehat{\delta }\psi `$ $`=`$ $`i\gamma ^\mu D_\mu \mathrm{\Phi }\eta _c(e^2)^{1/2}M\mathrm{\Phi }\eta _c,\widehat{\delta }M=\overline{\eta }_c\chi ,\widehat{\delta }\mathrm{\Phi }=\overline{\eta }_c\psi `$ (67)
The spinor supercharges generating these transformations can be shown to be
$`Q`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{e\mathrm{\Phi }_0}}{\displaystyle }d^2x[({\displaystyle \frac{1}{2}}ϵ^{\mu \nu \lambda }F_{\mu \nu }\gamma _\lambda +i\overline{)}M{\displaystyle \frac{e}{2}}(|\mathrm{\Phi }|^2\mathrm{\Phi }_{0}^{}{}_{}{}^{2}))\gamma ^0\mathrm{\Sigma }`$ (68)
$`+`$ $`(i(\overline{)}D\mathrm{\Phi })^{}{\displaystyle \frac{e}{2}}M\mathrm{\Phi }^{})\gamma ^0\psi ]`$
and
$`\overline{Q}`$ $`=`$ $`{\displaystyle \frac{\sqrt{2}}{e\mathrm{\Phi }_0}}{\displaystyle }d^2x[\overline{\mathrm{\Sigma }}\gamma ^0({\displaystyle \frac{1}{2}}ϵ^{\mu \nu \lambda }F_{\mu \nu }\gamma _\lambda i\overline{)}M{\displaystyle \frac{e}{2}}(|\mathrm{\Phi }|^2\mathrm{\Phi }_{0}^{}{}_{}{}^{2}))`$ (69)
$`+`$ $`\overline{\psi }\gamma ^0(i\overline{)}D\mathrm{\Phi }{\displaystyle \frac{e}{2}}M\mathrm{\Phi })]`$
and satisfy the $`N=2`$ algebra
$$\{Q_\alpha ,\overline{Q}^\beta \}=2(\gamma _0)_{\alpha }^{}{}_{}{}^{\beta }P^0+\delta _{\alpha }^{}{}_{}{}^{\beta }Z$$
(70)
where $`\alpha ,\beta =1,2`$ and
$$P^0=E=\frac{1}{2e^2\varphi _{0}^{}{}_{}{}^{2}}d^2x\left[\frac{1}{2}F_{ij}^2+|D_i\mathrm{\Phi }|^2+\frac{e^2}{4}(|\mathrm{\Phi }|^2\mathrm{\Phi }_{0}^{}{}_{}{}^{2})^2\right]$$
(71)
while the central charge is given by:
$$Z=\frac{1}{e^2\varphi _{0}^{}{}_{}{}^{2}}d^2x\left[\frac{e}{2}ϵ^{ij}F_{ij}(|\mathrm{\Phi }|^2\mathrm{\Phi }_{0}^{}{}_{}{}^{2})+iϵ^{ij}(D_i\mathrm{\Phi })(D_j\mathrm{\Phi })^{}\right]$$
(72)
Here we have considered static configurations with $`A_0=0`$ so that $`i,j=1,2`$. Moreover, we have put $`M`$ and all fermions to zero to restrict the supersymmetric model to the original $`U(1)`$ model. One can easily see that the central charge, as given by (72), coincides with the magnetic flux,
$$Z=_i\left(\frac{1}{e}A_j+\frac{i}{e^2\mathrm{\Phi }_{0}^{}{}_{}{}^{2}}\mathrm{\Phi }^{}D_j\varphi \right)ϵ^{ij}=\frac{2\pi }{e}n$$
(73)
It is now easy to find the Bogomol’nyi bound from the supersymmetry algebra (70). Indeed, since the anticommutators in (70) are Hermitian, one has:
$$\{Q_\alpha ,\overline{Q}^\beta \}\{Q^\alpha ,\overline{Q}_\beta \}0$$
(74)
or using (70),
$$E|Z|$$
(75)
In order to explicitly obtain Bogomol’nyi equations (saturating the energy bound) from the supersymmetry algebra, one considers
$$Q_I=\frac{Q_++iQ_{}}{\sqrt{2}}$$
(76)
$$Q_{II}=\frac{\overline{Q}^++i\overline{Q}^{}}{\sqrt{2}}$$
(77)
where we have defined $`Q_\pm `$ from
$$Q=\left(\begin{array}{c}Q_+\\ Q_{}\end{array}\right)$$
(78)
$$\overline{Q}=\left(\overline{Q}^+\overline{Q}^{}\right)$$
(79)
Now, suppose that a field configuration $`|B`$ saturates the Bogomol’nyi bound derived from (74). Then, one necessarily has
$$\left(Q_I\pm Q_{II}\right)|B=0$$
(80)
or, using (76)-(79) and (68)-(69)
$`ϵ^{ij}F_{ij}`$ $`=`$ $`\pm e(|\mathrm{\Phi }|^2\mathrm{\Phi }_{0}^{}{}_{}{}^{2})`$
$`iϵ_{ij}D^i\mathrm{\Phi }`$ $`=`$ $`\pm (D_j\mathrm{\Phi })^{}`$ (81)
which are nothing but the equations (59) once the axially symmetric ansatz (60) is imposed.
## 6 Conclusions
A new vortex solution was shown to exist in $`SU(3)`$ gauge theory with two adjoint Higgs bosons. This can be contrasted with the the solution found in Ref. that requires three adjoint Higgs bosons. At a critical value of the Higgs self-coupling (where the gauge and Higgs masses coincide) the Hamiltonian has an exact lower bound and the Higgs and gauge fields satisfy first order Bogomol’nyi type field equations. The field equations for two Higgs and two gauge compenents also decouple at the critical couplings and both of the decoupled sets are equivalent to an $`SU(2)`$ vortex model at critical copupling. That model, as it has been shown here, is equivalent to an Abelian Higgs model at critical coupling. Thus, the critical $`SU(3)`$ model is ultimately equivalent to a pair of critical Abelian Higgs models. This relationship connects our models to supersymmetry. The supersymetric version of our model implies that the vortex mass per unit length is bounded by the $`N=2`$ SUSY central charge, which, at the same time equals to the magnetic flux of the vortex. In this respect, we expect that the non Abelian vortices discussed here could play a relevant role in the confinement scenario arising in strongly coupled supersymmetric theories .
Acknowledments
F.A.S will like to thank J. Edelstein, G. Lozano and C. Núñez for discussions and helpfull comments. F.A.S is partially supported by CICBA as Investigador, and through grants CONICET (PIP 4330/96), and ANPCYT (PICT 97/2285). P.S. is supported in part by the U.S. Department of Energy through grant #DE FG02-84ER-40153.
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# 1 INTRODUCTION
## 1 INTRODUCTION
There is overwhelming evidence that perturbation series in quantum field theory and related disciplines diverge (see for example and references therein). Consequently, resummation techniques, which allow to associate a finite value to a divergent series, are needed if divergent perturbation series are to be used for numerical purposes. The best known and most often applied resummation techniques are Padé approximants and the Borel method . Recently, another summation method – the so-called delta transformation (Eq. (8.4-4) of ) – has gained some prominence as a summation method in various domains of physics . There is evidence that the delta transformation is able to sum divergent series whose coefficients $`c_n`$ grow essentially like $`n!`$, $`(2n)!`$ and even $`(3n)!`$. This is not achievable by employing Padé approximants .
In the case of both Padé approximants and the delta transformation, only the numerical values of a finite number of partial sums of the divergent input series are needed as input data. Within the framework of the Borel method, it is also necessary to know the large order asymptotics of the coefficients of the divergent power series. Consequently, the Borel method is slightly less general than the other two methods.
Rational approximations to (divergent) power series have an interesting feature: It is possible to extrapolate higher-order coefficients that were not used for the construction of the rational approximants. This feature, which was apparently first observed by Gilewicz , has so far been used quite extensively in the case of Padé approximants for divergent perturbation expansions from quantum field theory and related disciplines. We refer to the investigations by Elias, Steele and Chishtie et al. , by Samuel, Gardi, Karliner, Ellis et al. and by Cvetič et al. . An analogous approach also works in the case of the delta transformation . Recursive techniques for the Padé prediction of unknown series coefficients were developed recently .
It is the intention of this article to discuss – augmenting previous investigations – how additional information on the large-order asymptotics of the perturbative coefficients can be incorporated effectively into resummation and prediction schemes. Our emphasis will be on the Borel-Padé method which was introduced by Graffi, Grecchi, and Simon . This is a variant of the Borel method that uses Padé approximants for performing the analytic continuation of the Borel transformed series to a neighborhood of the positive real semiaxis. Further details on the Borel-Padé method as well as on other resummation methods can be found in Section 2.
We would like to mention here the investigations (based on the Borel method) by Fischer and by Caprini and Fischer regarding the resummation of divergent perturbative expansions in quantum field theory. In , Caprini and Fischer use asymptotic information (location of the first IR and UV renormalon poles) for the construction of conformal mappings in the Borel plane. Here, we describe alternative ways in which the asymptotics of the coefficients of the perturbative expansion can be utilized to enhance the effectiveness of resummation procedures and perturbative predictions. These improvements, which provide an alternative to the methods presented in , are not necessarily restricted to the first IR and UV renormalons, but can take advantage of the location of all known poles in the Borel plane. Regarding asymptotic properties of perturbative coefficients in quantum field theory we also refer to .
Cvetič and Yu in consider the resummation of the *real* part of the (one-loop) QED effective action (or vacuum-to vacuum amplitude) in the presence of background electric and magnetic fields, of which the exact nonperturbative answer is known . In the presence of an electric field, the QED effective action acquires an imaginary part proportional to the pair production amplitude for real electron-positron (or lepton-antilepton) pairs. The real part of the effective action does not constitute the full physical solution. The full, complex-valued answer requires an integration in the complex plane, for example along the special contour introduced in . Only in the case of the pure magnetic field, which we consider in this article, the QED effective action is entirely real. For this particular example, the delta transformation employed in or the Borel-Padé Cauchy principal value method used in provide the full, i.e. entirely real, not complex physical solution. In view of the above, we again consider here only the resummation of the pure magnetic field. The electric field or the electric field combined with a magnetic field should be treated along the ideas introduced in .
We briefly discuss the relation of resummation and perturbative predictions: There is evidence (see for example and references therein) that the *resummation* of divergent series is ambiguous, especially if these series do not fulfill a Carleman condition (see or Theorems XII.17 and XII.19 in ). However, the *prediction* of perturbative coefficients apparently does not suffer from such ambiguities. Perturbative predictions should be possible even in those cases where the perturbation series is evaluated “on the cut” in the complex plane. For the prediction of unknown perturbation series coefficients, the nonanalytic contributions, which are responsible for the ambiguities, are irrelevant. Only the analytic part of the function, which is represented by a divergent power series, matters. The terms of an otherwise more problematic divergent nonalternating series can be predicted just as well as the terms of an alternating series. That is to say, the resummation of a divergent series is not necessarily unique (see the different integration contours in ), but the perturbative coefficients, which can be extrapolated and predicted via rational approximants, are uniquely determined even though the complex integration along the different contours in leads to different results for the nonperturbative, nonanalytic contributions.
This paper is organized as follows. In Section 2 we present a short account of various resummation methods for divergent series. The exploitation of additional available asymptotic information in resummation algorithms is discussed in Section 3. Asymptotically optimized perturbative predictions (i.e., predictions of higher-order unknown perturbative coefficients) are discussed in Sections 4 and 5 for various example cases. We conclude with a discussion of the results in Section 6.
## 2 A REVIEW OF RESUMMATION METHODS
In this Section, we provide a condensed description of the resummation methods under consideration in this article. Let
$$f(z)\underset{\nu =0}{\overset{\mathrm{}}{}}\gamma _\nu z^\nu $$
(1)
be a (formal) power series for some function $`f`$. Then, we define the $`(\kappa ,\lambda )`$-generalized Borel integral transform according to
$`f(z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}t^{\lambda 1}^{(\kappa ,\lambda )}(f;zt^\kappa )\mathrm{exp}\left(t\right)dt`$ (2)
$`=`$ $`z^{\lambda /\kappa }{\displaystyle _0^{\mathrm{}}}s^{\lambda 1}^{(\kappa ,\lambda )}(f;s^\kappa )\mathrm{exp}\left(s/z^{1/k}\right)ds.`$ (3)
Here,
$$^{(\kappa ,\lambda )}(f;z)=\underset{\nu =0}{\overset{\mathrm{}}{}}\frac{\gamma _\nu }{\mathrm{\Gamma }(\kappa \nu +\lambda )}z^\nu $$
(4)
is the $`(\kappa ,\lambda )`$-generalized Borel transformed series of the power series (1) for $`f(z)`$. For $`\kappa =\lambda =1`$, we recover the usual formulas for the Borel transformation:
$`f(z)`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}(f;zt)\mathrm{exp}\left(t\right)dt`$ (5)
$`=`$ $`{\displaystyle \frac{1}{z}}{\displaystyle _0^{\mathrm{}}}(f;s)\mathrm{exp}\left(s/z\right)ds,`$ (6)
$`(f;z)`$ $`=`$ $`^{(1,1)}(f;z)={\displaystyle \underset{\nu =0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\gamma _\nu }{\nu !}}z^\nu .`$ (7)
There exists an extensive literature on the Borel method in general and on physical applications in special. Any attempt to provide something resembling a reasonably complete bibliography would clearly be beyond the scope of this article. Let us just mention that recent monographs on the Borel method and related topic were published by Shawyer and Watson and Sternin and Shatalov .
Let us assume that the coefficients $`\gamma _n`$ of the power series (1) possess the following large order asymptotics,
$$\gamma _nA\mathrm{\Gamma }(\kappa n+\lambda )B^n,n\mathrm{},$$
(8)
where $`A`$, $`B`$, $`\kappa `$, and $`\lambda `$ are suitable constants. We say that a $`(\kappa ^{},\lambda ^{})`$-generalized Borel method for some power series (1) is *asymptotically optimized* if the parameters $`\kappa ^{}`$ and $`\lambda ^{}`$ agree with the parameters $`\kappa `$ and $`\lambda `$ in the large-order asymptotics (8) for the series coefficient $`\gamma _n`$. Thus, the leading (hyper)factorial growth of $`\gamma _n`$ is exactly canceled out in this case. If we know in addition the parameter $`B`$ in (8), then we can immediately deduce that the asymptotically optimized Borel transformed series (4) possesses a pole at $`z=1/B`$.
The most difficult computational problem, which normally occurs in the context of a Borel summation process, is the construction of an analytic continuation for the Borel transformed series (4). If the coefficients $`\gamma _n`$ of the power series (1) satisfy (8), then the $`(\kappa ,\lambda )`$-generalized Borel transformed series (4) has a *nonzero* but *finite* radius of convergence. In order to be able to do the integration, we now need an analytic continuation which extends $`^{(\kappa ,\lambda )}(f;z)`$ from the interior of its circle of convergence to a neighborhood which contains the whole positive semiaxis. In this article, we emphasize the Borel-Padé method which was introduced by Graffi, Grecchi, and Simon and which accomplishes the analytic continuation by converting the partial sums of the Borel transformed series (4) to Padé approximants. It may happen that the Padé approximants thus constructed exhibit poles along the positive real axis. In such a case, the function $`f`$ is – strictly speaking – not Borel-Padé summable. However, by further generalizing the integrals (2) and (3) – either via a principal-value prescription , or by employing conformal mappings , or via the special integration contours used in – it may nevertheless be possible to associate a finite value to the divergent power series (1) for $`f`$.
In the case of Padé approximants we use the notation and conventions of the monograph by Baker and Graves-Morris . Thus, a Padé approximant $`[l/m]_f(z)`$ to $`f(z)`$ corresponds to the ratio of two polynomials $`P_l(z)`$ and $`Q_m(z)`$, which are of degrees $`l`$ and $`m`$, respectively, in $`z`$:
$$[l/m]_f(z)=\frac{P_l(z)}{Q_m(z)}=\frac{p_0+p_1z+\mathrm{}+p_lz^l}{1+q_1z+\mathrm{}+q_mz^m}.$$
(9)
The polynomials $`P_l(z)`$ and $`Q_m(z)`$ are constructed so that the Taylor expansion of the Padé approximation agrees with the original input series (1) up to terms of order $`l+m`$ in $`z`$,
$$f(z)[l/m]_f(z)=\mathrm{O}\left(z^{l+m+1}\right),z0.$$
(10)
This asymptotic error estimate leads to a system of linear equations by means of which the coefficients $`p_0`$, $`p_1`$, …, $`p_l`$ and $`q_1`$, $`q_2`$, …, $`q_m`$ in (9) can be computed. However, there are also several algorithms which permit a recursive computation of Padé approximants. A discussion of the merits and weaknesses of the various computational approaches can for instance be found in Section II.3 of the book by Cuyt and Wuytack .
An example of such a recursive algorithm is provided by Wynn’s epsilon algorithm :
$`ϵ_1^{(n)}`$ $`=`$ $`0,ϵ_0^{(n)}=s_n,n𝑵_0,`$ (11)
$`ϵ_{k+1}^{(n)}`$ $`=`$ $`ϵ_{k1}^{(n+1)}+\mathrm{\hspace{0.17em}1}/[ϵ_k^{(n+1)}ϵ_k^{(n)}],k,n𝑵_0.`$ (12)
Wynn showed that if the input data $`s_n`$ for the epsilon algorithm are the partial sums of the (formal) power series (1) for some function $`f(z)`$ according to
$$s_n=f_n(z)=\underset{\nu =0}{\overset{n}{}}\gamma _\nu z^\nu ,$$
(13)
then the elements $`ϵ_{2k}^{(n)}`$ with even subscripts are Padé approximants to $`f`$ according to
$$ϵ_{2k}^{(n)}=[n+k/k]_f(z).$$
(14)
In contrast, the elements $`ϵ_{2k+1}^{(n)}`$ with odd subscripts are only auxiliary quantities which diverge if the whole process converges.
If one tries to sum a divergent power series or to accelerate the convergence of a slowly convergent series by converting its partial sums to Padé approximants, it is usually a good idea to use either diagonal Padé approximants, whose numerator and denominator polynomials have equal degrees, or – if this is not possible – to use Padé approximants with numerator and denominator polynomials whose degrees differ as little as possible. If we use the epsilon algorithm for the computation of the Padé approximants, then Eq. (14) implies that we then obtain the following staircase sequence in the Padé table (see Eq. (4.3-7) of ):
$$[0/0],[1/0],[1/1],\mathrm{},[\nu /\nu ],[\nu +1/\nu ],[\nu +1/\nu +1],\mathrm{}.$$
(15)
This staircase sequence exploits the available information optimally if the partial sums $`f_m(z)`$ with $`m0`$ are computed successively and if after the computation of each new partial sum the element of the epsilon table with the highest possible even transformation order is computed. With the help of the notation $`[[x]]`$ for the integral part of $`x`$, this staircase sequence can be written compactly as follows:
$$ϵ_{2[[n/2]]}^{(n2[[n/2]])}=\left[n[[n/2]]/[[n/2]]\right]_f(z),n=0,1,2,\mathrm{}.$$
(16)
The asymptotic error estimate (10) implies that all series coefficients, which are employed for the computation of the Padé approximant $`[l/m]_f(z)`$, are recovered by a Taylor expansion. Consequently, the higher order derivatives of the Padé approximant provide predictions for “unknown” series coefficients, i.e. to those series coefficients that were not used for computation of $`[l/m]_f(z)`$.
There is an enormous amount of literature on Padé approximants in general as well as on their application in theoretical physics. Let us just mention that the popularity of Padé approximants in theoretical physics can be traced back to a review by Baker , that the monograph by Baker and Graves-Morris is the currently most complete source of information on Padé approximants, and that an account of the historical development of Padé approximants and related topics is given in a monograph by Brezinski .
The intense research on Padé approximants during the last decades of course also showed that Padé approximants suffer – like all other numerical techniques – from certain limitations and weaknesses. For example, Padé approximants are in principle limited to convergent or divergent power series, but cannot help in the case of many other types of slowly convergent or divergent sequences. Moreover, Padé approximants are either not useful or cannot be applied at all in the case of power series whose coefficients $`\gamma _n`$ grow like $`(2n)!`$ or even $`(3n)!`$ . Consequently, the intense research on Padé approximants also stimulated research on related techniques, the so-called sequence transformations.
Let us assume that $`\{s_n\}_{n=0}^{\mathrm{}}`$ is a sequence, whose elements may for instance be the partial sums of an infinite series according to $`s_n=_{k=0}^na_k`$. A *sequence transformation* is a rule which maps a sequence $`\{s_n\}_{n=0}^{\mathrm{}}`$ to a new sequence $`\{s_n^{}\}_{n=0}^{\mathrm{}}`$ with hopefully better numerical properties. In this terminology, Padé approximants are just a special class of sequence transformations since they transform the partial sums of a (formal) power series to a doubly indexed sequence of rational approximants.
If $`\{s_n\}_{n=0}^{\mathrm{}}`$ either converge to some limit $`s`$ as $`n\mathrm{}`$ or can be summed to the generalized limit $`s`$ in the case of divergence, then a sequence element $`s_n`$ can for all $`n0`$ be partitioned into the (generalized) limit $`s`$ and a remainder $`r_n`$ according to
$$s_n=s+r_n.$$
(17)
Normally, a sequence transformation will not be able to determine the (generalized) limit $`s`$ of $`\{s_n\}_{n=0}^{\mathrm{}}`$ *exactly*. Thus, the elements of the transformed sequence $`\{s_n^{}\}_{n=0}^{\mathrm{}}`$ can also be partitioned into the (generalized) limit $`s`$ and a transformed remainder $`r_n^{}`$ according to
$$s_n^{}=s+r_n^{},$$
(18)
and the transformed remainders will in general be different from zero for all finite values of $`n`$.
In the literature on convergence acceleration it is said that a sequence transformation *accelerates convergence* if the transformed remainders $`\{r_n^{}\}_{n=0}^{\mathrm{}}`$ vanish more rapidly than the original remainders $`\{r_n\}_{n=0}^{\mathrm{}}`$ according to
$$\underset{n\mathrm{}}{lim}\frac{r_n^{}}{r_n}=\underset{n\mathrm{}}{lim}\frac{s_n^{}s}{s_ns}=\mathrm{\hspace{0.33em}0},$$
(19)
and a divergent sequences, whose remainders $`r_n`$ do not vanish as $`n\mathrm{}`$, is summed to its generalized limit $`s`$ if the transformed remainders $`r_n^{}`$ vanish as $`n\mathrm{}`$.
Thus, a sequence transformation essentially tries to eliminate the remainders $`r_n`$ from the sequence elements $`s_n`$ as effectively as possible. Since, however, an in principle unlimited variety of different remainders can occur, it is necessary to make some assumptions – either explicitly or implicitly – which provide the basis for the construction of a sequence transformation. A detailed discussion of the construction of sequence transformations as well as many examples can be found in the book by Brezinski and Redivo Zaglia or in .
Normally, the assumptions being made are incorporated into the transformation scheme via model sequences, whose remainders possess a particular simple structure and can be expressed by a finite number of terms:
$$\stackrel{~}{s}_n=\stackrel{~}{s}+\underset{k=0}{\overset{k1}{}}\stackrel{~}{c}_j\phi _j(n).$$
(20)
Here, the $`\stackrel{~}{c}_j`$ are unspecified coefficients, and the $`\phi _j(n)`$ are by assumption known functions of $`n`$.
The elements of this model sequence contain $`k+1`$ unknown, the (generalized) limit $`\stackrel{~}{s}`$ and the $`k`$ coefficients $`\stackrel{~}{c}_j`$ with $`0jk1`$. Since all unknowns occur *linearly*, it is possible to construct a sequence transformation $`𝒯`$ – if necessary via Cramer’s rule – which is *exact* for the elements of this model sequence according to
$$𝒯=𝒯(\stackrel{~}{s}_n,\stackrel{~}{s}_{n+1},\mathrm{},\stackrel{~}{s}_{n+k},)=\stackrel{~}{s},$$
(21)
if applied to the numerical values of $`k+1`$ sequence elements $`\stackrel{~}{s}_n,\stackrel{~}{s}_{n+1},\mathrm{},\stackrel{~}{s}_{n+k}`$.
Of course, simple model sequences of that kind normally do not occur in practical problems. However, their elements provide at least for sufficiently large values of $`k`$ reasonably accurate approximations to the elements of the more realistic sequence
$$s_n=s+\underset{j=0}{\overset{\mathrm{}}{}}c_j\phi _j(n).$$
(22)
If we now apply this sequence transformation $`𝒯`$ to the numerical values of $`k+1`$ elements of the sequence (22), then we have no reason to assume that $`𝒯`$ might produce its exact (generalized) limit $`s`$. However, a more detailed mathematical analysis of the transformation process normally reveals that $`𝒯`$ eliminates the first $`k`$ terms $`c_j\phi _j(n)`$ with $`0jk1`$. Thus, the transformed remainder $`r_n^{}`$ starts with $`\phi _k(n)`$ instead of $`\phi _0(n)`$, which for sufficiently large values of $`k`$ normally constitutes a significant improvement.
Most sequence transformations can be constructed on the basis of model sequences of the type of Eq. (20). For example, Wynn could show that his epsilon algorithm is exact for model sequences whose remainders can be expressed as a linear combination of exponential terms according to
$$\stackrel{~}{s}_n=\stackrel{~}{s}+\underset{j=0}{\overset{k1}{}}\stackrel{~}{c}_j\lambda _j^n.$$
(23)
Concerning the $`\lambda _j`$ it is only assumed that they are different from zero and one and ordered according to magnitude, i.e., $`\lambda _j0,1`$ and $`|\lambda _0|>|\lambda _1|>|\lambda _{k1}|>0`$. Thus, if the numerical values of $`2k+1`$ elements $`\stackrel{~}{s}_n`$, $`\stackrel{~}{s}_{n+1}`$, …, $`\stackrel{~}{s}_{n+2k}`$ of this model sequence are available, then the epsilon algorithm is exact according to
$$ϵ_{2k}^{(n)}=\stackrel{~}{s}.$$
(24)
Moreover, Wynn constructed in Theorems 16 and 17 of asymptotic expansions ($`n\mathrm{}`$) for the transformed remainders $`r_n^{}`$ created by the application of $`ϵ_{2k}^{(n)}`$ to the elements of the sequence
$$s_n=s+\underset{j=0}{\overset{\mathrm{}}{}}c_j\lambda _j^n,$$
(25)
which is an obvious generalization of the model sequence (23). He showed that the transformed remainders $`r_n^{}`$ are proportional to $`\lambda _k^n`$ which corresponds to an elimination of the first $`k`$ exponential terms $`c_j\lambda _j^n`$ on the right-hand side of (25). Since the $`\lambda _j`$ are by assumption ordered in magnitude, this constitutes a significant achievement. Consequently, Wynn’s epsilon algorithm is *asymptotically optimal* for sequences of the type of (25). This means that no other sequence transformation, which also uses only the numerical values of the elements of the sequence (25) as input data, can produce a better asymptotic ($`n\mathrm{}`$) truncation error.
Levin introduced a class of sequence transformations which are exact for model sequences of the following type:
$$\stackrel{~}{s}_n=\stackrel{~}{s}+\omega _nz_n.$$
(26)
Here, $`\omega _n`$ is an estimate for the truncation error $`\stackrel{~}{r}_n`$, and $`z_n`$ is a correction term. Levin assumed that $`z_n`$ can be expressed as a truncated power series in $`1/(n+\zeta )`$ where $`\zeta `$ is a positive shift parameter. In Sections 7 – 9 of , several other sequence transformations were constructed which are also exact for the model sequence (26) but make different assumptions about the correction terms $`z_n`$. The remainder estimates $`\omega _n`$ introduce additional degrees of freedom in the construction of the sequence transformation as compared to Padé approximants. One may draw an analogy between sequence transformations and Padé approximants on the one hand and the Gaussian integration and the Simpson rule on the other hand; the variable integration nodes and weight factors of the Gaussian integration yield additional degrees of freedom which may be used in order to construct a potentially much more powerful algorithm for numerical integration.
In the following text, we will concentrate on sequence transformations which assume that $`z_n`$ can be expressed as a truncated factorial series (Section 8 of ):
$$z_n=\underset{j=0}{\overset{k1}{}}\stackrel{~}{c}_j/(\zeta +n)_j.$$
(27)
Here, $`(n+\zeta )_j=\mathrm{\Gamma }(n+\zeta +j)/\mathrm{\Gamma }(n+\zeta )`$ is a Pochhammer symbol, and $`\zeta `$ is a positive shift parameter. The assumption (27) implies that the sequence transformations derived in this way are particularly well suited for sequences satisfying
$$s_n=s+\omega _n\underset{j=0}{\overset{\mathrm{}}{}}c_j/(\zeta +n)_j.$$
(28)
It is a priori not obvious that the ratio $`[s_ns]/\omega _n=r_n/\omega _n`$ can be expressed as a factorial series. Nevertheless, this assumption leads to powerful sequence transformations which are apparently particularly well suited for the summation of factorially divergent series . Let $`s_n=_{k=0}^na_k`$ be the partial sums of an infinite series. If the $`a_k`$ strictly alternate in sign and decrease monotonously in magnitude, then the best simple estimate for the truncation error $`r_n=_{k=n+1}^{\mathrm{}}a_k`$ is the first term $`a_{n+1}`$ neglected in the partial sum $`s_n`$. Moreover, the first term neglected is also the best simple remainder estimate for many factorially divergent alternating series (see for example Theorem 13-2 of ). The mathematical structure of a factorially divergent series is expected for the perturbative expansions in quantum field theory. Further arguments supporting the general applicability of the delta transformation to the series of the mathematical structure as expected for quantum field theory will be discussed in .
If we now combine the assumption that $`z_n`$ should be a truncated factorial series according to (27) with the remainder estimate
$$\omega _n=\mathrm{\Delta }s_n=a_{n+1},$$
(29)
which was introduced by Smith and Ford , we obtain the delta transformation which is defined by the following ratio of finite sums (Eq. (8.4-4) of ):
$`\delta _k^{(n)}(\zeta ,s_n)`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }^k\left[(\zeta +n)_{k1}s_n/\mathrm{\Delta }s_n\right]}{\mathrm{\Delta }^k\left[(\zeta +n)_{k1}/\mathrm{\Delta }s_n\right]}}`$ (30)
$`=`$ $`{\displaystyle \frac{{\displaystyle \underset{j=0}{\overset{k}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right){\displaystyle \frac{(\zeta +n+j)_{k1}}{(\zeta +n+k)_{k1}}}{\displaystyle \frac{s_{n+j}}{\mathrm{\Delta }s_{n+j}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right){\displaystyle \frac{(\zeta +n+j)_{k1}}{(\zeta +n+k)_{k1}}}{\displaystyle \frac{1}{\mathrm{\Delta }s_{n+j}}}}}.`$
Here, the same notation as in is used. Thus, $`\mathrm{\Delta }`$ stands for the difference operator defined by $`\mathrm{\Delta }g(n)=g(n+1)g(n)`$, $`(a)_n=\mathrm{\Gamma }(a+n)/\mathrm{\Gamma }(a)`$ is a Pochhammer symbol, $`k`$ and $`n`$ are nonnegative integers, and $`\zeta `$ is a shift parameter which has to be positive to allow $`n=0`$ in Eq. (30). The most obvious choice, which is always used in this article, is $`\zeta =1`$.
In Section 8.3 of , a simple recursive scheme is described which permits – depending upon the initial values – the recursive calculation of either the numerator or the denominator sum of $`\delta _k^{(n)}(\zeta ,s_n)`$.
In the context of quantum field theory and related disciplines, the delta transformation (30) may be used for the summation of divergent perturbation expansions which are power series in some coupling constant. Thus, if we replace the input data $`s_n`$ in (30) by the partial sums $`f_n(z)=_{\nu =0}^n\gamma _\nu z^\nu `$ of the (formal) power series (1) for $`f(z)`$, we obtain a rational expression, whose numerator and denominator polynomials are of degrees $`k+n`$ and $`k`$, respectively, in $`z`$:
$$\delta _k^{(n)}(\zeta ,f_n(z))=\frac{{\displaystyle \underset{j=0}{\overset{k}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right){\displaystyle \frac{(\zeta +n+j)_{k1}}{(\zeta +n+k)_{k1}}}{\displaystyle \frac{z^{kj}f_{n+j}(z)}{\gamma _{n+j+1}}}}{{\displaystyle \underset{j=0}{\overset{k}{}}}(1)^j\left({\displaystyle \genfrac{}{}{0pt}{}{k}{j}}\right){\displaystyle \frac{(\zeta +n+j)_{k1}}{(\zeta +n+k)_{k1}}}{\displaystyle \frac{z^{kj}}{\gamma _{n+j+1}}}}.$$
(31)
If the coefficients $`\gamma _n`$ of the power series for $`f(z)`$ are all different from zero, the rational function (31) satisfies the asymptotic error estimate (Eq. (4.29) of )
$$f(z)\delta _k^{(n)}(\zeta ,f_n(z))=O(z^{k+n+2}),z0.$$
(32)
This estimate, which is formally very similar to the analogous estimate (10) for Padé approximants, implies that all terms of the formal power series, which are used for construction of the rational approximant $`\delta _k^{(n)}(\zeta ,f_n(z))`$, are reproduced exactly by a Taylor expansion around $`z=0`$. Moreover, the higher order derivatives provide just like in the Padé case predictions for those coefficients $`\gamma _{n+k+2}`$, $`\gamma _{n+k+3}`$, …, that were not used for the construction of the rational function.
As already discussed, the power of the delta transformation or of other Levin-type transformations results from the fact that an explicit estimate for the truncation error is incorporated into the transformation scheme. The truncation error estimate used by the delta transformation is the first term $`\gamma _{n+1}z^{n+1}`$ neglected in the partial sum $`f_n(z)=_{\nu =0}^n\gamma _\nu z^\nu `$. Consequently, for a proper application of the delta transformation *all* coefficients $`\gamma _n`$ of the power series for $`f`$ with $`n1`$ have to be different from zero because otherwise the estimate for the truncation error makes no sense. This restriction also follows directly from the ratio (31), where undefined expressions occur if coefficients $`\gamma _n`$ with $`n1`$ are zero (cf. Eq. (11) in which entails divisions by zero).
## 3 ASYMPTOTICALLY OPTIMIZED RESUMMATION
The problem of the resummation of divergent perturbative expansions in quantum field theory and related disciplines has been discussed in a number of recent publications, for example in . We investigate here asymptotically optimized resummation methods, i.e. methods which utilize information about large-order asymptotics of perturbative coefficients with the intention of enhancing the rate of convergence of the resummation algorithm.
We discuss here possible improvements of the Borel-Padé method on the basis of potentially available information about the large-order asymptotics of perturbative coefficients. We concentrate on the particular model example discussed recently by Dunne and Hall , by Cvetič and Yu in and by ourselves in . We discuss the QED effective action $`S_\mathrm{B}`$ in the presence of a constant background magnetic field. The exact nonperturbative result for $`S_\mathrm{B}`$ can be expressed as a proper-time integral:
$$S_\mathrm{B}=\frac{e^2B^2}{8\pi ^2}\underset{0}{\overset{\mathrm{}}{}}\frac{\mathrm{d}s}{s^2}\left\{\mathrm{coth}s\frac{1}{s}\frac{s}{3}\right\}\mathrm{exp}\left(\frac{m_\mathrm{e}^2}{eB}s\right).$$
(33)
Here, $`B`$ is the magnetic field strength, and $`m_\mathrm{e}`$ and $`e`$ are the mass and the charge of the electron, respectively (this result is given for example in Eq. (4-123) in ).
The integral representation (33) for $`S_\mathrm{B}`$ can be expressed as a strictly alternating perturbation series in the effective coupling coupling $`g_\mathrm{B}=e^2B^2/m_\mathrm{e}^4`$:
$$S_\mathrm{B}=\frac{2e^2B^2}{\pi ^2}g_\mathrm{B}\underset{n=0}{\overset{\mathrm{}}{}}c_ng_\mathrm{B}^n,$$
(34)
where
$$c_n=\frac{(1)^{n+1}\mathrm{\hspace{0.17em}4}^n\left|_{2n+4}\right|}{(2n+4)(2n+3)(2n+2)}.$$
(35)
Here, $`_{2n+4}`$ is a Bernoulli number. Thus, the perturbation expansion (34) has the remarkable feature that an unlimited number of series coefficients $`c_n`$ are known analytically. Consequently, this series is particularly well suited as a model system for studying resummation methods.
Next, we utilize the fact that a Bernoulli number with even index can be expressed by a Riemann zeta function according to (see Section 23.2. on p. 807 of )
$$\left|_{2n}\right|=\frac{2(2n)!}{(2\pi )^{2n}}\zeta (2n).$$
(36)
Inserting this into (34) and (35) yields
$`S_\mathrm{B}`$ $`=`$ $`{\displaystyle \frac{e^2B^2}{8\pi ^2}}g_\mathrm{B}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n(2n+1)!{\displaystyle \frac{2\zeta (2n+4)}{\pi ^{2n+4}}}g_\mathrm{B}^n`$ (37)
$`=`$ $`{\displaystyle \frac{e^2B^2}{8\pi ^2}}g_\mathrm{B}S_\mathrm{B}^{},`$ (38)
where in the last line we define implicitly the scaled function $`S_\mathrm{B}^{}`$ which is also considered, e.g., in Table 2 of . It is a direct consequence of the Dirichlet series (Section 23.2. of )
$$\zeta (s)=\underset{m=0}{\overset{\mathrm{}}{}}(m+1)^s,$$
(39)
that we have the inequality $`1\zeta (2n+4)\zeta (4)`$ for all nonnegative integers $`n`$. Consequently, the zeta function does not contribute to the factorial divergence of the perturbation series (34). Thus, the factorial $`(2n+1)!`$ on the right-hand side of (37) implies that the perturbation series for $`S_\mathrm{B}`$ diverges for every coupling $`g_\mathrm{B}0`$. Furthermore, it is clear from the representation (37) that an asymptotically optimized $`(\kappa ,\lambda )`$-generalized Borel resummation scheme for $`S_\mathrm{B}`$ according to (2)–(4) requires the parameter setting $`\kappa =\lambda =2`$.
We now discuss the construction of the asymptotically optimized Borel transform explicitly. We start from the scaled series,
$$S_\mathrm{B}^{}(g_\mathrm{B})=\underset{n=0}{\overset{\mathrm{}}{}}(16c_n)g_\mathrm{B}^n=\underset{n=0}{\overset{\mathrm{}}{}}(1)^n(2n+1)!\frac{2\zeta (2n+4)}{\pi ^{2n+4}}g_\mathrm{B}^n.$$
(40)
The $`(2,2)`$-generalized Borel transformed series of $`S_\mathrm{B}^{}`$ is given by
$$^{(2,2)}(S_\mathrm{B}^{};z)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{16c_n}{(2n+1)!}g_\mathrm{B}^n=\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\frac{2\zeta (2n+4)}{\pi ^{2n+4}}z^n.$$
(41)
This series can be brought into a form which clearly shows the singularity structure of the Borel transformed series. For that purpose, we replace the Riemann zeta function by its Dirichlet series according to (39) and interchange the order of the two infinite nested summations:
$$^{(2,2)}(S_\mathrm{B}^{};z)=\frac{2}{\pi ^4}\underset{m=0}{\overset{\mathrm{}}{}}(m+1)^4\underset{n=0}{\overset{\mathrm{}}{}}\left\{z/\left[\pi (m+1)\right]^2\right\}^n.$$
(42)
Thus, $`^{(2,2)}(S_\mathrm{B}^{};z)`$ is essentially a superposition of geometric series with arguments $`z/\left[\pi (m+1)\right]^2`$. If we now use $`_{k=0}^{\mathrm{}}(x)^k=1/(1+x)`$, we obtain:
$$^{(2,2)}(S_\mathrm{B}^{};z)=\frac{2}{\pi ^4}\underset{m=0}{\overset{\mathrm{}}{}}\frac{(m+1)^4}{1+z/\left[\pi (m+1)\right]^2}.$$
(43)
This representation shows that the poles of the Borel transformed series $`^{(2,2)}(S_\mathrm{B}^{};z)`$ are located along the negative real axis according to
$$z=n^2\pi ^2,$$
(44)
where $`n`$ is a nonzero positive integer. Moreover, we obtain the following representations for the QED effective action $`S_\mathrm{B}`$ as a $`(2,2)`$-generalized Borel integral according to Eqs. (2) and (3):
$`S_\mathrm{B}`$ $`=`$ $`{\displaystyle \frac{e^2B^2}{8\pi ^2}}g_\mathrm{B}{\displaystyle _0^{\mathrm{}}}t^{(2,2)}(S_\mathrm{B}^{};g_\mathrm{B}t^2)\mathrm{exp}(t)dt`$ (45)
$`=`$ $`{\displaystyle \frac{e^2B^2}{8\pi ^2}}{\displaystyle _0^{\mathrm{}}}s^{(2,2)}(S_\mathrm{B}^{};s^2)\mathrm{exp}(s/g_\mathrm{B}^{1/2})ds.`$ (46)
Cvetič and Yu in use a $`(2,2)`$-generalized Borel transformed series for the QED effective action, constructed according to Eq. (45). As explained in Section 1, this transformation is asymptotically optimized in the sense that the leading factorial growth of the perturbative coefficients in Eq. (37) is divided out. It could appear from the Eqs. (4) and (8) in that a $`(1,1)`$-generalized Borel transform is used where the $`n`$th perturbative coefficient is divided by a factor of $`n!=\mathrm{\Gamma }(n+1)`$. This is, however, not the case. Note that in , the perturbative expansions are written in a very peculiar parameter, which is
$$\stackrel{~}{b}=eB/m_\mathrm{e}^2.$$
(47)
In normal QED terminology, this would correspond to an expansion in powers of $`\sqrt{\alpha }`$. By consequence, all even-order perturbative coefficients vanish in the analysis presented in Ref. . In this context, one may to note that the expression for the delta transformation according to Eq. (11) in is actually undefined since it involves divisions by zero. The expansion parameter used here is $`g_\mathrm{B}=e^2B^2/m_\mathrm{e}^4=\stackrel{~}{b}^2`$; this parameter is also used in and in .
Due to the special, mathematically compact form of the perturbative coefficients in Eq. (35), the asymptotically optimized Borel transform (42) is simply a superposition of geometric series. Such a simple mathematical structure cannot be expected to be of general importance concerning series occurring in quantum field theory. Note that this series can be brought in a form which clearly shows that Wynn’s epsilon algorithm, which computes Padé approximants according to (14), is optimal, as discussed in Section 2. For that purpose, we rewrite the $`n`$th partial sum of the series (43) as follows:
$$\underset{\nu =0}{\overset{n}{}}(1)^\nu \frac{2\zeta (2\nu +4)}{\pi ^{2\nu +4}}z^\nu =^{(2,2)}(S_\mathrm{B}^{};x)\underset{\nu =n+1}{\overset{\mathrm{}}{}}(1)^\nu \frac{2\zeta (2\nu +4)}{\pi ^{2\nu +4}}z^\nu .$$
(48)
On substituting the Dirichlet series (39) into the infinite series on the right-hand side and interchanging the order of summations, we obtain
$$\underset{\nu =0}{\overset{n}{}}(1)^\nu \frac{2\zeta (2\nu +4)}{\pi ^{2\nu +4}}z^\nu =^{(2,2)}(S_\mathrm{B}^{};z)+\frac{2(1)^n}{\pi ^4}\underset{m=0}{\overset{\mathrm{}}{}}\frac{z(m+1)^6}{\pi ^2+z/(m+1)^2}\left(\frac{z}{[\pi (m+1)]^2}\right)^n.$$
(49)
Thus, the partial sum of the Borel transformed series (43) possesses the following general structure:
$$s_n=s+(1)^n\underset{j=0}{\overset{\mathrm{}}{}}c_j\lambda _j^n.$$
(50)
This sequence is obviously a special case of the sequence (25), for which Wynn’s epsilon algorithm is – as discussed in Section 2 – asymptotically optimal.
The conclusions drawn by Cvetič and Yu in appear to be restricted, at least in part, to the particular model example studied in their paper. In this context it should be emphasized that the superiority of the delta transformation over Padé approximants if applied directly to factorially divergent series, cannot be assumed to persist after the Borel transformation by which the leading factorial divergence is divided out. I.e., the delta transformation is more powerful than the Padé technique for factorially divergent series, but this finding by no means allows us to conclude that, or in fact has any connection to the assumption that the combined Borel-delta technique should be numerically superior to the Borel-Padé method. This consideration is relevant for the interpretation of the conclusions drawn by Cvetič and Yu with regard to the variant of the Borel method proposed in , which uses the delta transformation for the analytic continuation of the Borel transformed series and which is called the Borel-Weniger method by the authors of .
We return now to the discussion of further improvements of the asymptotically optimized Borel-Padé method. The leading asymptotics do not only permit to modify (optimize) the Borel transform accordingly, but indeed it is the leading large-order asymptotics which determine the location of the poles in the Borel plane. In view of Eq. (43), the singularities of the function $`^{(2,2)}(S_\mathrm{B}^{};z)`$ defined in (41) are at $`z=n^2\pi ^2`$, i.e. along the negative real axis. This is where one would expect them to lie in (distant) analogy to the renormalon theory .
Using the information on the location of the poles, it is possible to construct further improved Padé approximants. To this end, we utilize the known location of the poles in order to construct improved Padé approximants to the function $`^{(2,2)}(S_\mathrm{B}^{};z)`$. Normally, the Borel integral (43) would be evaluated with the upper- or lower-diagonal Padé approximants to $`^{(2,2)}(S_\mathrm{B}^{};z)`$ in the integrand. We use upper-diagonal Padé approximants here, as they can be computed by Wynn’s epsilon algorithm according to (14). We denote by $`𝒫_n(z)`$ the upper-diagonal Padé approximant,
$$𝒫_n(z)=\left[[[(n+1)/2]]/[[n/2]]\right]_{^{(2,2)}(S_\mathrm{B}^{})}(z).$$
(51)
In the upper-diagonal case, we evaluate the transforms $`𝒯S_{\mathrm{B},n}^{}`$ where
$`𝒯S_{\mathrm{B},n}^{}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}t𝒫_n\left(g_\mathrm{B}t^2\right)\mathrm{exp}(t)dt,`$ (52)
and observe numerical convergence of the transform at large transformation order $`n`$. When the location of the poles is known, we may improve the convergence of the transforms by the following replacement,
$$𝒫_n(z)𝒫_n^{}(z)=\frac{𝒬_n(z)}{\underset{i=1}{\overset{[[n/2]]}{}}(1+z/(n^2\pi ^2))},$$
(53)
where $`𝒬_n(z)`$ is the $`[[[(n+1)/2]]/0]`$-Padé approximant to the function $`_n(z)`$,
$$𝒬_n(z)=\left[[[(n+1)/2]]/0\right]__n(z),$$
(54)
and $`_n(z)`$ is given by
$$_n(z)=\underset{i=1}{\overset{[[n/2]]}{}}(1+z/(n^2\pi ^2))^{(2,2)}(S_\mathrm{B}^{};z).$$
(55)
The asymptotic enhancement is possible only if additional asymptotic information is available on the perturbative coefficients. Such information may be available (renormalon poles), but this is not necessarily provided. In this context it should be noted that there is currently no general proof of the assumption that the renormalon poles are the only relevant poles in the Borel plane , but the factorial divergence of the perturbative coefficients is a commonly accepted assumption .
The pole-structure improved transforms are obtained from (52) by the replacement (53),
$`𝒯^{}S_{\mathrm{B},n}^{}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}t𝒫_n^{}\left(g_\mathrm{B}t^2\right)\mathrm{exp}(t)dt.`$ (56)
Similar improvement of the convergence of transforms can also be expected in those cases where the final evaluation of the Borel integral proceeds in the complex plane along the integration contours introduced in .
Further improvement of the rate of convergence is possible by taking the transforms $`𝒯^{}S_{\mathrm{B},n}^{}`$ as input data to the epsilon algorithm (11) in order to accelerate the convergence of the sequence of the pole-structure improved transforms $`\{𝒯^{}S_{\mathrm{B},n}^{}\}_{n=0}^{\mathrm{}}`$. The application of the epsilon algorithm defined in Eq. (11) to the pole-structure improved transforms results in a sequence of upper-diagonal Padé approximants which we denote by
$$𝒯^{\prime \prime }S_{\mathrm{B},n}^{}=ϵ_{2[[n/2]]}^{(n2[[n/2]])}.$$
(57)
As input data for the epsilon algorithm, we use
$$s_n=𝒯^{}S_{\mathrm{B},n}^{}.$$
(58)
In a second epsilon transformation we may in turn employ the $`𝒯^{\prime \prime }S_{\mathrm{B},n}^{}`$ as input data for a further application of the epsilon algorithm,
$$𝒯^{\prime \prime \prime }S_{\mathrm{B},n}^{}=ϵ_{2[[n/2]]}^{(n2[[n/2]])},$$
(59)
where we use $`s_n=𝒯^{\prime \prime }S_{\mathrm{B},n}^{}`$ as input data. This results in a sequence of pole-structure and doubly epsilon-improved transforms $`𝒯^{\prime \prime \prime }S_{\mathrm{B},n}^{}`$. The application of the epsilon algorithm further enhances the rate of convergence of the pole-structure improved transforms.
In Table 1 we present numerical data for the Borel transforms calculated according to Eq. (52) (in passing we note that these correspond the method proposed in ) and the transforms (59). The first 15 transforms calculated according to Eq. (52) exhibit convergence to 8 significant digits, whereas the pole-structure and epsilon improved transforms coincide with the exact result to within 18 significant digits.
The delta transformation is a general-purpose transformation which has been proven to be applicable to a wide variety of alternating factorially divergent series . In the context of the delta transformation, additional asymptotic information could be used in order to modify the remainder estimates $`\omega _n`$ defined in Eq. (29) (see also Eqs. (7.3-8), (8.2-7), (8.4-1) and (8.4-4) of ). Also, we note a rescaling of the perturbative coefficients as a potential source for further improvements . Work along these lines is currently in progress and will be presented elsewhere .
## 4 ASYMPTOTICALLY OPTIMIZED PREDICTIONS
We refer here to the predictions of unknown higher-order perturbative coefficients as perturbative predictions or perturbative extrapolations. As outlined in Section 2 and , these extrapolations are obtained by reexpanding certain rational approximants in powers of the coupling parameter. The next higher-order term obtained after the reexpansion can then be interpreted as a prediction for that perturbative coefficient. The rational approximants discussed in Section 2 fulfill accuracy-through-order relations, i.e., upon reexpansion in the coupling parameter, all the perturbative terms used for the construction of the rational approximant are reproduced \[see Eq. (10) for Padé approximants and Eq. (32) for the delta transformation\]. The coefficients of the Borel transformed series are related to those of the input series by Eq. (4). Therefore, we can either predict the perturbative coefficients of the original series or the coefficients of the Borel transformed series.
We consider here the asymptotic improvement of three different prediction methods: (i) asymptotically optimized predictions based on the combination of Borel and Padé techniques, (ii) predictions based on the delta transformation, and (iii) the direct application of Padé approximants to the perturbation series. We consider here the following improvements beyond reexpansion of the rational approximant,
1. A-posteriori corrections. These are further corrections to the perturbative predictions obtained by estimating not only the coefficient, but also the probable error in making that estimate. Similar methods in the context of Padé approximants have been introduced, e.g., in .
2. Fixing poles. In the context of the asymptotically improved Borel-Padé method, it is possible to improve predictions if the leading and subleading large-order asymptotics of the perturbative predictions are known. These asymptotics determine the location of the poles in the Borel plane, which can be put in by hand (see also Eqs. (53)–(55) in Section 3 and, in part, ).
3. Renormalization group. The renormalization group can be used to enhance perturbative predictions for certain classes of diagrams; this has been used e.g. in for the anomalous magnetic moment of the muon.
It is natural to assume that combinations of these techniques should be investigated where appropriate.
The basic idea of a-posteriori corrections is as follows. The errors made in the “prediction” of lower-order coefficients are available by the time we come to higher order, so they may be utilized for an estimate of the error which is to be expected in a prediction of the next higher-order coefficient. We denote the predictions which are obtained by extrapolating the coefficients and the “prediction errors” (in contrast to the coefficients alone) by the term a-posteriori improved predictions because the further correction due to the extrapolated error is applied after the reexpansion of the rational approximant which yields the “first-order” prediction. The a-posteriori improvement of predictions is useful in both lower and higher orders of perturbation theory. In higher orders, the transient, pre-asymptotic behavior of the perturbative coefficients has died away, and the extrapolations of the coefficients as well as the a-posterori corrections become more accurate. A particular merit of the a-posteriori corrections is the fact that they can be applied to any of the prediction algorithms proposed above, in order to achieve an improvement of the prediction beyond the reexpansion of the rational approximant used.
We will be investigating in the sequel the Borel transform of the QED effective action defined in (43). In order to investigate the extrapolation of coefficients of the Borel transform, we define auxiliary quantities (coefficients) $`\widehat{c}_{2n+1}`$ by the relation
$$\widehat{c}_{2n+1}=\frac{16c_n}{(2n+1)!},$$
(60)
where the $`c_n`$ are defined in Eq. (35). We additionally set $`\widehat{c}_{2n}=0`$ for even-order coefficients. In terms of the coefficients $`c_j(p)`$ which are defined in Eq. (5) in , the $`\widehat{c}_{2n+1}`$ are given by $`\widehat{c}_{2n+1}=(1)^nc_{2n+1}(0)`$. The coefficient $`\widehat{c}_{2n+1}`$, written in terms of the Bernoulli numbers, reads \[see also Eq. (35)\]
$$\widehat{c}_{2n+1}=(1)^n\frac{2^{2n+4}|_{2n+4}|}{(2n+4)!}.$$
(61)
We concentrate here on the coefficient $`\widehat{c}_{13}`$ defined in Eq. (60) and we discuss how the prediction for the coefficient $`\widehat{c}_{13}`$ can be improved on the basis of a-posteriori corrections and other asymptotic improvements. We define correction factors $`\xi _n`$ by
$$\widehat{c}_n=\xi _n\overline{c}_n$$
(62)
where $`\widehat{c}_n`$ is the exact $`n`$th order coefficient and $`\overline{c}_n`$ is the estimate obtained by reexpanding the rational approximant which is used for the prediction. In the case of Borel-Padé approximants, these would be Padé approximants applied to the Borel transform of the QED effective action (43). Specifically, for the prediction of the $`n`$th perturbative coefficient, these are the approximants $`[[(n1)/2]]/[[n/2]]`$ for the lower-diagonal and $`[[n/2]]/[[(n1)/2]]`$ for the upper-diagonal case. For the prediction of $`\widehat{c}_{13}`$, the previous errors made in the “prediction” of $`\widehat{c}_7`$, $`\widehat{c}_9`$ and $`\widehat{c}_{11}`$ can be analyzed. Note that the exact values of $`\widehat{c}_7`$, $`\widehat{c}_9`$ and $`\widehat{c}_{11}`$ must be assumed as available, exploitable information by the time we try to predict $`\widehat{c}_{13}`$.
An estimate for $`\xi _{13}`$ can be obtained for example by fitting the natural logarithms of the quantities $`\xi _71`$, $`\xi _91`$ and $`\xi _{11}1`$ with a linear model in order to obtain an estimate for $`\xi _{13}1`$. This linear fit of the logarithms of the $`\xi _i`$ is based on the empirical observation that relative errors of the predictions decrease exponentially in higher order, a phenomenon which has been observed in a number of applications, including variants of anharmonic oscillators. In the context of Padé approximants, a similar error dependence has been conjectured (see ). The leading coefficient of the decay of the relative errors may depend on the problem considered and on the extrapolation scheme used, but the exponential improvement of perturbative predictions in higher order appears to be a rather general feature. Details on this point will be presented elsewhere .
Using a linear least-squares fit of the respective logarithms $`\mathrm{ln}(\xi _71)`$, $`\mathrm{ln}(\xi _91)`$ and $`\mathrm{ln}(\xi _{11}1)`$, we obtain an estimate of $`\xi _{13}1=2.47\times 10^6`$ in the case of upper-diagonal Borel-Padé approximants. This leads to the data presented in Table 2. Note that even the crude linear model for the $`\mathrm{ln}(\xi _i1)`$ used here already doubles the accuracy of the prediction of the coefficient $`\widehat{c}_{13}`$ as compared to the plain Borel-Padé prediction used, for example, in . Note also that an averaging of upper-and lower-diagonal Borel-Padé approximants does not improve the situation in favor of the plain Borel-Padé extrapolation. Further improvements of the a-posteriori corrected predictions is possible with more elaborate extrapolation schemes .
We would like to mention that knowledge of the leading asymptotics of the perturbative coefficients is required for the construction of an asymptotically optimized Borel-Padé transformation; this information is not available for many of the phenomenologically interesting series currently investigated . It is helpful to note that the delta transformation is (like Padé) a rather general-purpose method for the prediction of perturbative coefficients, and that knowledge of large-order asymptotics is not required for its construction or application in a particular case. This property is helpful especially in cases of practical interest where little is rigorously known about the large-order asymptotics of the perturbative coefficients, and where only a limited number of perturbative coefficients are available. A number of practically interesting examples of delta-based predictions were discussed in , and it was observed that the delta transformation yields more accurate predictions than the Padé technique in many cases.
We now turn to a discussion of topologically new effects in higher orders of perturbation theory and improvements of predictions based on the renormalization group. Topologically new effects have caused problems for perturbative predictions in the past. We refer to the quartic Casimirs in the QCD beta function and to light-by-light scattering graphs in the tenth order anomalous magnetic moment of the muon . Analogous considerations might hold for the perturbation series investigated in . The topologically new effects cannot be taken into account by straight extrapolations, nor by renormalization-group improved extrapolations, which lead to resummation of certain classes of diagrams. For the anomalous magnetic moment of the muon discussed in , the contribution of the topologically new light-by-light scattering diagrams originally analyzed in could not be reproduced by renormalization-group techniques introduced in . By reexpansion of the delta approximant to the perturbation series for the muon anomaly, an estimate of $`a_\mu ^{(10)}=711`$ has been obtained . If an a-posteriori correction based on a combination of the delta transformation and Padé approximants is added to this prediction, then the estimate for the 10th order coefficient changes to $`a_\mu ^{(10)}970`$ . This improved estimate is in excellent agreement with the analytically obtained approximate result of $`930(170)`$ from , comprising the topogically new effects which are present in five-loop order.
Now it is of course not permissible to conclude that topologically new effects can be taken into account in the general case by considering a-posteriori corrections, and/or that a-posteriori corrections necessarily give an accurate estimate for the size of these problematic, topologically new effects. On the other hand, the reverse proposition, which is that perturbative predictions are scientifically unsound because of their general inability to include topologically new effects, does not appear to be generally valid, either.
## 5 THE ANOMALOUS DIMENSION
We consider the divergent perturbation series for the $`\gamma `$ function (anomalous dimension) of the Yukawa coupling as studied in . Specifically, we consider the resummation of the perturbation series for the anomalous dimension of a fermion field with a Yukawa interaction $`g\overline{\psi }\sigma \psi `$ at $`d_c=4`$, which is given in Eq. (17) in . This calculation comprises an evaluation of the contribution of all nested self-energy diagrams to the anomalous dimension $`\gamma `$ function of the Yukawa theory up to the 30-loop level (an analogous analysis is performed in for the $`(4ϵ)`$–dimensional $`\varphi ^4`$ theory). We restrict the discussion here to the Yukawa case. With the convention
$$a=\frac{g^2}{4\pi ^2},$$
(63)
the result for the anomalous dimension $`\gamma `$ function as considered in reads,
$$\stackrel{~}{\gamma }_{\mathrm{hopf}}(a)\underset{n=1}{\overset{\mathrm{}}{}}(1)^n\frac{\stackrel{~}{G}_n}{2^{2n1}}a^n.$$
(64)
The perturbative coefficients $`\stackrel{~}{G}_n`$ are listed in Table 3. The coefficients grow factorially in absolute magnitude; in little change is observed in the quantities
$$\stackrel{~}{S}_n=\frac{\stackrel{~}{G}_n}{2^{n1}\mathrm{\Gamma }(n+1/2)}$$
(65)
for large $`n`$. The evaluation confirms in a concrete, 30-loop calculation the assumption originally put forth by Dyson that the convergence radius of the quantum field theoretic perturbative expansion is zero. For a large number of quantum field theoretic observables like the anomalous magnetic moment of the muon (see Section 4) only a few perturbative terms are known. Although rapid growth of the perturbative coefficients is observed even in relatively low order (see the large number of examples discussed in ), one may argue that the factorial growth of the coefficients has not been demonstrated in concrete calculations, and that it is unclear if severe cancellations between different classes of diagrams occur in higher order.
The 30-loop calculation by Broadhurst and Kreimer may indicate in one particular example at least, that the factorial divergence is likely to persist, and that possible cancellations due to the renormalization or between different sets of diagrams do not contradict the concept of ultimate factorial divergence of the perturbative coefficients. From the observation made in that the $`S_n`$ defined in Eq. (65) change little at large $`n`$, one may tentatively infer the leading factorial divergence of the perturbative coefficients,
$$\stackrel{~}{G}_n2^{n1}\mathrm{\Gamma }(n+1/2)n\mathrm{}.$$
(66)
This asymptotic behavior, of course, leads to a vanishing radius of convergence of the perturbative expansion (64).
As observed by Broadhurst and Kreimer, the perturbation series (64) can be resummed with the help of an asymptotically improved Borel-Padé technique. From Eq. (22) in , it is clear that a $`(1,1/2)`$-generalized Borel-Padé transformation is used by Broadhurst and Kreimer \[for the definition of generalized Borel-Padé transformations see Eqs. (2)–(4) in this article\]. This is not completely obvious because the transformation as used by Broadhurst and Kreimer has been modified additionally such as to normalize the first coefficient of the Borel transform (not of the input series) to unity, and the transformation is additionally rewritten such as to reflect the vanishing coefficient of zeroth order in $`a`$ in Eq. (64). From the leading asymptotics in Eq. (66), the singularity of the $`(1,1/2)`$-generalized Borel transform closest to the origin can be inferred. This singularity was explicitly “put in by hand” by Broadhurst and Kreimer (see Eq. (22) in ).
It is also possible to resum the alternating divergent series (64) by a delta transformation, even at large coupling. At a large Yukawa coupling of $`g=30`$, we obtain a relative accuracy of 6 significant figures in the resummed results with a plain, unmodified delta transformation.
We add here a remark on the relation of the asymptotically optimized Borel-Padé based predictions to those obtained using the delta transformation. We consider the relative accuracy of perturbative predictions for the coefficient $`\stackrel{~}{G}_{30}`$ of the perturbation series defined in Eq. (64) using various methods. The coefficient $`\stackrel{~}{G}_{30}`$ is known (see Table 3), therefore we merely check the accuracy to which this coefficient can be reproduced by considering the first 29 perturbative coefficients of the series (64). With an asymptotically optimized $`(1,1/2)`$-generalized Borel-Padé technique, the coefficient $`\stackrel{~}{G}_{30}`$ can be reproduced with a relative accuracy of $`5\times 10^{16}`$. The delta transformation, without any modifications, leads to a prediction with a relative error of $`3\times 10^{16}`$; this result is more accurate than the prediction provided by the asymptotically optimized Borel-Padé transformation. In accordance with the results of Section 4, the Borel-Padé transformation can be significantly enhanced by including the pole closest to the origin (see Eq. (66) above and Eq. (22) in ). When this pole is included, a prediction is obtained with a relative error of $`4\times 10^{17}`$. We do not consider a-posteriori improvements to either of these predictions, here.
## 6 CONCLUSION
We have considered the resummation of divergent perturbation series and the prediction of unknown higher-order perturbative coefficients (perturbative predictions or perturbative extrapolations). We have mentioned and discussed the following resummation prescriptions,
* the direct application of Padé approximants to a divergent series,
* the direct application of the nonlinear (delta) sequence transformation,
* asymptotically improved variants of the Borel-Padé technique.
The direct application of Padé approximants is less efficient in the resummation of divergent perturbation series than both the delta transformation and the combined Borel and Padé techniques. The combined, asymptotically improved Borel and Padé techniques, and the delta transformation are complementary. On the one hand, it can hardly be overemphasized that the asymptotically improved Borel-Padé technique is less general than the delta transformation because it depends on the availability of information on the leading large-order asymptotics of the coefficients. By contrast, there is considerable evidence that the plain, unmodified delta transformation can sum factorially divergent alternating series which diverge as strongly as $`(3n)!`$ . This is beyond the power of directly applied Padé approximants and also beyond the power of the $`(1,1)`$-generalized Borel-Padé transformation (“usual”) Borel-Padé transformation defined in Eq. (4).
If additional information is available on the input series, then the asymptotically optimized Borel-Padé technique is rather attractive. As discussed in Section 3, it is possible to enhance the rate of convergence simply by utilizing the location of known poles in the Padé approximants to the Borel transform of the input series. These improvements are not restricted to the first UV and IR renormalon poles, but, as shown in Eqs. (53)–(55) and exemplified by the numerical results in Table 1, can take advantage of an in principle unlimited number of poles in the Borel plane. Other techniques for possible improvements of the Borel-Padé algorithm have been described in . Note that it is also possible to generalize the Borel-Padé technique to those cases where there are poles along the positive real axis in which case the Borel integral in Eq. (2) is actually undefined . Using the special integration contours in , it is even possible to derive nonperturbative imaginary parts from real, not complex, perturbative coefficients. Concerning nonperturbative effects in quantum field theory we also refer to the recent investigation .
As discussed in Section 3, it is possible to accelerate the convergence of the resulting Borel-Padé transforms by subsequent application of Wynn’s epsilon algorithm (Borel-Padé-Wynn technique). These techniques lead to an improved rate of convergence. Note that the use of explicit information of the location of the poles in the Borel plane and the subsequent improvement of the convergence of the transforms by Wynn’s epsilon algorithm should also lead to accelerated convergence in the case of the complex integrations discussed in .
Because all the resummation prescriptions discussed above fulfill accuracy-through-order relations, they can be used to predict perturbative coefficients (we refer to this procedure as perturbative predictions or perturbative extrapolations). The straightforward predictions are obtained by reexpansion of the rational approximant in powers of the coupling parameter. That is to say, we consider here perturbative predictions based on
* the reexpansion of Padé approximants directly applied to the perturbative (input) series,
* the reexpansion of nonlinear (delta) sequence transformations directly applied to the perturbative (input) series,
* and the reexpansion of Padé approximants applied to the asymptotically improved Borel transform of the input series.
As it has been demonstrated in and , the predictions based on the delta transformation and on the combined Borel and Padé techniques yield better results for the perturbative coefficients of the QED effective action than the Padé approximants alone. In we also presented a number of more realistic and practically interesting examples in which the delta transformation leads to better predictions than the Padé approximants.
Note that there is currently no general proof of the assumption that the renormalon poles are the only relevant poles in the Borel plane , but the factorial divergence of the perturbative coefficients is a rather commonly accepted assumption . We should therefore assume that at least asymptotically, the perturbation series in quantum field theory approximate factorially divergent series. This is also confirmed by the concrete 30-loop calculation presented in . For many factorially divergent series, the delta transformation produces better numerical results than Padé approximants (see, e.g., Ch. 13 of and ). Therefore, the delta transformation can be expected to provide a competitive alternative to Padé approximants. We again refer to the large number of recent publications on Padé-based extrapolations in quantum field theory .
We consider here mainly the following asymptotic improvements of perturbative extrapolations
* a-posteriori corrections based on estimates not only for the next higher-order coefficients, but also for the error which is to be expected in the estimation of that coefficient and
* the use of known (renormalon) poles in order to fix the denominatior structure of Padé approximants in the context of the Borel-Padé method.
The general idea of a-posteriori corrections is the following. In higher orders of perturbation theory a number of lower-order coefficients are available which may be used in order to construct rational approximants. These coefficients can, apart from being useful for the construction the approximants itself, also be utilized in order to obtain an estimate for the expected error in the perturbative prediction. To this end, the extrapolation procedure is applied to the known lower-order terms in the perturbation series. A comparison of the “predictions” for the known lower-order with their exact results gives lower-order correction factors which may be extrapolated to higher order. This immediately leads to a correction factor for the next higher-order perturbative extrapolation (see Section 4).
Using a-posteriori corrections and the pole structure, we obtain improved results for the perturbative predictions of the model problem studied in (see Table 2). The a-posteriori corrections also lead to an improved estimate for the 10th order anomalous magnetic moment of the muon and bring the a-posteriori corrected prediction in close agreement with an analytically obtained estimate . The renormalization-group analysis can lead to a resummation of certain classes of Feynman diagrams, but it does not lead to an understanding of topologically new effects which occur in higher orders of perturbation theory. This phenomenon has lead to problems in perturbative predictions in the past, especially in those cases where these predictions were improved on the basis of a renormalization group analysis. We refer to the analysis by Kataev and Starshenko on the 10th order anomalous magnetic moment of the muon and various investigations on the QCD beta function . Notably, as discussed in Section 4, the a-posteriori corrected prediction appears to be consistent with the topologically new effects observed in 10th order of perturbation theory and calculated approximately in .
We would like to stress here again that the concept of a-posteriori corrections is rather general and can be applied to all prediction algorithms mentioned above. Specifically, we refer to the investigation on significant improvements which can be achieved in the context of Padé-based predictions with this technique. The reduction of the magnitude of the a-posteriori correction is an attractive feature of the predictions based on the delta transformation.
In Sections 5 we show that for the more realistic 30-loop series calculated by Broadhurst and Kreimer in , even the asymptotically optimized Borel-Padé technique cannot quite match the accuracy obtainable by the plain delta transformation. It is only when an additional pole is explicitly put in by hand in the Padé transformations that the combined, asymptotically improved Borel-Padé technique becomes more accurate than the plain delta transformation. We stress here that the construction of the asymptotically improved Borel-Padé technique requires in itself a knowledge of the large-order asymptotics of the perturbative coefficients. Such information is not available in general cases. Specifically, in those cases where only a small number of coefficients are known the leading asymptotics cannot be reliably inferred from empirical approaches, either.
It has been the purpose of this article to clarify how resummation algorithms and perturbative extrapolations can be improved if additional information is available on a particular input series. We have explained in Sections 35 various algorithms by which the resummation of divergent series and the prediction of perturbative coefficients can be improved on the basis of additionally available asymptotic information on a given series; these improvements can be applied to the Borel-Padé based techniques and to the delta transformation techniques .
## ACKNOWLEDGMENTS
U.J. acknowldges helpful conversations with P.J. Mohr. E.J.W. acknowledges support from the Fonds der Chemischen Industrie. G.S. acknowledges continued support from BMBF, GSI and DFG.
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# Rotating Bianchi type 𝑉 dust models generalizing the 𝑘=-1 Friedmann model
## 1 Statement of the problem and summary of the paper.
This paper is a continuation of a series of papers on rotating dust models in relativity<sup>1-3</sup>. The initial motivation for this research was the desire to find a rotating generalization of the Friedmann models. In spite of much effort spent on investigating solutions of Einstein equations with a rotating matter source, no such generalization has been found so far; see literature surveys in Refs. 3 and 4. Refs. 1, 2 and 3 provided a complete classification scheme for hypersurface-homogeneous rotating perfect fluid models with zero acceleration. Unlike in previous approaches, nothing was assumed about the position of the symmetry orbits in spacetime; the classification includes also timelike and null orbits, and so it is the farthest-reaching application of the Bianchi classification to rotating and nonaccelerating perfect fluid models in relativity. The models split into 3 general classes: I, in which two of the Killing fields are everywhere spanned on the vector fields of velocity $`u^\alpha `$ and rotation $`w^\alpha (`$Ref. 1); II, in which only one Killing field is spanned on $`u^\alpha `$ and $`w^\alpha (`$Ref. 2); and III, in which all Killing fields are linearly independent of $`u^\alpha `$ and $`w^\alpha (`$Ref. 3). The many particular cases arise because of several possible alignments or misalignments among the 3 Killing fields and $`u^\alpha `$ and $`w^\alpha `$.
By the Bianchi type of the symmetry algebra and by the relation of the velocity field to the symmetry orbits it can be recognized in which cases generalizations of the Friedmann models can be expected. Two such candidate cases were found in class II, and five more in class III. Those of class III were prohibitively complicated, but one of the cases of class II allowed for some progress, and this one is presented in the present paper. It is the Bianchi type $`V`$ subcase of the case 1.2.2.2, given by eq. (5.19) in Ref. 2. The other candidate case found in class II, eq. (5.10) in Ref. 2, can reproduce only the de Sitter or the Einstein model in the limit of zero rotation, this is seen from the time-dependence of the metric. Hence, it is not interesting for cosmology and therefore disregarded here. In sec. 2, the metric is simplified by a coordinate transformation, and a first integral of the Einstein equations is found. With zero value of this integral, coordinate transformations can be used to eliminate two components of the metric tensor, and the number of nontrivial Einstein equations is reduced to 7. Although there are only 4 functions + matter density to be determined by these 7 equations, the set later turns out to be self-consistent. In sec. 3, it is shown that the $`k=1`$ Friedmann models are contained among the solutions that result in the limit of zero rotation. In sec. 4, the Einstein equations are reduced to a set $`S`$ of 3 second-order equations to determine 3 metric components + a quadrature $`Q`$ to determine the fourth component ($`g_{33}`$). Of the Einstein equations derived in sec. 2, one is fulfilled identically in consequence of the set $`\{SQ\}`$, one turns out to be a constraint imposed on the initial data, and the one that determines the matter-density turns out to provide a first integral. The constraint and the first integral are second-degree polynomials in the first derivatives of the unknown functions whose coefficients depend on the unknown functions. The first integral determines $`g_{33}`$ algebraically in terms of the other components, and so it is a replacement for the quadrature $`Q`$. It is also shown that the set $`S`$ cannot be obtained as the Euler-Lagrange equations from a variational principle of the Hilbert type. Finally, it is shown in sec. 4 how the set $`\{SQ\}`$ reproduces the Friedmann equations in the limit of zero rotation and zero shear. In sec. 5, Lie point-symmetries of the set are found: there is a two-dimensional symmetry group that allows one to reduce one second-order equation to a first-order equation plus a quadrature. However, this reduction provides no real progress toward solving the set $`S`$; the first-order equation is still a member of a complicated set. In sec. 6, a method of systematic search for polynomial first-order first integrals of a set of ordinary differential equations is applied to the set $`S`$ of sec. 4. It is shown that no first integrals that are polynomials of degree 1 or 2 in the first derivatives exist. The same method is used to reveal the existence of a possible constraint on initial data, which is of degree 1 in first derivatives, that is preserved by the set $`S`$. However, the constraint necessarily implies zero matter-density, and so it is not interesting for cosmology.
Calculations that are of secondary importance for the main text, but are difficult to reproduce, are described in the appendices.
## 2 The Einstein equations, their first integral and implications of the zero value of this integral.
The subject of the present paper are the Einstein equations for the Bianchi type V subcase of case 1.2.2.2 of Ref. 2. For reference, the initial formulae are recalled in their original notation.
The Bianchi type V symmetry results when $`c=0`$ in eqs. (5.19) of Ref. 2 and when, in addition, $`j=a`$ in eqs. (5.16). Hence, the metric is:
$$\mathrm{d}s^2=\mathrm{d}t^2+2y\mathrm{d}t\mathrm{d}x+y^2h_{11}\mathrm{d}x^2+2h_{12}\mathrm{d}x\mathrm{d}y+2y^2h_{13}\mathrm{d}x\mathrm{d}z$$
$$+(h_{22}/y^2)\mathrm{d}y^2+2h_{23}\mathrm{d}y\mathrm{d}z+y^2h_{33}\mathrm{d}z^2,$$
(2.1)
where the coordinates are $`\{x^\alpha \}=\{x^0,x^1,x^2,x^3\}=\{t,x,y,z\}`$, and $`h_{ij},i,j=1,2,3`$ are unknown functions of the variable
$$v=\mathrm{e}^ty^{C_2/a},$$
(2.2)
$`a`$ and $`C_2`$ being arbitrary constants. The velocity field $`u^\alpha `$, the rotation field $`w^\alpha `$ and the Killing fields $`k_{(i)}^{}{}_{}{}^{\alpha }`$, $`i=1,2,3`$ are given by:
$$u^\alpha =\delta _{}^{\alpha }{}_{0}{}^{},w^\alpha =(\rho /y)\delta _{}^{\alpha }{}_{0}{}^{},k_{(1)}^{}{}_{}{}^{\alpha }=\delta _{}^{\alpha }{}_{1}{}^{},k_{(3)}^{}{}_{}{}^{\alpha }=\delta _{}^{\alpha }{}_{3}{}^{},$$
$$k_{(2)}^{}{}_{}{}^{\alpha }=C_2\delta _{}^{\alpha }{}_{0}{}^{}+a(x\delta _{}^{\alpha }{}_{1}{}^{}y\delta _{}^{\alpha }{}_{2}{}^{}+z\delta _{}^{\alpha }{}_{3}{}^{}),$$
(2.3)
where $`\rho `$ is the matter-density of dust. The rotation tensor $`\omega _{\alpha \beta }`$ has only one algebraically independent nonzero component:
$$\omega _{12}=\frac{1}{2},$$
(2.4)
and therefore the coordinates used here are ill-suited for considering the limit $`\omega 0`$. From the first equation in (2.3) it can be seen that the coordinates are comoving.
As shown in Ref. 1, it follows from the equations of motion and from the equation of conservation of the number of particles that:
$$g:=\mathrm{det}(g_{\alpha \beta })=(y/\rho )^2,$$
(2.5)
where $`\rho `$ is the mass-density.
This is the form in which the metric resulted from the Killing equations in Ref. 2. It is advantageous to transform the coordinates as follows:
$$t=t^{}(C_2/a)\mathrm{ln}y^{},x=x^{}C_2/(ay^{}),(y,z)=(y^{},z^{}).$$
(2.6)
The result is equivalent to substituting $`C_2=0`$ and $`a=1`$ in eqs. (2.1) - (2.4), i.e. the forms of the metric (2.1), of the vector fields $`u^\alpha `$, $`w^\alpha `$, $`k_{(1)}^{}{}_{}{}^{\alpha }`$ and $`k_{(3)}^{}{}_{}{}^{\alpha }`$ in (2.3) and of the rotation tensor $`\omega _{\alpha \beta }`$ in (2.4) do not change (although the new $`h_{ij}^{}`$ in (2.1) will be linear combinations of the old $`h_{ij}`$), while the new $`k_{(2)}^{}{}_{}{}^{\alpha }`$ basis vector will be:
$$k_{(2)}^{}{}_{}{}^{\alpha }=x\delta _{}^{\alpha }{}_{1}{}^{}y\delta _{}^{\alpha }{}_{2}{}^{}+z\delta _{}^{\alpha }{}_{3}{}^{},$$
(2.7)
and the argument of $`h_{ij}`$ will now be $`v=\mathrm{e}^t^{}`$, i.e. the $`h_{ij}`$ are from now on unknown functions of the time-coordinate $`t`$.
The isometry corresponding to (2.7) is:
$$t^{}=t,(x^{},z^{})=\mathrm{e}^\tau (x,z),y^{}=\mathrm{e}^\tau y,$$
(2.8)
where $`\tau `$ is the group parameter.
It is convenient to parametrize the metric as follows:
$$\mathrm{d}s^2=(\mathrm{d}t+y\mathrm{d}x)^2(yK_{11}\mathrm{d}x)^2(K/y)^2(\mathrm{d}y+y^2h\mathrm{d}x)^2K_{33}^{}{}_{}{}^{2}[yg\mathrm{d}x+(f/y)\mathrm{d}y+y\mathrm{d}z]^2,$$
(2.9)
where $`K_{11}`$, $`K`$, $`K_{33}`$, $`h`$, $`f`$, and $`g`$ are unknown functions of $`t`$. The components of the Einstein tensor referred to below are tetrad components $`G_{IJ}=e_{}^{\alpha }{}_{I}{}^{}e_{}^{\beta }{}_{J}{}^{}G_{\alpha \beta }`$, i.e. projections of the coordinate components $`G_{\alpha \beta }`$ onto the orthonormal tetrad $`e^I:=e_{}^{I}{}_{\alpha }{}^{}\mathrm{d}x^\alpha `$ implied by (2.9):
$$e^0=\mathrm{d}t+y\mathrm{d}x,e^1=yK_{11}\mathrm{d}x,e^2=(K/y)(\mathrm{d}y+y^2h\mathrm{d}x),$$
$$e^3=K_{33}[yg\mathrm{d}x+(f/y)\mathrm{d}y+y\mathrm{d}z],$$
(2.10)
where $`e_{}^{\alpha }{}_{I}{}^{}`$ is the inverse matrix to $`e_{}^{I}{}_{\alpha }{}^{}`$, i.e. $`e_{}^{\alpha }{}_{I}{}^{}e_{}^{I}{}_{\beta }{}^{}=\delta _{}^{\alpha }{}_{\beta }{}^{}`$, $`e_{}^{\alpha }{}_{J}{}^{}e_{}^{I}{}_{\alpha }{}^{}=\delta _{}^{I}{}_{J}{}^{}`$. In the parametrization (2.9), the determinant of the metric is:
$$g=(yK_{11}KK_{33})^2.$$
(2.11)
The tetrad components of the Einstein tensor corresponding to the metric (2.9) are given in the Appendix A. As seen from there, two combinations of those equations are of first order, they are $`K_{11}G_{03}+G_{13}=0`$, i.e.:
$$(\frac{3}{2}K_{33}/K_{11})[(K_{11}^{}{}_{}{}^{2}1)K^2f,_t+h(hf,_tg,_t)]=0$$
(2.12)
and $`K_{11}G_{02}+G_{12}=0`$, i.e.:
$$(K_{11}K)^1[\frac{3}{2}K^2hh,_t+\frac{1}{2}hK_{11}K_{11},_t+(K_{11}^{}{}_{}{}^{2}1)(2K,_t/KK_{33},_t/K_{33})]=0.$$
(2.13)
As shown in Appendix B, the case $`h=0`$ does not lead to interesting developments, so we shall proceed further under the assumption:
$$h0.$$
(2.14)
Then, eq. (2.12) implies:
$$g,_t=[h+(K_{11}^{}{}_{}{}^{2}1)/(hK^2)]f,_t$$
(2.15)
With this, the equations $`G_{03}=G_{13}=G_{23}=0`$ turn out to be equivalent, and they can be written as follows:
$$\frac{1}{2}(\frac{K_{11}^{}{}_{}{}^{2}1}{h}\frac{K_{33}^{}{}_{}{}^{3}f,_t}{K_{11}K}),_t+\frac{K_{33}^{}{}_{}{}^{3}f,_t}{K_{11}K}=0.$$
(2.16)
This invites the introduction of the new variable $`u(t)`$ by $`u,_t=h/(K_{11}^{}{}_{}{}^{2}1)`$, and then (2.16) becomes:
$$\left(\frac{K_{33}^{}{}_{}{}^{3}f,_u}{K_{11}K}\right),_u2\frac{K_{33}^{}{}_{}{}^{3}f,_u}{K_{11}K}=0,$$
(2.17)
which has the first integral $`K_{33}^{}{}_{}{}^{3}f,_u/(K_{11}K)=C\mathrm{e}^{2u}`$, $`C=`$ const, i.e.:
$$f,_t=C\mathrm{e}^{2u}hK_{11}K/[K_{33}^{}{}_{}{}^{3}(K_{11}^{}{}_{}{}^{2}1)].$$
(2.18)
From here on, we shall follow only the special case $`C=0`$, which is a solution of the Einstein equations, but not a general one: it is a subcase chosen ad hoc for further progress with integration. Then, from (2.18) and (2.15) $`f=`$ const, $`g=`$ const, and from (2.10) the coordinate transformation $`z^{}=z+f/y+gx`$ leads to
$$f=g=0$$
(2.19)
without changing any of the other formulae for $`g_{\alpha \beta }`$, $`u^\alpha `$, $`w^\alpha `$, $`\omega _{\alpha \beta }`$ or $`k_{(i)}^{}{}_{}{}^{\alpha }`$.
The Einstein equations $`G_{03}=G_{13}=G_{23}=0`$ are now fulfilled identically. We are left with 7 equations of the set (A.1) – (A.10) in Appendix A that should determine the 4 functions $`K_{11}`$, $`K`$, $`K_{33}`$ and $`h`$, and the matter density $`\rho `$ in addition. It will turn out in sec. 4 that the 7 equations are dependent just in the way needed to make the problem self-consistent and determinate.
## 3 The Friedmann limit of the metric.
As already stated, the coordinates used in sec. 2 are ill-suited for considering the limit $`\omega 0`$. It will be shown in the present section that this limit can be calculated after a coordinate transformation and a reparametrization of the metric. This is just a demonstration of existence, and it is not claimed that the limit $`\omega 0`$ thus obtained is unique (i.e. another nonrotating limit might be obtained starting from a different coordinate transformation). However, we will be satisfied to show that a limit exists in which the $`k<0`$ Friedmann model is contained.
Since $`\omega _{12}=\omega _{21}=\frac{1}{2}`$ are the only nonzero components of the rotation tensor, a natural coordinate transformation to consider is:
$$y=\omega _0y^{}.$$
(3.1)
where $`\omega _0`$ is a constant. After the transformation:
$$\omega _{12}^{}=\frac{1}{2}\omega _0=\omega _{21}^{}$$
(3.2)
(all other $`\omega _{\alpha \beta }=0`$), and the limit of zero rotation is $`\omega _00`$. However, before this limit is taken, the metric functions in (2.9) must be reparametrized or else the limit will be singular. The following reparametrizations will do the job:
$$K_{11}=\stackrel{~}{K}_{11}/\omega _0,K_{33}=\stackrel{~}{K}_{33}/\omega _0,f=\stackrel{~}{f}\omega _0.$$
(3.3)
The transformation (3.1) and the reparatmetrization (3.3) result in the following metric:
$$\mathrm{d}s^2=(\mathrm{d}t+\omega _0y^{}\mathrm{d}x)^2(y^{}\stackrel{~}{K}_{11}\mathrm{d}x)^2K^2(\mathrm{d}y^{}/y^{}+\omega _0y^{}h\mathrm{d}x)^2$$
$$\stackrel{~}{K}_{33}^{}{}_{}{}^{2}[y^{}g\mathrm{d}x+(\stackrel{~}{f}/y^{})\mathrm{d}y^{}+y^{}\mathrm{d}z]^2$$
(3.4)
whose limit $`\omega _00`$ (with primes and tildes omitted) is:
$$\mathrm{d}s^2=\mathrm{d}t^2(yK_{11}\mathrm{d}x)^2(K/y)^2\mathrm{d}y^2K_{33}^{}{}_{}{}^{2}[yg\mathrm{d}x+(f/y)\mathrm{d}y+y\mathrm{d}z]^2,$$
(3.5)
This is more than sufficiently general to accomodate the $`k=1`$ Friedmann model that results when $`g=f=0`$ and $`K_{11}=K=K_{33}:=R(t)`$, where $`R(t)`$ is the Friedmann scale factor. The resulting coordinates are none of the standard ones, but are related by $`y=\mathrm{e}^u`$ to one of the sets used in the literature (see eq. (1.3.15) in Ref. 5).
The fact that (3.5), the limit $`\omega _00`$ of (2.9), is still more general than the Friedmann metric means that (3.5) has nonzero shear, i.e. shear survives the transition $`\omega 0`$.
However, one possible problem still lies ahead. It was proven above that the $`k=1`$ Friedmann model is contained among the solutions of the set (A.1) – (A.10). What is still needed is an explicit solution with the property that it has nonzero rotation in general, but reproduces the $`k=1`$ Friedmann model in the limit $`\omega 0`$. Experience with the Einstein equations in other cases shows that sometimes, while integrating the equations, one encounters mutually exclusive alternatives $`A`$ and $`B`$ such that it is no longer possible to recover $`B`$ as a limit of $`A`$ after the integration is completed. A well-known example are the two subfamilies ($`\beta ^{}=0`$ and $`\beta ^{}0`$) of the Szekeres-Szafron<sup>6-7</sup> cosmological models; see Ref. 5 for more on this point. (Only recently, a reformulation of the two classes was invented that allows to recover the $`\beta ^{}=0`$ family from the other one, see Ref. 8). Hence, it may still happen that among the explicit solutions, the rotating dust model and the Friedmann $`k=1`$ model will turn out to be mutually exclusive subfamilies. This uncertainty will persist until an explicit solution is found.
It will be shown at the end of sec. 4 that the explicitly written out Einstein equations do allow a continuous limiting transition $`\omega 0,\sigma 0`$, and in the limit they reproduce exactly the Friedmann equations.
## 4 The independent Einstein equations.
We shall now proceed with the subcase (2.19). Eq. (2.12) is then fulfilled identically. Eq. (2.13) does not change, and it can be more conveniently rewritten if $`K_{11}`$ is parametrized as follows:
$$K_{11}=\mathrm{cosh}(F).$$
(4.1)
Then, from (2.13):
$$K_{33},_t=K_{33}[\frac{3}{2}K^2hh,_t/\mathrm{sinh}^2(F)+\frac{1}{2}h/\mathrm{sinh}^2(F)+2K,_t/K\mathrm{cosh}(F)F,_t/\mathrm{sinh}(F)].$$
(4.2)
When this is substituted into the remaining equations (A.1) – (A.10), the function $`K_{33}`$ disappears from the set completely, i.e. we are left with 6 equations to determine $`h`$, $`K`$, $`F`$ and the matter-density plus the quadrature implied by (4.2) that allows one to calculate $`K_{33}`$ once $`h(t)`$, $`K(t)`$ and $`F(t)`$ are known.
Since (2.13) is now satisfied, the equations $`G_{02}=0`$ and $`G_{12}=0`$ are equivalent, and they can be written as:
$$h,_{tt}=\frac{3}{2}K^2hh,_{t}^{}{}_{}{}^{2}/\mathrm{sinh}^2(F)5K,_th,_t/K+(2\mathrm{cosh}^2(F)1)F,_th,_t/\mathrm{sinh}(F)\mathrm{cosh}(F)$$
$$+hh,_t/\mathrm{sinh}^2(F)+K,_t/K^3+F,_t/K^2\mathrm{cosh}(F)\mathrm{sinh}(F)\frac{1}{2}h/(K\mathrm{sinh}(F))^2$$
(4.3)
This is used to eliminate $`h,_{tt}`$ from the other Einstein equations. The equation $`G_{01}=0`$ can then be solved for $`F,_{tt}`$ (the solution is given in Appendix C) and this is used to eliminate $`F,_{tt}`$ from the diagonal components of the Einstein tensor (all the non-diagonal Einstein equations have been used up at this point). After such a substitution, the following identity is fulfilled:
$$G_{11}+G_{33}2G_{22}0,$$
(4.4)
i.e. one of the three equations $`G_{11}=G_{22}=G_{33}=\mathrm{\Lambda }`$ can be discarded because it is a consequence of the remaining two. We choose to discard $`G_{33}=\mathrm{\Lambda }`$.
Then, $`K,_{tt}`$ can be calculated from $`G_{22}G_{11}=0`$. The result is:
$$K,_{tt}=\frac{1}{4}K^3\mathrm{sinh}^2(F)h,_{t}^{}{}_{}{}^{2}\frac{3}{2}K^3h\mathrm{cosh}(F)\mathrm{sinh}^3(F)F,_th,_t$$
$$\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)K,_tF,_t+2\mathrm{cosh}(F)\mathrm{sinh}^1(F)K,_tF,_tK\mathrm{cosh}^2(F)\mathrm{sinh}^2(F)F,_{t}^{}{}_{}{}^{2}$$
$$\frac{3}{4}Kh,_t+\frac{3}{2}K^3h^2\mathrm{sinh}^4(F)h,_t+\frac{3}{4}K^3h^2\mathrm{sinh}^2(F)h,_t$$
$$\frac{3}{2}hK,_th\mathrm{sinh}^2(F)K,_tKh\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)F,_t+\frac{3}{2}Kh\mathrm{cosh}^3(F)\mathrm{sinh}^3(F)F,_t$$
$$\frac{1}{4}K^1\mathrm{cosh}^2(F)\mathrm{sinh}^2(F)\frac{1}{2}Kh^2\mathrm{sinh}^4(F)\frac{1}{4}Kh^2\mathrm{sinh}^2(F)$$
(4.5)
This is used to eliminate $`K,_{tt}`$ from the right-hand side of the equation determinig $`F,_{tt}`$ (see Appendix C), and the result is:
$$F,_{tt}=\frac{3}{4}K^2\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)h,_{t}^{}{}_{}{}^{2}\frac{3}{2}Kh\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)K,_th,_t$$
$$+2K^2\mathrm{cosh}^1(F)\mathrm{sinh}(F)K,_{t}^{}{}_{}{}^{2}K^1K,_tF,_t\mathrm{cosh}(F)\mathrm{sinh}^1(F)F,_{t}^{}{}_{}{}^{2}$$
$$+\frac{3}{4}K^2h^2\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)h,_t+\frac{3}{2}K^2h^2\mathrm{cosh}^1(F)\mathrm{sinh}^3(F)h,_t$$
$$+\frac{1}{2}\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)h,_t\frac{3}{4}\mathrm{cosh}^1(F)\mathrm{sinh}(F)h,_t$$
$$\frac{5}{2}K^1h\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)K,_t\frac{3}{2}K^1h\mathrm{cosh}^1(F)\mathrm{sinh}(F)K,_t$$
$$+h\mathrm{sinh}^2(F)F,_t+\frac{3}{2}hF,_t\frac{1}{4}K^2\mathrm{cosh}^1(F)\mathrm{sinh}(F)\frac{3}{4}K^2\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)$$
$$\frac{1}{2}h^2\mathrm{cosh}^1(F)\mathrm{sinh}^3(F)\frac{1}{4}h^2\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)$$
(4.6)
With (4.2), (4.3), (4.5) and (4.6) all substituted into (A.5), the equation $`G_{11}=\mathrm{\Lambda }`$ reduces to the following form:
$$G_{11}=\frac{1}{4}K^2\mathrm{cosh}^2(F)h,_{t}^{}{}_{}{}^{2}+\frac{3}{2}Kh\mathrm{cosh}^2(F)K,_th,_t+\frac{3}{2}K^2h\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)F,_th,_t$$
$$2K^2\mathrm{cosh}^2(F)\mathrm{sinh}^2(F)K,_{t}^{}{}_{}{}^{2}2K^1\mathrm{cosh}^1(F)\mathrm{sinh}(F)F,_tK,_t+F,_{t}^{}{}_{}{}^{2}$$
$$+\frac{3}{2}K^2h^2\mathrm{cosh}^2(F)h,_t3K^2h^2\mathrm{sinh}^2(F)h,_t\frac{1}{2}\mathrm{cosh}^2(F)h,_t+\frac{3}{2}h,_t$$
$$+\frac{5}{2}K^1h\mathrm{cosh}^2(F)K,_t+3K^1hK,_t\frac{5}{2}h\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)F,_t3h\mathrm{cosh}^1(F)\mathrm{sinh}(F)F,_t$$
$$+\frac{1}{4}K^2\mathrm{cosh}^2(F)+\frac{3}{2}K^2+\frac{1}{2}h^2\mathrm{cosh}^2(F)+h^2\mathrm{sinh}^2(F)=\mathrm{\Lambda }.$$
(4.7)
Now it may be verified that $`G_{11}=`$ const is preserved by eqs. (4.3), (4.5) and (4.6). This is done as follows. The derivative $`\frac{\mathrm{d}}{\mathrm{d}t}G_{11}`$ is calculated, and $`h,_{tt}`$, $`K,_{tt}`$ and $`F_{tt}`$ that reappear are eliminated using (4.3), (4.5) and (4.6). Then, $`K,_{t}^{}{}_{}{}^{2}`$ is found from (4.7) and used to eliminate $`K,_{t}^{}{}_{}{}^{3}`$ and $`K,_{t}^{}{}_{}{}^{2}`$ from $`\frac{\mathrm{d}}{\mathrm{d}t}G_{11}`$. The result is the identity $`\frac{\mathrm{d}}{\mathrm{d}t}G_{11}0`$. This means that, in virtue of the other field equations, if $`G_{11}=\mathrm{\Lambda }`$ holds at any given time, then it will remain constant at all other times. Hence, $`G_{11}=\mathrm{\Lambda }`$ is a limitation imposed by the Einstein equations on the initial data for eqs. {(4.3), (4.5), (4.6)}, and it defines the cosmological constant in terms of the other constants that will appear after (4.3), (4.5) and (4.6) are solved. If $`\mathrm{\Lambda }=0`$, then $`G_{11}=0`$ reduces the number of arbitrary constants by 1.
Hence, with (4.4), we are left with only four equations: (4.2), (4.7) and any two equations from the set $`S=\{(4.3),(4.5),(4.6)\}`$, to determine the four functions $`K_{33}`$, $`h`$, $`K`$ and $`F`$. The third equation in $`S`$ is implied by the remaining two together with (4.7). The only field equation that has not yet been used up is:
$$G_{00}=(8\pi G/c^4)\rho \mathrm{\Lambda }.$$
(4.8)
This may be expected to simply define the matter-density in terms of the metric functions. However, in the formulation used in this paper, matter-density enters the equations in two ways: as a source term in $`G_{00}`$ above, and also through (2.5). From (2.5) and (2.11) it follows that $`\rho `$ must be related to the other functions by:
$$\rho =(K_{11}KK_{33})^1.$$
(4.9)
Together with (4.8) and (4.1) this implies that the following must hold:
$$[(G_{00}+\mathrm{\Lambda })\mathrm{cosh}(F)KK_{33}],_t0.$$
(4.10)
Indeed, this is an identity. This is verified as follows. First, (4.1), (4.2), (4.3), (4.5) and (4.6) are substituted into (A.1) (with $`f=g=0`$) to eliminate all second derivatives. Then, (4.10) is calculated, and (4.2), (4.3), (4.5) and (4.6) are used to eliminate $`K_{33},_t`$ and all second derivatives again. Finally, (4.7) is used to eliminate $`K,_{t}^{}{}_{}{}^{3}`$ and $`K,_{t}^{}{}_{}{}^{2}`$ from the left-hand side of (4.10). In the end, the identity (4.10) results. Hence, (2.5) and (4.8) are consistent in virtue of the other field equations, and moreover $`(G_{00}+\mathrm{\Lambda })\mathrm{cosh}(F)KK_{33}=C=\mathrm{const}`$ (with second derivatives of $`h`$, $`K`$ and $`F`$ eliminated by (4.3), (4.5) and (4.6) and with $`K,_{t}^{}{}_{}{}^{2}`$ eliminated by (4.7)) is the following first integral of the Einstein equations:
$$K_{33}[3Kh\mathrm{sinh}(F)F,_t\frac{3}{2}Kh^2\mathrm{cosh}^1(F)+\frac{3}{2}K\mathrm{cosh}^1(F)\mathrm{sinh}^2(F)h,_t$$
$$+3h\mathrm{cosh}^1(F)\mathrm{sinh}^2(F)K,_t\frac{3}{2}K^1\mathrm{sinh}^2(F)\mathrm{cosh}^1(F)\frac{3}{2}K^3h^2\mathrm{cosh}^1(F)h,_t]=C.$$
(4.11)
Note that, from (4.8) and (4.9), $`C=8\pi G/c^40`$, and so (4.11) determines $`K_{33}`$ algebraically. Hence, (4.11) can replace (4.2) as the definition of $`K_{33}`$. Thereby, the problem of this paper was reduced to the following procedure:
1. Find the most general solution of the set $`\{(4.3),(4.5),(4.6)\}`$. It will contain 6 arbitrary constants $`\{C_1,\mathrm{},C_6\}`$.
2. Impose (4.7) on the $`\{h,K,F\}`$ found in the previous step. This will be just a definition of $`\mathrm{\Lambda }`$ in terms of $`\{C_1,\mathrm{},C_6\}`$ or, when $`\mathrm{\Lambda }=0`$, an additional constraint imposed on $`\{C_1,\mathrm{},C_6\}`$.
3. Calculate $`K_{33}`$ from (4.11), with $`C=8\pi G/c^4`$.
4. Calculate the matter-density from (4.9).
As shown in Ref. 9, an efficient method to find first integrals of a set of equations exists if the set can be obtained from a Lagrangian. Unfortunately, the problem of determinig whether a given set of equations is derivable from a lagrangian is rather complicated and unsolved in general <sup>10</sup>. It is known that the Einstein equations for class B Bianchi metrics may not admit a lagrangian, even though the general Einstein equations do (see Ref. 11 for an explanation). It is shown in Appendix D that eqs. $`\{4.3),(4.5),(4.6)\}`$ do not follow from the most natural lagrangian conceivable in this case: a second-degree polynomial in the first derivatives of $`h`$, $`K`$ and $`F`$, with coefficients being functions of $`h`$, $`K`$ and $`F`$.
For further reference, let us consider the limit of zero rotation in (4.2) – (4.3) and (4.5) – (4.7). After the reparametrization (3.3) we have:
$$\mathrm{cosh}(F)=\stackrel{~}{K_{11}}/\omega _0,\mathrm{sinh}(F)=\sqrt{\stackrel{~}{K}_{11}^{}{}_{}{}^{2}/\omega _{0}^{}{}_{}{}^{2}1},$$
$$F,_t=\stackrel{~}{K_{11,t}}/\sqrt{\stackrel{~}{K}_{11}^{}{}_{}{}^{2}\omega _{0}^{}{}_{}{}^{2}},$$
(4.12)
and then (4.2) in the limit $`\omega _00`$ becomes:
$$\stackrel{~}{K}_{33,t}=\stackrel{~}{K}_{33}(2K,_t/K\stackrel{~}{K_{11,t}}/\stackrel{~}{K_{11}}),$$
(4.13)
which is an identity in the Friedmann limit $`\stackrel{~}{K_{11}}=K=\stackrel{~}{K}_{33}=R(t)`$.
Eq. (4.3) could in fact be discarded in the limit $`\omega _00`$. This is because eq. (4.3) was derived from (A.3), and those terms in (A.3) that lead to (4.3) are all multiplied by $`\omega _{0}^{}{}_{}{}^{2}`$ after the reparametrization (3.3). The off-diagonal component of (3.4) that is proportional to $`h`$ will vanish with any $`h`$ when $`\omega _00`$. Nevertheless, (4.3) gives a result consistent with the other equations in this limit. The limiting form of it is:
$$h,_{tt}=5h,_tK,_t/K+2\stackrel{~}{K_{11,t}}h,_t/\stackrel{~}{K_{11}}+K,_t/K^3.$$
(4.14)
The limit $`\omega _00`$ of (4.5) is:
$$K,_{tt}=K\stackrel{~}{K}_{11,t}^{}{}_{}{}^{2}/\stackrel{~}{K}_{11}^{}{}_{}{}^{2}\frac{3}{4}Kh,_t\frac{3}{2}hK,_t+\frac{3}{2}Kh\stackrel{~}{K_{11,t}}/\stackrel{~}{K_{11}}1/(4K)+2K,_t\stackrel{~}{K_{11,t}}/\stackrel{~}{K_{11}}.$$
(4.15)
The same limit of (4.6) is:
$$\stackrel{~}{K}_{11,tt}/\stackrel{~}{K_{11}}=2K,_{t}^{}{}_{}{}^{2}/K^2K,_t\stackrel{~}{K_{11,t}}/(K\stackrel{~}{K_{11}})+\frac{3}{2}h\stackrel{~}{K_{11,t}}/\stackrel{~}{K_{11}}\frac{3}{2}h,_t\frac{3}{2}hK,_t/K1/(4K^2).$$
(4.16)
In the Friedmann limit $`\stackrel{~}{K_{11}}=K=R(t)`$, eqs. (4.15) and (4.16) become identical:
$$R,_{tt}/R=R,_{t}^{}{}_{}{}^{2}/R^2\frac{3}{4}h,_t1/(4R^2).$$
(4.17)
Finally, the limit $`\omega _00`$ of (4.7) is:
$$2K,_{t}^{}{}_{}{}^{2}/K^22K,_t\stackrel{~}{K_{11,t}}/(K\stackrel{~}{K_{11}})+\stackrel{~}{K}_{11,t}^{}{}_{}{}^{2}/\stackrel{~}{K}_{11}^{}{}_{}{}^{2}+\frac{3}{2}h,_t+3hK,_t/K3h\stackrel{~}{K_{11,t}}/\stackrel{~}{K_{11}}$$
$$+3/(2K^2)=\mathrm{\Lambda }.$$
(4.18)
The Friedmann limit of this is:
$$3R,_{t}^{}{}_{}{}^{2}/R^2+\frac{3}{2}h,_t+3/(2R^2)=\mathrm{\Lambda }.$$
(4.19)
Finding $`h,_t`$ from (4.19) and substituting it in (4.17) we obtain:
$$R,_{tt}/R=R,_{t}^{}{}_{}{}^{2}/(2R^2)+1/(2R^2)\mathrm{\Lambda }/2,$$
(4.20)
which is exactly one of the Friedmann equations. Incidentally, the $`h,_t`$ found from (4.19), if substituted in (4.14), leads to (4.20) again. Hence, in the Friedmann limit, (4.14) follows from (4.19) and (4.17), and need not be discarded.
Note that also (4.11) has a meaningful Friedmann limit. In order to make this limit finite, it must be assumed that:
$$C=\stackrel{~}{C}/\omega _{0}^{}{}_{}{}^{2},$$
(4.21)
and then the limit $`\omega _00`$ of (4.11) is:
$$K_{33}[\mathrm{\Lambda }K\stackrel{~}{K_{11}}+2K,_{t}^{}{}_{}{}^{2}\stackrel{~}{K_{11}}/K+2K,_t\stackrel{~}{K_{11,t}}K\stackrel{~}{K}_{11,t}^{}{}_{}{}^{2}/\stackrel{~}{K_{11}}3\stackrel{~}{K_{11}}/K]=\stackrel{~}{C}.$$
(4.22)
In the Friedmann limit this becomes:
$$R(\mathrm{\Lambda }R^2+3R,_{t}^{}{}_{}{}^{2}3)=\stackrel{~}{C}.$$
(4.23)
Recalling the Friedmann formula for the mass-density, with $`k=1`$:
$$3R,_{t}^{}{}_{}{}^{2}/R^23/R^2+\mathrm{\Lambda }=(8\pi G/c^2)\rho ,$$
(4.24)
we recognize in (4.23) the familiar mass-conservation formula of the Friedmann model, $`\rho R^3=c^2\stackrel{~}{C}/(8\pi G)=\mathrm{const}`$.
## 5 The Lie point-symmetries of the equations (4.3), (4.5) and (4.6).
Point symmetries of (sets of) differential equations are transformations in the space of the independent + dependent variables that leave the set of solutions of the equations unchanged. The point symmetries that form Lie groups (if they exist for a given set of equations) can help in transforming apparently intractable equations into solvable ones by adapting the variables suitably to the generators of the symmetries. The background philosophy and many of the methods are analogous to simplifying the Einstein equations by adapting the coordinates to the Killing vector fields (if such exist). It is assumed that the readers are familiar with this latter procedure. The basic definitions and theorems concerning point symmetries are presented in detail in Refs. 9 and 10.
Eqs. (4.3), (4.5) and (4.6) are of the following form:
$$\frac{\mathrm{d}^2z^i}{\mathrm{d}t^2}=W_{}^{i}{}_{jk}{}^{}\frac{\mathrm{d}z^j}{\mathrm{d}t}\frac{\mathrm{d}z^k}{\mathrm{d}t}+V_{}^{i}{}_{j}{}^{}\frac{\mathrm{d}z^j}{\mathrm{d}t}+U^i,$$
(5.1)
where $`i=0,1,2`$; $`(z^0,z^1,z^2)=(h,K,F)`$ and $`W_{}^{i}{}_{jk}{}^{}`$, $`V_{}^{i}{}_{j}{}^{}`$ and $`U^i`$ are functions of the $`z^i`$, but not of $`t`$. (Incidentally, the independence of $`t`$ of all these coefficients immediately implies one group of symmetries, $`tt^{}=t+s`$, where $`s`$ is the group parameter. This group will emerge from the calculation below.) Let the following be a one-dimensional group of point transformations:
$$t^{}=t^{}(t,\{z^j\},\tau ),z^i=z^i(t,\{z^j\},\tau ),$$
(5.2)
where $`\tau `$ is the group parameter and $`\tau =\tau _0`$ corresponds to the identity (so that $`t^{}(t,\{z^j\},\tau _0)t`$, etc.). The generators of this group (the field of vectors tangent to the orbits of the group (5.2)) are then:
$$X=\xi \frac{}{t}+\eta ^j\frac{}{z^j},$$
(5.3)
where:
$$\left[\begin{array}{c}\xi \\ \eta ^j\end{array}\right]=\frac{\mathrm{d}}{\mathrm{d}\tau }\left[\begin{array}{c}t^{}\\ z^j\end{array}\right]_{\tau =\tau _0}.$$
(5.4)
The generator $`X`$ is extended to arbitrary derivatives $`\frac{\mathrm{d}^kz}{\mathrm{d}t^k}:=\stackrel{(k)}{z}`$ by the recursive formulae:
$$\stackrel{(0)j}{\eta }=\eta ^j,\stackrel{(k)j}{\eta }=\frac{\mathrm{d}\stackrel{(k1)j}{\eta }}{\mathrm{d}t}\frac{\mathrm{d}^kz^j}{\mathrm{d}t^k}\frac{\mathrm{d}\xi }{\mathrm{d}t},$$
(5.5)
and by:
$$\stackrel{(k)}{X}=\xi \frac{}{t}+\eta ^j\frac{}{z^j}+\stackrel{(1)j}{\eta }\frac{}{\stackrel{(1)j}{z}}+\mathrm{}+\stackrel{(k)j}{\eta }\frac{}{\stackrel{(k)j}{z}}.$$
(5.6)
The derivatives $`\frac{\mathrm{d}}{\mathrm{d}t}`$ in (5.5) are total derivatives, i.e.
$$\frac{\mathrm{d}}{\mathrm{d}t}f(t,\{z^i\},\{\stackrel{(1)i}{z}\},\mathrm{},\{\stackrel{(k)i}{z}\})=\frac{f}{t}+\frac{\mathrm{d}z^j}{\mathrm{d}t}\frac{f}{z^j}+\underset{p=1}{\overset{k}{}}\stackrel{(p+1)j}{z}\frac{}{\stackrel{(p)j}{z}},$$
and the order $`n`$ to which the generator $`X`$ has to be extended is equal to the highest order of derivatives in the set (5.1) ($`n=2`$ in our case). A generator of a point symmetry obeys then:
$$\stackrel{(n1)}{X}\mathrm{\Omega }^i=\frac{\mathrm{d}\stackrel{(n1)i}{\eta }}{\mathrm{d}t}\mathrm{\Omega }^i\frac{\mathrm{d}\xi }{\mathrm{d}t},$$
(5.7)
where $`\mathrm{\Omega }^i`$ is the right-hand side of (5.1). (The right-hand side of (5.7) is the $`\stackrel{(n)i}{\eta }`$ as given by (5.5), but with $`\frac{\mathrm{d}^nz^i}{\mathrm{d}t^n}`$ replaced by $`\mathrm{\Omega }^i`$ from (5.1)). Eqs. (5.7) must be identities in all the derivatives $`\stackrel{(1)i}{z},\mathrm{},\stackrel{(n1)i}{z}`$, and so they imply several separate equations to be obeyed by the $`\xi `$ and $`\eta ^i`$.
The procedure in finding and exploiting point symmetries is thus the following:
1. Find the general solution of (5.7) for $`X`$. Since the generators form a Lie algebra (see Ref. 9), the most general $`X`$ will be spanned on a finite number of basis vector fields $`X_{(k)}`$.
2. Read off the basis $`X_{(k)}`$ from that solution.
3. Adapt the variables $`\{t^{}(t,\{z^j\}),z^i(t,\{z^j\})\}`$ to the basis fields $`X_{(k)}`$ so as to maximally simplify the equations.
For our equations (5.1), eqs. (5.7) imply the following four relations:
$$\xi ,_{kl}+W_{}^{j}{}_{kl}{}^{}\xi ,_j=0,$$
(5.8)
$$\eta ^i,_{kl}=W_{}^{i}{}_{kl,s}{}^{}\eta ^s+2W_{}^{i}{}_{s(l}{}^{}\eta ^s,_{k)}W_{}^{s}{}_{kl}{}^{}\eta ^i,_s+\delta _{}^{i}{}_{(l}{}^{}V_{}^{s}{}_{k)}{}^{}\xi ,_s+V^i,_{(l}\xi ,_{k)}+2\frac{^2\xi }{tz^{(k}}\delta _{}^{i}{}_{l)}{}^{},$$
(5.9)
where parentheses on indices denote symmetrization,
$$\frac{^2\eta ^i}{tz^k}=W_{}^{i}{}_{ks}{}^{}\frac{\eta ^s}{t}+\frac{1}{2}V_{}^{i}{}_{k,s}{}^{}\eta ^s+\frac{1}{2}V_{}^{i}{}_{s}{}^{}\eta ^s,_k\frac{1}{2}V_{}^{s}{}_{k}{}^{}\eta ^i,_s$$
$$+\frac{1}{2}V_{}^{i}{}_{k}{}^{}\frac{\xi }{t}+U^i\xi ,_k+\frac{1}{2}\delta _{}^{i}{}_{k}{}^{}U^s\xi ,_s+\frac{1}{2}\delta _{}^{i}{}_{k}{}^{}\frac{^2\xi }{t^2},$$
(5.10)
$$\frac{^2\eta ^i}{t^2}=V_{}^{i}{}_{s}{}^{}\frac{\eta ^s}{t}+U^i,_s\eta ^sU^s\eta ^i,_s+2U^i\frac{\xi }{t}.$$
(5.11)
The general solution of these equations (with $`W_{}^{i}{}_{kl}{}^{}`$, $`V_{}^{i}{}_{k}{}^{}`$ and $`U^i`$ read off from (4.3), (4.5) and (4.6)) is:
$$X=A\frac{}{t}+B(t\frac{}{t}h\frac{}{h}+K\frac{}{K}),$$
(5.12)
where $`A`$ and $`B`$ are arbitrary constants. The proof that this is the most general solution is laborious but straightforward, it is given in Appendix E. Hence, our set of equations has a two-dimensional symmetry group whose generators are:
$$X_{(1)}=\frac{}{t},X_{(2)}=t\frac{}{t}h\frac{}{h}+K\frac{}{K},$$
(5.13)
and the corresponding finite symmetry transformations are:
$$t^{}=t+\tau _1,(h^{},K^{},F^{})=(h,K,F);$$
$$t^{}=\mathrm{e}^{\tau _2}t,h^{}=\mathrm{e}^{\tau _2}h,K^{}=\mathrm{e}^{\tau _2}K,F^{}=F,$$
(5.14)
where $`\tau _1`$ and $`\tau _2`$ are the group parameters. The first symmetry was self-evident, as already mentioned, and the second one can be verified by inspection of the equations (4.3), (4.5) and (4.6).
Unfortunately, these symmetries do not lead to any discernible simplification of the set $`S=\{(4.3),(4.5),(4.6)\}`$. In variables adapted to the generator $`X_{(1)}`$, the independent variable is $`K`$, and $`t(K)`$ is one of the functions. The set (5.1) thus transformed is of first order in $`\varphi (K):=dt/dK`$, but the first-order equation is still a member of a complicated set and none of the equations separates out. Moreover, after the transformed set is algebraically solved for $`t,_{KK}`$, $`h,_{KK}`$ and $`F,_{KK}`$, the right-hand sides become polynomials of third degree in $`t,_k`$, $`h,_K`$ and $`F,_K`$.
The variables adapted to the generator $`X_{(2)}`$ are $`(t^{},h^{},K^{})`$, where:
$$t=\mathrm{e}^K^{}t^{},K=\mathrm{e}^K^{},h=\mathrm{e}^K^{}h^{}.$$
(5.15)
In these variables, the set (5.1) becomes of first order in $`\psi (t^{})=K^{},_t^{}`$. However, after it is solved for $`h^{},_{t^{}t^{}}`$, $`K^{},_{t^{}t^{}}`$ and $`F,_{t^{}t^{}}`$, the right-hand sides of $`h^{},_{t^{}t^{}}`$ and $`F,_{t^{}t^{}}`$ contain rational functions of the form $`W/(1+t^{}K^{},_t^{})`$, where $`W`$ is a monomial of second degree in some of the $`h^{},_t^{}`$, $`K^{},_t^{}`$ and $`F,_t^{}`$. Neither equation separates out. It is not possible to adapt the variables to both the generators simultaneously because the group is nonabelian. This author was not able to make any use of the new variables.
## 6 First integrals that are polynomials in $`(h,_t,K,_t,F,_t)`$.
Suppose that the set $`\stackrel{~}{S}=\{(4.3),(4.5),(4.6),(4.7)\}`$ has a first integral of the form:
$$I:=Q_{ij}\dot{z}^i\dot{z}^j+L_i\dot{z}^i+E=C=\mathrm{const},$$
(6.1)
where $`C`$ is an arbitrary constant, $`Q_{ij}=Q_{ji}`$, $`L_i`$ and $`E`$ are unknown functions of $`(h,K,F)`$, $`i,j=1,2,3`$, $`z^1=h,z^2=K,z^3=F`$. Then $`\frac{\mathrm{d}I}{\mathrm{d}t}0`$ in virtue of $`\stackrel{~}{S}`$, i.e. using (5.1) to eliminate $`\ddot{z}^i`$:
$$(2Q_{ij}\dot{z}^j+L_i)(W_{}^{i}{}_{kl}{}^{}\dot{z}^k\dot{z}^l+V_{}^{i}{}_{k}{}^{}\dot{z}^k+U^i)+Q_{ij,k}\dot{z}^i\dot{z}^j\dot{z}^k+L_{i,j}\dot{z}^i\dot{z}^j+E,_i\dot{z}^i=0.$$
(6.2)
In showing that (6.2) is zero, (4.7) must be used. Eq. (4.7) may be safely used to eliminate $`F,_{t}^{}{}_{}{}^{3}`$ and $`F,_{t}^{}{}_{}{}^{2}`$, but not the remaining $`F,_t`$. This is because $`F,_t`$ found from (4.7) would be of the form:
$$F,_t=P(h,_t,K,_t)+\sqrt{\mathrm{\Delta }(h,_t,K,_t)},$$
(6.3)
where $`P`$ and $`\mathrm{\Delta }`$ are polynomials of degree 2 in $`h,_t`$ and $`K,_t`$. If $`\mathrm{\Delta }`$ were a square of a first-degree polynomial, then (6.3) could be used to eliminate $`F,_t`$ altogether from (6.2). However, $`\mathrm{\Delta }`$ being a square implies an additional equation obeyed by $`h,_t`$ and $`K,_t`$ (the discriminant of $`\mathrm{\Delta }`$ must be zero). Hence, if $`h,_t`$ and $`K,_t`$ are to be treated as independent, then $`F,_t`$ is linearly independent of $`h,_t`$ and $`K,_t`$. Then the coefficients of $`F,_t`$ in (6.2) must sum up to zero anyway, and eliminating $`F,_t`$ is of no use.
Knowing this, it can be verified that first integrals of the form (6.1) do not exist. The calculations are conceptually straightforward, but lead through horrible intermediate expressions, so they are not reported here. The hypothesis that (6.1) is a first integral uniquely leads to an equation that is equivalent to (4.7).
The same method may be used to test whether our set of equations admits a constraint that would be a polynomial of degree 1 or 2 in the first derivatives. The only difference with respect to the procedure of looking for a first integral is that in verifying whether (6.2) is zero, eq. (6.1) is used, too. If a nontrivial solution of (6.2) with this additional simplification is found, then it means that the derivative of (6.1) by $`t`$ is zero if (6.1) holds for any fixed $`t`$. Then, such (6.1) is a constraint preserved by the set $`S`$. However, even this attempt has not led to useful results. Constraints of degree 2, i.e. those with $`Q_{ij}0`$, lead to prohibitively complicated equations and could not be investigated. One constraint of the form (6.1) with $`Q_{ij}=0`$ was found, but it is equivalent to the square bracket in (4.11) being zero, and so implies zero matter density. Again, the details are not reported because they contain complicated equations, but no ingenious ideas. This result proves the usefulness of the method – a sensible constraint was revealed – but the solution with zero density is not interesting for cosmology, and thus not necessarily worth investigating.
The zero-density constraint was found without using eq. (4.7). Eq. (4.7) would reduce the number of unknown functions by one, but the resulting set of equations is prohibitively complicated and no progress was achieved.
## 7 Summary of results.
It was shown that the Einstein equations for the metric (2.9) with $`f=g=0`$ are self-consistent and solvable. They reduce to the set $`S=\{(4.3),(4.5),(4.6)\}`$ to determine $`h`$, $`K`$ and $`K_{11}=\mathrm{cosh}(F)`$, and (4.11) to determine $`K_{33}`$ (where $`C=8\pi G/c^4`$). The matter density is found from (4.9). The first derivatives of the functions obeying the set $`S`$ must obey (4.7).
The Friedmann solution with $`k=1`$ is contained among the solutions of this set, as shown in eqs. (4.12) – (4.24). Unfortunately, no explicit example of a more general solution could be found. Attempts to follow ad hoc Ansatzes produced uninteresting results. The Ansatz $`K=K_{33}`$ led to the deSitter solution in disguise, in which the t-lines had nonzero rotation. The Ansatz $`K_{11}=K/C`$ ($`C=`$ const), which is consistent with the Friedmann limit, led to such complicated equations that it could not even be verified if they are not contradictory. The assumption of zero shear implies zero expansion, in virtue of the theorem $`(\sigma =0)(\omega \theta =0)`$ that holds for dust (see Ref. 12).
The set $`S`$ was shown to have a two-dimensional group of point symmetries, given by (5.14), and to admit no Lagrangian of the Hilbert type. It was also verified that no first integrals of the form (6.1) exist.
The progress achieved in this paper was the reduction of the problem of existence of a rotating generalization of the $`k=1`$ Friedmann model to the technical problem of finding an explicit solution of the set $`S`$. The solvability of the set $`S`$ may be taken for granted because the Friedmann model itself was shown to be one of its solutions. It is still unknown, though, whether a continuous family of solutions exists labeled by by the parameter $`\omega `$ (rotation) such that the limit $`\omega 0`$ taken in the explicit solution leads to the $`k=1`$ Friedmann model.
A similar analysis as done here should be done for the other promising cases identified in Ref. 3.
Acknowledgements. The algebraic manipulations for this paper were carried out using the computer algebra system Ortocartan<sup>13,14</sup> that was for this purpose extended by several new programs<sup>15</sup>. The author is grateful to J. Kijowski, J. Jezierski, K. Rosquist, R. Jantzen, M. A. H. MacCallum and C. Uggla for useful comments and instructions on the Lagrangian methods.
Appendix A
The Einstein equations for the metric (2.9).
As explained in sec. 2 (after eq. (2.9)), these are the projections of the Einstein tensor on the forms of the orthonormal tetrad (2.10), thus for example $`G_{03}`$ below is equal to $`e_{I}^{}{}_{}{}^{\alpha }e_{J}^{}{}_{}{}^{\beta }G_{\alpha \beta }`$ with I = 0 and J = 3, where $`G_{\alpha \beta }`$ are the coordinate components of the Einstein tensor.
$$G_{00}=2K_{11}^{}{}_{}{}^{3}hK_{11},_t+K_{11}^{}{}_{}{}^{3}K^1K_{11},_tK,_t+K_{11}^{}{}_{}{}^{3}K_{33}^{}{}_{}{}^{1}K_{11},_tK_{33},_t$$
$$+\frac{3}{4}K_{11}^{}{}_{}{}^{2}K^2\frac{1}{4}K_{11}^{}{}_{}{}^{2}K^2K_{33}^{}{}_{}{}^{2}f,_{t}^{}{}_{}{}^{2}3K_{11}^{}{}_{}{}^{2}K^1hK,_t$$
$$K_{11}^{}{}_{}{}^{2}K^1K_{33}^{}{}_{}{}^{1}K,_tK_{33},_tK_{11}^{}{}_{}{}^{2}K^1K,_{tt}\frac{1}{4}K_{11}^{}{}_{}{}^{2}K^2h,_{t}^{}{}_{}{}^{2}$$
$$3K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{1}hK_{33},_tK_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{1}K_{33},_{tt}+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}hf,_tg,_t$$
$$\frac{1}{4}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}h^2f,_{t}^{}{}_{}{}^{2}\frac{1}{4}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}g,_{t}^{}{}_{}{}^{2}3K_{11}^{}{}_{}{}^{2}h^2\frac{5}{2}K_{11}^{}{}_{}{}^{2}h,_t$$
$$+K_{11}^{}{}_{}{}^{1}K^1K_{11},_tK,_t+K_{11}^{}{}_{}{}^{1}K_{33}^{}{}_{}{}^{1}K_{11},_tK_{33},_t\frac{1}{4}K^2K_{33}^{}{}_{}{}^{2}f,_{t}^{}{}_{}{}^{2}$$
$$+K^1K_{33}^{}{}_{}{}^{1}K,_tK_{33},_t3K^2$$
$`(A.1)`$
$$G_{01}=2K_{11}^{}{}_{}{}^{2}hK_{11},_tK_{11}^{}{}_{}{}^{2}K^1K_{11},_tK,_tK_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{1}K_{11},_tK_{33},_t$$
$$+\frac{1}{2}K_{11}^{}{}_{}{}^{1}K^2+\frac{1}{2}K_{11}^{}{}_{}{}^{1}K^2K_{33}^{}{}_{}{}^{2}f,_{t}^{}{}_{}{}^{2}+K_{11}^{}{}_{}{}^{1}K^1hK,_t$$
$$+K_{11}^{}{}_{}{}^{1}K^1K,_{tt}+K_{11}^{}{}_{}{}^{1}K_{33}^{}{}_{}{}^{1}hK_{33},_t+K_{11}^{}{}_{}{}^{1}K_{33}^{}{}_{}{}^{1}K_{33},_{tt}+\frac{3}{2}K_{11}^{}{}_{}{}^{1}h,_t$$
$`(A.2)`$
$$G_{02}=\frac{1}{2}K_{11}^{}{}_{}{}^{3}KK_{11},_th,_t\frac{1}{2}K_{11}^{}{}_{}{}^{3}K^1K_{11},_t\frac{3}{2}K_{11}^{}{}_{}{}^{2}Khh,_t$$
$$\frac{1}{2}K_{11}^{}{}_{}{}^{2}KK_{33}^{}{}_{}{}^{1}K_{33},_th,_t\frac{1}{2}K_{11}^{}{}_{}{}^{2}Kh,_{tt}\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^2K,_t$$
$$+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^1h+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^1K_{33}^{}{}_{}{}^{1}K_{33},_t+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^1K_{33}^{}{}_{}{}^{2}hf,_{t}^{}{}_{}{}^{2}$$
$$\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^1K_{33}^{}{}_{}{}^{2}f,_tg,_t\frac{3}{2}K_{11}^{}{}_{}{}^{2}K,_th,_tK_{11}^{}{}_{}{}^{1}K^1K_{11},_t$$
$$+2K^2K,_tK^1K_{33}^{}{}_{}{}^{1}K_{33},_t$$
$`(A.3)`$
$$G_{03}=\frac{1}{2}K_{11}^{}{}_{}{}^{3}K_{33}hK_{11},_tf,_t+\frac{1}{2}K_{11}^{}{}_{}{}^{3}K_{33}K_{11},_tg,_t+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K_{33}hf,_{tt}$$
$$\frac{3}{2}K_{11}^{}{}_{}{}^{2}K_{33}hg,_t+\frac{3}{2}K_{11}^{}{}_{}{}^{2}K_{33}h^2f,_t\frac{1}{2}K_{11}^{}{}_{}{}^{2}K_{33}g,_{tt}+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K_{33}f,_th,_t$$
$$+\frac{3}{2}K_{11}^{}{}_{}{}^{2}hK_{33},_tf,_t\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^2K_{33}f,_t+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^1K_{33}hK,_tf,_t$$
$$\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^1K_{33}K,_tg,_t\frac{3}{2}K_{11}^{}{}_{}{}^{2}K_{33},_tg,_t+\frac{3}{2}K^2K_{33}f,_t$$
$`(A.4)`$
$$G_{11}=\frac{1}{4}K_{11}^{}{}_{}{}^{2}K^2\frac{1}{4}K_{11}^{}{}_{}{}^{2}K^2K_{33}^{}{}_{}{}^{2}f,_{t}^{}{}_{}{}^{2}+K_{11}^{}{}_{}{}^{2}K^1hK,_t$$
$$+K_{11}^{}{}_{}{}^{2}K^1K_{33}^{}{}_{}{}^{1}K,_tK_{33},_t+\frac{1}{4}K_{11}^{}{}_{}{}^{2}K^2h,_{t}^{}{}_{}{}^{2}+K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{1}hK_{33},_t$$
$$\frac{1}{2}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}hf,_tg,_t+\frac{1}{4}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}h^2f,_{t}^{}{}_{}{}^{2}+\frac{1}{4}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}g,_{t}^{}{}_{}{}^{2}$$
$$+K_{11}^{}{}_{}{}^{2}h^2\frac{1}{2}K_{11}^{}{}_{}{}^{2}h,_t\frac{1}{4}K^2K_{33}^{}{}_{}{}^{2}f,_{t}^{}{}_{}{}^{2}K^1K_{33}^{}{}_{}{}^{1}K,_tK_{33},_t$$
$$K^1K,_{tt}K_{33}^{}{}_{}{}^{1}K_{33},_{tt}+K^2$$
$`(A.5)`$
$$G_{12}=\frac{1}{2}K_{11}^{}{}_{}{}^{2}KK_{11},_th,_t+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^1K_{11},_t+\frac{1}{2}K_{11}^{}{}_{}{}^{1}KK_{33}^{}{}_{}{}^{1}K_{33},_th,_t$$
$$+\frac{1}{2}K_{11}^{}{}_{}{}^{1}Kh,_{tt}\frac{3}{2}K_{11}^{}{}_{}{}^{1}K^2K,_t+\frac{1}{2}K_{11}^{}{}_{}{}^{1}K^1K_{33}^{}{}_{}{}^{1}K_{33},_t$$
$$\frac{1}{2}K_{11}^{}{}_{}{}^{1}K^1K_{33}^{}{}_{}{}^{2}hf,_{t}^{}{}_{}{}^{2}+\frac{1}{2}K_{11}^{}{}_{}{}^{1}K^1K_{33}^{}{}_{}{}^{2}f,_tg,_t+\frac{3}{2}K_{11}^{}{}_{}{}^{1}K,_th,_t$$
$`(A.6)`$
$$G_{13}=\frac{1}{2}K_{11}^{}{}_{}{}^{2}K_{33}hK_{11},_tf,_t\frac{1}{2}K_{11}^{}{}_{}{}^{2}K_{33}K_{11},_tg,_t\frac{1}{2}K_{11}^{}{}_{}{}^{1}K_{33}hf,_{tt}$$
$$+\frac{1}{2}K_{11}^{}{}_{}{}^{1}K_{33}g,_{tt}\frac{1}{2}K_{11}^{}{}_{}{}^{1}K_{33}f,_th,_t\frac{3}{2}K_{11}^{}{}_{}{}^{1}hK_{33},_tf,_tK_{11}^{}{}_{}{}^{1}K^2K_{33}f,_t$$
$$\frac{1}{2}K_{11}^{}{}_{}{}^{1}K^1K_{33}hK,_tf,_t+\frac{1}{2}K_{11}^{}{}_{}{}^{1}K^1K_{33}K,_tg,_t+\frac{3}{2}K_{11}^{}{}_{}{}^{1}K_{33},_tg,_t$$
$`(A.7)`$
$$G_{22}=K_{11}^{}{}_{}{}^{3}hK_{11},_tK_{11}^{}{}_{}{}^{3}K_{33}^{}{}_{}{}^{1}K_{11},_tK_{33},_t+\frac{1}{4}K_{11}^{}{}_{}{}^{2}K^2\frac{1}{4}K_{11}^{}{}_{}{}^{2}K^2K_{33}^{}{}_{}{}^{2}f,_{t}^{}{}_{}{}^{2}$$
$$\frac{3}{4}K_{11}^{}{}_{}{}^{2}K^2h,_{t}^{}{}_{}{}^{2}+2K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{1}hK_{33},_t+K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{1}K_{33},_{tt}+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}hf,_tg,_t$$
$$\frac{1}{4}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}h^2f,_{t}^{}{}_{}{}^{2}\frac{1}{4}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}g,_{t}^{}{}_{}{}^{2}+K_{11}^{}{}_{}{}^{2}h^2+\frac{3}{2}K_{11}^{}{}_{}{}^{2}h,_t$$
$$K_{11}^{}{}_{}{}^{1}K_{33}^{}{}_{}{}^{1}K_{11},_tK_{33},_tK_{11}^{}{}_{}{}^{1}K_{11},_{tt}+\frac{1}{4}K^2K_{33}^{}{}_{}{}^{2}f,_{t}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{1}K_{33},_{tt}+K^2$$
$`(A.8)`$
$$G_{23}=\frac{1}{2}K_{11}^{}{}_{}{}^{3}K^1K_{33}K_{11},_tf,_t+\frac{1}{2}K_{11}^{}{}_{}{}^{2}KK_{33}hf,_th,_t\frac{1}{2}K_{11}^{}{}_{}{}^{2}KK_{33}g,_th,_t$$
$$+\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^2K_{33}K,_tf,_tK_{11}^{}{}_{}{}^{2}K^1K_{33}hf,_t\frac{1}{2}K_{11}^{}{}_{}{}^{2}K^1K_{33}f,_{tt}$$
$$\frac{3}{2}K_{11}^{}{}_{}{}^{2}K^1K_{33},_tf,_t+\frac{1}{2}K_{11}^{}{}_{}{}^{1}K^1K_{33}K_{11},_tf,_t\frac{1}{2}K^2K_{33}K,_tf,_t$$
$$+\frac{1}{2}K^1K_{33}f,_{tt}+\frac{3}{2}K^1K_{33},_tf,_t$$
$`(A.9)`$
$$G_{33}=K_{11}^{}{}_{}{}^{3}hK_{11},_tK_{11}^{}{}_{}{}^{3}K^1K_{11},_tK,_t\frac{1}{4}K_{11}^{}{}_{}{}^{2}K^2+\frac{3}{4}K_{11}^{}{}_{}{}^{2}K^2K_{33}^{}{}_{}{}^{2}f,_{t}^{}{}_{}{}^{2}$$
$$+2K_{11}^{}{}_{}{}^{2}K^1hK,_t+K_{11}^{}{}_{}{}^{2}K^1K,_{tt}\frac{1}{4}K_{11}^{}{}_{}{}^{2}K^2h,_{t}^{}{}_{}{}^{2}+\frac{3}{2}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}hf,_tg,_t$$
$$\frac{3}{4}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}h^2f,_{t}^{}{}_{}{}^{2}\frac{3}{4}K_{11}^{}{}_{}{}^{2}K_{33}^{}{}_{}{}^{2}g,_{t}^{}{}_{}{}^{2}+K_{11}^{}{}_{}{}^{2}h^2+\frac{3}{2}K_{11}^{}{}_{}{}^{2}h,_tK_{11}^{}{}_{}{}^{1}K^1K_{11},_tK,_t$$
$$K_{11}^{}{}_{}{}^{1}K_{11},_{tt}\frac{3}{4}K^2K_{33}^{}{}_{}{}^{2}f,_{t}^{}{}_{}{}^{2}K^1K,_{tt}+K^2$$
$`(A.10)`$
Since the source in the Einstein equations is dust with a cosmological constant, and since the zero-th tetrad vector is just the velocity vector, the above components should obey the following equations:
$$G_{00}=(8\pi G/c^4)\rho \mathrm{\Lambda },$$
$$G_{11}=G_{22}=G_{33}=\mathrm{\Lambda },$$
$$\mathrm{nondiagonal}G_{IJ}=0,$$
$`(A.11)`$
where $`\rho `$ is the dust energy-density and $`\mathrm{\Lambda }`$ is the cosmological constant.
Appendix B
Consequences of $`h=0`$ in the Einstein equations.
With $`h=0`$, eq. (2.12) becomes:
$$K_{33}f,_t(K_{11}^{}{}_{}{}^{2}1)/(K_{11}K^2)=0$$
$`(B.1)`$
We can immediately discard the solution $`K_{33}=0`$ because then $`det(g_{\alpha \beta })=0`$. When $`K_{11}^{}{}_{}{}^{2}=1`$, the limit $`\omega 0`$ of the resulting metric will necessarily have either nonzero shear or zero expansion (see sec. 3 for a calculation of this limit), and so no generalization of the Friedmann models can be expected here. Hence, the only consequence of (B.1) that is worth pursuing is:
$$f,_t=0.$$
$`(B.2)`$
Then, with $`h=0`$, eq. (A.6) implies:
$$K_{11}K_{33}/K^3=\mathrm{const}.$$
$`(B.3)`$
However, in the limit $`\omega 0`$ this again implies either nonzero shear or zero expansion, i.e. no Friedmann limit.
Appendix C
The result for $`F,_{tt}`$ from $`G_{01}=0`$.
When (4.2), (4.2) and (4.3) are substituted in (A.2), the following formula results for $`F,_{tt}`$:
$$F,_{tt}=\frac{3}{2}Kh\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)K,_th,_t+3h\mathrm{cosh}^2(F)F,_t\frac{7}{2}h\mathrm{sinh}^2(F)F,_t3hF,_t$$
$$+\frac{1}{2}K^2\mathrm{cosh}^1(F)\mathrm{sinh}(F)+2K^2\mathrm{cosh}^1(F)\mathrm{sinh}(F)K,_{t}^{}{}_{}{}^{2}$$
$$+\frac{1}{2}K^1h\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)K,_t+3K^1h\mathrm{cosh}^1(F)\mathrm{sinh}(F)K,_t$$
$$+3K^1\mathrm{cosh}^2(F)K,_tF,_t+3K^1\mathrm{cosh}^1(F)\mathrm{sinh}(F)K,_{tt}7K^1K,_tF,_t$$
$$+\frac{9}{2}K^2h\mathrm{sinh}^2(F)F,_th,_t3K^2h^2\mathrm{cosh}^1(F)\mathrm{sinh}^3(F)h,_t$$
$$\frac{3}{2}K^2h^2\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)h,_t\frac{3}{2}K^2\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)h,_{t}^{}{}_{}{}^{2}$$
$$+h^2\mathrm{cosh}^1(F)\mathrm{sinh}^3(F)+\frac{1}{2}h^2\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)+2\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)F,_{t}^{}{}_{}{}^{2}$$
$$+\frac{1}{2}\mathrm{cosh}^1(F)\mathrm{sinh}^1(F)h,_t+2\mathrm{cosh}^1(F)\mathrm{sinh}(F)F,_{t}^{}{}_{}{}^{2}+\frac{3}{2}\mathrm{cosh}^1(F)\mathrm{sinh}(F)h,_t.$$
It will be modified later because it contains $`K,_{tt}`$ on the right-hand side, while the final equations that will be dealt with should have no second derivatives on the right-hand sides.
Appendix D
Nonexistence of a Hilbert-type Lagrangian for the set {(4.3), (4.5), (4.6)}.
Eqs. (4.3), (4.5) and (4.6) can be written in the form:
$$\frac{\mathrm{d}^2z^i}{\mathrm{d}t^2}=W_{}^{i}{}_{jk}{}^{}\frac{\mathrm{d}z^j}{\mathrm{d}t}\frac{\mathrm{d}z^k}{\mathrm{d}t}+V_{}^{i}{}_{j}{}^{}\frac{\mathrm{d}z^j}{\mathrm{d}t}+U^i,$$
$`(D.1)`$
where $`i=0,1,2`$; $`z^0=h`$, $`z^1=K`$, $`z^2=F`$ and $`W_{}^{i}{}_{jk}{}^{}`$, $`V_{}^{i}{}_{j}{}^{}`$ and $`U^i`$ are functions of $`(h,K,F)`$ (but not of $`t`$). Note that the set (D.1) is covariant with respect to arbitrary transformations $`z^iz^i=f^i(\{z^j\})`$: the first derivatives $`\frac{\mathrm{d}z^j}{\mathrm{d}t}`$ transform then like a contravariant vector, and so do the terms $`U^i`$, the coefficients $`V_{}^{i}{}_{j}{}^{}`$ transform like a mixed tensor, and the coefficients $`(W_{}^{i}{}_{jk}{}^{})`$ transform like components of an affine connection. (The nontensorial terms in the transformed $`(W_{}^{i}{}_{jk}{}^{})`$ arise from $`\frac{\mathrm{d}^2z^i}{\mathrm{d}t^2}`$). The most natural ansatz for a lagrangian for (D.1) is:
$$L=Q_{ij}\frac{\mathrm{d}z^i}{\mathrm{d}t}\frac{\mathrm{d}z^j}{\mathrm{d}t}+L_i\frac{\mathrm{d}z^i}{\mathrm{d}t}+\mathrm{\Phi },$$
$`(D.2)`$
where $`Q_{ij}`$, $`L_i`$ and $`\mathrm{\Phi }`$ are functions of $`(h,K,F)`$. Such a lagrangian would result from the Hilbert lagrangian by taking out a complete divergence and integrating the result with respect to the spatial variables. The Euler-Lagrange equations implied by (D.2) are:
$$Q_{is}\frac{\mathrm{d}^2z^s}{\mathrm{d}t^2}=(Q_{ki,l}\frac{1}{2}Q_{kl,i})\frac{\mathrm{d}z^k}{\mathrm{d}t}\frac{\mathrm{d}z^l}{\mathrm{d}t}+\frac{1}{2}(L_{k,i}L_{i,k})\frac{\mathrm{d}z^k}{\mathrm{d}t}+\frac{1}{2}\mathrm{\Phi },_i$$
$`(D.3)`$
If these are to be equivalent to (D.1), then the following must hold:
$$Q_{is}W_{}^{s}{}_{kl}{}^{}=\frac{1}{2}(Q_{ki,l}+Q_{li,k}Q_{kl,i}),$$
$`(D.4)`$
$$Q_{is}V_{}^{s}{}_{k}{}^{}=\frac{1}{2}(L_{k,i}L_{i,k}),$$
$`(D.5)`$
$$Q_{is}U^s=\frac{1}{2}\mathrm{\Phi },_i.$$
$`(D.6)`$
Eqs. (D.4) imply that $`(W_{}^{i}{}_{jk}{}^{})`$ must be Christoffel symbols constructed from the metric $`Q_{ij}`$, eqs. (D.5) imply that $`\frac{1}{2}L_i`$ must be a vector potential for the tensor field $`Q_{is}V_{}^{s}{}_{k}{}^{}`$, and eqs. (D.6) imply that $`\mathrm{\Phi }/2`$ must be a scalar potential for the vector field $`Q_{is}U^s`$. All of these are strong conditions and they may be impossible to fulfil in many cases.
Indeed, for our eqs. {(4.3), (4.5), (4.6)}, the solution of (D.4) turns out to be $`Q_{ij}0`$, i.e. the Lagrangian (D.2) does not exist. This is an outline of the proof.
After eqs. (D.4) are written out in the form
$$Q_{ij,k}=W_{}^{s}{}_{ik}{}^{}Q_{sj}W_{}^{s}{}_{jk}{}^{}Q_{is},$$
$`(D.7)`$
with $`W_{}^{s}{}_{kl}{}^{}`$ read off from {(4.3), (4.5), (4.6)}, the following two equations follow, among other results:
$$Q_{11,F}+\frac{1}{4}(K\mathrm{cosh}(F)/\mathrm{sinh}(F))Q_{11,K}=(2\mathrm{cosh}^2(F)1)Q_{11}/(\mathrm{cosh}(F)\mathrm{sinh}(F)),$$
$`(D.8)`$
$$Q_{22,K}+(2\mathrm{cosh}^2(F)1)\mathrm{sinh}(F)Q_{22,F}/(2K\mathrm{cosh}^3(F))=(3\mathrm{cosh}^2(F)1)Q_{22}/(K\mathrm{cosh}^2(F)).$$
$`(D.9)`$
The solutions of these are:
$$Q_{11}=q_{11}(h,\frac{K^4}{\mathrm{sinh}(F)})\frac{1}{\mathrm{cosh}(F)\mathrm{sinh}(F)},$$
$`(D.10)`$
$$Q_{22}=K\mathrm{sinh}^2(F)q_{22}(h,\frac{K\sqrt{2\mathrm{cosh}^2(F)1}}{\mathrm{sinh}^2(F)}),$$
$`(D.11)`$
where $`q_{ij}`$ are arbitrary functions of their two arguments. The equation $`Q_{11,K}=\mathrm{}`$ is then solved with the result:
$$Q_{12}=K^5q_{11,w}/\mathrm{sinh}^3(F),$$
$`(D.12)`$
where $`w=K^4/\mathrm{sinh}(F)`$ is the second argument of $`q_{11}`$, and the equation $`Q_{22,K}=\mathrm{}`$ implies:
$$\frac{1}{v}q_{22,v}=q_{11,w}K^3/(\mathrm{cosh}(F)\mathrm{sinh}^2(F)),$$
$`(D.13)`$
where $`v`$ is the second argument of $`q_{22}`$. The left-hand side of (D.13) is an invariant of the operator $`(2K\mathrm{cosh}^3(F)/\mathrm{sinh}(F))\frac{}{K}+(2\mathrm{cosh}^2(F)1)\frac{}{F}`$, while $`q_{11,w}`$ is an invariant of the operator $`\frac{1}{4}K\frac{}{K}+(\mathrm{sinh}(F)/\mathrm{cosh}(F))\frac{}{F}`$. Application of these two operators to (D.13) leads to $`q_{22,v}=q_{11,w}=0`$, which implies $`Q_{12}=0`$. With this, the remaining equations (D.7) quickly lead to $`Q_{ij}0`$, which means that the Lagrangian (D.2) does not exist in this case.
Since the Euler-Lagrange equations (D.4) are covariant with respect to arbitrary transformations of the Lagrangian variables (in our case $`hh^{}(h,K,F)`$, etc.), and equations of the form (D.1) are covariant, too, the conclusion that a Lagrangian of the form (D.2) exists (or does not exist) is coordinate-independent, i.e. having shown that eqs. $`\{(4.3),(4.5),(4.6)\}`$ do not follow from a Lagrangian (D.2) in our variables $`\{h,K,F\}`$, we know that no such Lagrangian will exist in any other variables.
Appendix E
The general solution of eqs. (5.8) – (5.11).
Eqs. (5.8) have the form:
$$\xi _{;kl}=0,$$
$`(E.1)`$
where ; is the covariant derivative in which $`(W_{}^{i}{}_{kl}{}^{})`$ play the role of the connection coefficients. (They appear in this role for a second time already, see Appendix D.) The integrability conditions of (E.1) are:
$$R_{}^{s}{}_{ijk}{}^{}\xi ,_s=0,$$
$`(E.2)`$
where $`R_{}^{s}{}_{ijk}{}^{}=R_{}^{s}{}_{ikj}{}^{}`$ is the curvature tensor corresponding to the connection$`(W_{}^{i}{}_{kl}{}^{})`$. Eqs. (E.2) are 9 equations (labelled by the sets of indices $`(i,j,k)=(0,0,1);(0,0,2);(0,1,2)`$; etc) and they could have nontrivial solutions only if every subset of 3 equations chosen from among them had a zero determinant. Actually, of the 84 determinants only two vanish, and some of them will not vanish even if the functions $`h(t)`$, $`K(t)`$ and $`F(t)`$ are functionally dependent. Here is one determinant that will never vanish, it corresponds to $`\{(i,j,k)\}=\{(1,0,1);(1,1,2);(2,1,2)\}`$:
$$\mathrm{det}(E.2)=K^4[(\frac{189}{32}+\frac{3}{4}\mathrm{cosh}^6(F)+\frac{53}{16}\mathrm{cosh}^4(F)+\frac{5}{8}\mathrm{cosh}^2(F)\frac{11}{8}\mathrm{sinh}^2(F))].$$
Hence, the unique solution of (5.8) is:
$$\xi =\xi (t).$$
$`(E.3)`$
With $`\xi ,_i=0`$, eqs. (5.9) simplify somewhat, and the equation corresponding to $`(i,k,l)=(0,1,2)`$ becomes $`\eta ^0,_{KF}=2\eta ^0,_F/K`$, which has the solution:
$$\eta ^0=F^0(t,h,F)/K^2+G^0(t,h,K),$$
$`(E.4)`$
$`F^0`$ and $`G^0`$ being unknown functions. Then, eq. (5.9) with $`(i,k,l)=(0,1,1)`$ allows us to separate the variables $`K`$ and $`F`$, and its solution, substituted into (E.4), gives the result:
$$\eta ^0=M^0(t,h)\mathrm{sinh}^2(F)/K^2+J^0(t,h)/K^4+L^0(t,h),$$
$`(E.5)`$
where $`M^0`$, $`J^0`$ and $`L^0`$ are new unknown functions. With this, eq. (5.9) corresponding to $`(i,k,l)=(0,2,2)`$ implies $`J^0=0`$, and the one with $`(i,k,l)=(0,0,2)`$ solves as follows:
$$\eta ^1=\frac{3}{5}M^0Kh\mathrm{log}(\mathrm{sinh}(F))\frac{2}{5K}M^0,_h\mathrm{sinh}^2(F)+\frac{K(2\mathrm{cosh}^2(F)1)}{5\mathrm{cosh}(F)\mathrm{sinh}(F)}\eta ^2+F^1(t,h,K),$$
$`(E.6)`$
where $`F^1(t,h,K)`$ is a new unknown function, and $`\eta ^2`$ is still completely unknown. Then, for $`(i,k,l)=(0,0,1)`$, the equation (5.9) has the solution:
$$F^1=\frac{3}{5}M^0Kh\mathrm{log}K+G^1(t,h)K,$$
$`(E.7)`$
where $`G^1(t,h)`$ is a new unknown function.
When (E.6) and (E.7) are substituted into the (0, 0, 0) component of (5.9), an algebraic equation for $`\eta ^2`$ results, whose solution is:
$$\eta ^2=\frac{1}{3\mathrm{cosh}^2(F)+1}\{\frac{5}{3}\frac{M^0,_{hh}}{K^4h}\mathrm{cosh}(F)\mathrm{sinh}^5(F)\frac{5}{3}\frac{L^0,_{hh}}{K^2h}\mathrm{cosh}(F)\mathrm{sinh}^3(F)$$
$$+\frac{35}{6}\frac{M^0}{K^2h}\mathrm{cosh}(F)\mathrm{sinh}^3(F)+\frac{5}{2}\frac{L^0}{h}\mathrm{cosh}(F)\mathrm{sinh}(F)$$
$$+3M^0h\mathrm{cosh}(F)\mathrm{sinh}(F)[\mathrm{log}(\mathrm{sinh}(F))\mathrm{log}K]+5\mathrm{cosh}(F)\mathrm{sinh}(F)G^1$$
$$+\frac{1}{2K^2}M^0,_h\mathrm{cosh}(F)\mathrm{sinh}^3(F)+\frac{5}{2}L^0,_h\mathrm{cosh}(F)\mathrm{sinh}(F)$$
$$5\frac{M^0}{K^2h}\mathrm{cosh}(F)\mathrm{sinh}^3(F)[\mathrm{log}(\mathrm{sinh}(F))\mathrm{log}K]\frac{25}{3}\frac{G^1,_h}{K^2h}\mathrm{cosh}(F)\mathrm{sinh}^3(F)\}.$$
$`(E.8)`$
Both sides of the (1, 1, 1) component of (5.9) become then polynomials in $`(\mathrm{log}K)`$ and $`1/K`$, whose corresponding coefficients have to be respectively equal. The coefficients of $`K^1\mathrm{log}K`$ imply then $`M^0=0`$, and with this, only two other terms remain whose solutions are:
$$L^0=3C(t)h^{5/3}B(t)h,G^1=C(t)h^{8/3}+B(t).$$
$`(E.9)`$
where $`C(t)`$ and $`B(t)`$ are unknown functions. But this results in $`\eta ^2=0`$, $`\eta ^0=L^0`$, $`\eta ^1=G^1K`$. Then, the (1, 0, 0) component of (5.9) implies $`C(t)=0`$, and with this all the remaining equations in (5.9) are fulfilled. Thus the final solution of (5.9) is:
$$\eta ^0=B(t)h,\eta ^1=B(t)K,\eta ^2=0.$$
$`(E.10)`$
With $`\xi =\xi (t)`$ from (E.3) and $`\eta ^i`$ as above, any equation of the set (5.10) implies:
$$B=\mathrm{const},\xi =Bt+A,A=\mathrm{const},$$
$`(E.11)`$
and this satisfies all the remaining equations (5.10) and (5.11). Hence, the general solution of (5.8) – (5.11) is (5.12).
This result was derived under the tacit assumption that the functions $`h`$, $`K`$ and $`F`$ are functionally independent. In the course of solving the equations (4.3), (4.5) and (4.6), relations between these functions may appear. It happens sometimes that such relations are revealed by the symmetry equations as cases in which the symmetry group is larger than in the generic case (see e.g. Ref. 16 where special cases of larger symmetry of a single equation were revealed by the symmetry equations). This possibility has not been investigated for the equations (5.9) – (5.11). However, for the equation (5.8) the solution is always (E.3), even if the functions $`h`$, $`K`$ and $`F`$ are not independent, as shown in the paragraph containing (E.3).
References.
<sup>1</sup> A. Krasiński, J. Math. Phys. 39, 380 (1998).
<sup>2</sup> A. Krasiński, J. Math. Phys. 39, 401 (1998).
<sup>3</sup> A. Krasiński, J. Math. Phys. 39, 2148 (1998).
<sup>4</sup> A. Krasiński, Solutions of the Einstein field equations for a rotating perfect fluid. Part 3: A survey of models of rotating perfect fluid or dust. Preprint, Warsaw 1975.
<sup>5</sup> A. Krasiński, Inhomogeneous cosmological models. Cambridge Unviversity Press 1997.
<sup>6</sup> P. Szekeres, Commun. Math. Phys. 41, 55 (1975).
<sup>7</sup> D. A. Szafron, J. Math. Phys. 18, 1673 (1977).
<sup>8</sup> C. Hellaby, Class. Quant. Grav. 13, 2537 (1996).
<sup>9</sup> H. Stephani, Differential equations; their solution using symmetries. Cambridge University Press 1989.
<sup>10</sup> P. J. Olver, Applications of Lie groups to differential equations, 2nd edition. Springer 1993, pp. 350 – 379.
<sup>11</sup> M. A. H. Mac Callum, in General relativity, an Einstein centenary survey. Edited by S. W. Hawking and W. Israel. Cambridge University Press 1979, p. 552 – 553.
<sup>12</sup> G. F. R. Ellis, J. Math. Phys. 8, 1171 (1967).
<sup>13</sup> A. Krasiński, Gen. Rel. Grav. 25, 165 (1993).
<sup>14</sup> A. Krasiński, The newest release of the Ortocartan set of programs for algebraic calculations in relativity. Submitted for publication.
<sup>15</sup> A. Krasiński, M. Perkowski, The system Ortocartan – user’s manual. Fifth edition, 2000. Available by email or on a diskette.
<sup>16</sup> H. Stephani, J. Phys. A16, 3529 (1983).
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# OKHEP-00-06 Relativistic Coulomb Resummation in QCD
## Acknowledgments
The authors would like to thank A.N. Sissakian, D.V. Shirkov, and O.P. Solovtsova for interest in this work and valuable discussions. Partial support of the work by the US National Science Foundation, grant PHY-9600421, by the US Department of Energy, grant DE-FG-03-98ER41066, and by the RFBR, grants 99-01-00091, 99-02-17727, is gratefully acknowledged. The work of ILS was also supported in part by the University of Oklahoma, through its College of Arts and Science, the Vice President for Research, and the Department of Physics and Astronomy. He thanks the members of the high energy group of the University of Oklahoma for their warm hospitality.
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# 1 Orbifold String Theory
## 1 Orbifold String Theory
In 1985, Dixon, Harvey, Vafa and Witten considered string compactification on Calabi-Yau orbifolds (arising as global quotients $`X/G`$ by a finite group $`G`$) for the purpose of symmetry breaking \[DHVW\]. Later on, orbifold string theory became an important part of string theory. For example, orbifolds provide some of the simplest nontrivial models in string theory. Until very recently, only physical argument for mirror symmetry had been given for orbifold models \[GP\]. In fact, orbifolds are such a popular topic in string theory that a search on hep-th yields more than 200 papers whose title contains orbifold. The reason that orbifold string theory is interesting mathematically is that it contains information which we do not have in the smooth case. Roughly speaking, to have a consistent string theory, string Hilbert space has to contain factors called twisted sectors. Twisted sectors can be viewed as the contribution from singularities. All other quantities such as correlation functions have to contain the contribution from the twisted sectors. So far, the twisted sectors are best understood in the context of conformal field theory. It is our intention to initiate a program to investigate the new geometry and topology of orbifolds caused by the inclusion of twisted sectors. We should mention that the orbifold string theory construction has only been carried out for global quotients. However, it is well-known that most of Calabi-Yau orbifolds are not global quotients. It seems to also be important to be able to construct orbifold string theory over general Calabi-Yau orbifolds. Our orbifold quantum cohomology theory works over arbitrary symplectic or projective orbifolds. We hope that our construction will shed some light on the construction of orbifold string theory for general Calabi-Yau orbifolds.
An orbifold, by definition, is a singular space. One can try to desingularize a Calabi-Yau orbifold by the means of resolution or deformation. To preserve the Calabi-Yau condition, we have to restrict ourselves to the so-called crepant resolutions. It is natural to ask for the relation between orbifold string theory and string theory of its desingularization. In fact, this link provides some of the most interesting mathematics from orbifold string theory. In physics, orbifold string theory and ordinary string theory of its crepant resolution appear to be two members in a family of theories. This strongly suggests that there must be a relation between them. The strongest predication is that they are the same. Indeed, this is what physicists hope for. For quantum cohomology, this translates into the following orbifold string theory prediction:
(1.1) The quantum cohomology of a crepant resolution should be “isomorphic” to the “ quantum cohomology” of the orbifold.
Here, the “quantum cohomology” of the orbifold should be understood as orbifold quantum cohomology. The goal of our project is to establish a mathematical theory for orbifold quantum cohomology.
Actually, the above prediction is false in general. A counterexample is a $`K3`$-surface with ADE-singularities. But this is not the end of the story. Recall that the most general form of mirror symmetry fails for rigid Calabi-Yau 3-folds. But this did not stop research from unearthing the layer and layer of mathematical treasures from mirror symmetry. In fact, it is entirely possible that weaker forms or the current form $`(1.1)`$ for a more restrictive class of orbifolds are still true. The authors believe that this link to crepant resolutions will greatly enrich this subject. Therefore, it is useful to keep this strongest form of prediction in mind for the direction of future research.
The weakest form of the orbifold string theory prediction is to replace quantum cohomology by orbifold Euler number. Here, it has a good chance to hold. The orbifold Euler number is defined as the sum of Euler numbers over all sectors. It has a natural interpretation as the Euler characteristic of orbifold K-theory \[AR\]. A weak orbifold string theory prediction is that Euler number of a crepant resolution is the same as the orbifold Euler number of itself. On the mathematical side, a similar phenomenon was independently discovered earlier by John McKay, which is now known as McKay correspondence \[McK, Re\]. A version of McKay correspondence is stated as follows:
Let $`GSL(n,𝐂)`$ be a finite group, and $`\pi :YX=𝐂^n/G`$ be a crepant resolution, then there exist “natural” bijections between conjugacy classes of $`G`$ and a basis of $`H_{}(Y;𝐙)`$.
Based on these ideas, Batyrev-Dais \[BD\] proposed the so-called strong McKay correspondence and defined string-theoretic Hodge numbers.
The classical part of our orbifold quantum cohomology is a new cohomology of orbifolds which we call orbifold cohomology (see section 2). In the case of Gorenstein orbifolds, Batyrev-Dais’s string-theoretic Hodge number is just the Hodge number of our orbifold cohomology. The next level of the orbifold string theory prediction is to identify the orbifold cohomology group with the ordinary cohomology group of a crepant resolution. This is best described through the orbifold K-theory \[AR\]. It is unlikely that one can identify the cohomology ring structures because of quantum corrections. The third level would be the last level concerning quantum cohomology, which is the most challenging one. At this moment, it is not clear how to formulate the prediction without the risk of finding a simple counterexample.
In the following sections, we shall outline the construction of an orbifold Gromov-Witten theory, which obeys almost all of the axioms of ordinary Gromov-Witten theory, as it should according to physics.
## 2 Orbifold Cohomology
The ordinary quantum cohomology ring appears as a (quantum) deformation of the ordinary cohomology ring with the cup product. In the orbifold Gromov-Witten theory, the role of ordinary cohomologies is played by the so called orbifold cohomologies, which we shall describe in this section. One of our main results is the construction of an orbifold cup product on the total orbifold cohomology group, which makes it into a ring with unit. We will call the resulting ring the orbifold cohomology ring. The orbifold quantum cohomology ring is just the corresponding quantum deformation of the orbifold cohomology ring. Details are in \[CR1\].
Let $`X`$ be a closed, almost complex orbifold of dimension $`n`$, with almost complex structure $`J`$. For any $`pX`$, let $`(V_p,G_p,\pi _p)`$ be a uniformizing system, which can be taken so that $`V_p`$ is a small ball in $`𝐂^n`$ centered at the origin and $`G_p`$ acts on $`V_p`$ as a finite subgroup of $`U(n)`$. We consider the set
$$\stackrel{~}{X}:=\{(p,(g)_{G_p})|pX,gG_p\}.$$
$`(2.1)`$
Here $`(g)_{G_p}`$ represents the conjugacy class of $`g`$ in $`G_p`$. There is a locally constant function $`\iota :\stackrel{~}{X}𝐐`$ defined as follows: write $`g`$ as a diagonal matrix
$$diag(e^{2\pi im_{1,g}/m_g},\mathrm{},e^{2\pi im_{n,g}/m_g}),$$
where $`m_g`$ is the order of $`g`$ in $`G_p`$, and $`0m_{i,g}<m_g`$ for $`i=1,\mathrm{},n`$. We define
$$\iota (p,(g)_{G_p})=\underset{i=1}{\overset{n}{}}\frac{m_{i,g}}{m_g}.$$
$`(2.2)`$
Let $`I:\stackrel{~}{X}\stackrel{~}{X}`$ be the involution defined by $`I(p,(g)_{G_p})=(p,(g^1)_{G_p})`$.
Lemma 2.1: There is an equivalence relation among the $`(g)_{G_p}`$, and if we let $`T=\{(g)\}`$ be the set of equivalence classes and define $`X_{(g)}=\{(p,(g)_{G_p})\stackrel{~}{X}|(g)_{G_p}(g)\}`$, then each $`X_{(g)}`$ is naturally a closed, connected, almost complex orbifold, and $`\stackrel{~}{X}`$ is decomposed as a disjoint union $`_{(g)T}X_{(g)}`$. Furthermore, if we denote the value of the locally constant function $`\iota :\stackrel{~}{X}𝐐`$ by $`\iota _{(g)}`$, and let $`(g^1)`$ denote the image of $`(g)`$ under the involution $`I`$, and $`(1)`$ denote the equivalence class of the trivial element $`(1)_{G_p}`$, the following conditions are satisfied:
1. $`\iota :\stackrel{~}{X}𝐐`$ is integer-valued iff each $`G_p`$ is contained in $`SL(n,𝐂)`$.
2. $$\iota _{(g)}+\iota _{(g^1)}=dim_𝐂Xdim_𝐂X_{(g)}.$$
$`(2.3)`$
3. $`\iota _{(g)}0`$ for all $`(g)T`$, and $`\iota _{(g)}=0`$ iff $`(g)=(1)`$.
Note that for Calabi-Yau orbifolds, each $`\iota _{(g)}`$ is integer-valued. When $`X=Y/G`$ is a global quotient, $`\stackrel{~}{X}`$ can be identified with $`_{(g),gG}Y^g/C(g)`$.
The orbifold cohomologies are just direct sums of ordinary cohomologies of $`X_{(g)}`$ with degrees shifted by $`2\iota _{(g)}`$. More precisely,
Definition 2.2: Let $`X`$ be a closed almost complex orbifold with $`dim_𝐂X=n`$. For any rational number $`d[0,2n]`$, the orbifold cohomology group of degree $`d`$ is defined to be the direct sum
$$H_{orb}^d(X;𝐐)=_{(g)T}H^{d2\iota _{(g)}}(X_{(g)};𝐐).$$
$`(2.4)`$
We will call $`\iota _{(g)}`$ degree shifting numbers, which have been referred as fermion shift numbers in physics \[Z\]. The orbifold $`X_{(g)}`$ or its cohomology will be called a twisted sector if $`(g)(1)`$, and called the nontwisted sector if $`(g)=(1)`$. The construction of $`\stackrel{~}{X}`$ (cf. (2.1)) first appeared in \[Ka\].
The following theorem is proved in \[CR1\], whose construction is based on genus-zero, degree zero orbifold Gromov-Witten invariants.
Theorem 2.3: Let $`(X,J)`$ be a closed almost complex orbifold of dimension $`n`$. Then
1. There is a non-degenerate pairing $`<,>_{orb}:H_{orb}^d(X;𝐐)\times H_{orb}^{2nd}(X;𝐐)𝐐`$ extending the ordinary Poincaré pairing on the nontwisted sectors $`H^{}(X;𝐐)`$.
2. There is a cup product $`_{orb}:H_{orb}^p(X;𝐐)\times H_{orb}^q(X;𝐐)H_{orb}^{p+q}(X;𝐐)`$ for any $`0p,q2n`$ such that $`p+q2n`$, which has the following properties:
* The total orbifold cohomology group $`H_{orb}^{}(X;𝐐)=_{0d2n}H_{orb}^d(X;𝐐)`$ is a ring with unit $`e_X^0H_{orb}^0(X;𝐐)`$ under $`_{orb}`$, where $`e_X^0`$ is the Poincaré dual to the fundamental class $`[X]`$. In particular, $`_{orb}`$ is associative.
* Restricted to each $`H_{orb}^d(X;𝐐)\times H_{orb}^{2nd}(X;𝐐)H_{orb}^{2n}(X;𝐐)`$,
$$_X\alpha _{orb}\beta =<\alpha ,\beta >_{orb}.$$
$`(2.5)`$
* The cup product $`_{orb}`$ is invariant under deformations of $`J`$.
* When $`X`$ is of integral degree shifting numbers, the total orbifold cohomology group $`H_{orb}^{}(X;𝐐)`$ is integrally graded, and we have supercommutativity
$$\alpha _1_{orb}\alpha _2=(1)^{\mathrm{deg}\alpha _1\mathrm{deg}\alpha _2}\alpha _2_{orb}\alpha _1.$$
$`(2.6)`$
* Restricted to the nontwisted sectors, i.e., the ordinary cohomologies $`H^{}(X;𝐐)`$, the cup product $`_{orb}`$ equals the ordinary cup product on $`X`$.
We remark that there is an analogous construction using Dolbeault cohomology groups; for details see \[CR1\].
## 3 Good Map and Pull-Back Bundle
Now we come to one of the main issues in the construction of orbifold Gromov-Witten invariants. Recall that if $`f:XX^{}`$ is a $`C^{\mathrm{}}`$ map between manifolds and $`E`$ is a smooth vector bundle over $`X^{}`$, then there is a smooth pull-back vector bundle $`f^{}E`$ over $`X`$ and a bundle morphism $`\overline{f}:f^{}EE`$ which covers the map $`f`$. However, if instead we have a $`C^{\mathrm{}}`$ map $`\stackrel{~}{f}`$ between orbifolds $`X`$ and $`X^{}`$, and an orbibundle $`E`$ over the orbifold $`X^{}`$, the question whether there is a pull-back orbibundle $`E^{}`$ over $`X`$ and an orbibundle morphism $`\overline{f}:E^{}E`$ covering the map $`\stackrel{~}{f}`$ is a quite complicated issue: $`E^{}`$ might not exist, or even if it exists it might not be unique. Traditionally, a neighborhood of a smooth map into a manifold is described by smooth sections of the pull-back of the tangent bundle of the manifold. Hence understanding this question became the very first step in describing the moduli spaces of pseudo-holomorphic maps, or more precisely, the very first step in order to understand what would be the corresponding notion of “stable map” in the orbifold case.
By a $`C^{\mathrm{}}`$ map between orbifolds $`X`$ and $`X^{}`$ we mean an equivalence class of collections of local smooth liftings between local uniformizing systems of a continuous map from $`X`$ to $`X^{}`$. This notion is equivalent to the notion of V-maps in \[S\], where the notion of orbifold was first introduced under the name V-manifold. A brief review of orbifolds is given in \[CR1\], and a self-contained, elementary discussion of various aspects of differential geometry and global analysis on orbifolds is contained in \[CR2\].
Now we describe our key concept: the notion of good map. Let $`\stackrel{~}{f}:XX^{}`$ be a $`C^{\mathrm{}}`$ map between orbifolds $`X`$ and $`X^{}`$ whose underlying continuous map is denoted by $`f`$. Let $`𝒰=\{U_\alpha ;\alpha \mathrm{\Lambda }\}`$ be an open cover of $`X`$ and $`𝒰^{}=\{U_\alpha ^{};\alpha \mathrm{\Lambda }\}`$ be an open cover of the image $`f(X)`$ in $`X^{}`$, which satisfy the following conditions:
1. Each $`U_\alpha `$ (resp. $`U_\alpha ^{}`$) is uniformized by $`(V_\alpha ,G_\alpha ,\pi _\alpha )`$ (resp. by $`(V_\alpha ^{},G_\alpha ^{},\pi _\alpha ^{})`$).
2. If $`U_\alpha U_\beta `$ (resp. $`U_\alpha ^{}U_\beta ^{}`$), then there is a collection of smooth open embeddings $`(V_\alpha ,G_\alpha ,\pi _\alpha )(V_\beta ,G_\beta ,\pi _\beta )`$ (resp. $`(V_\alpha ^{},G_\alpha ^{},\pi _\alpha ^{})(V_\beta ^{},G_\beta ^{},\pi _\beta ^{})`$), which are called injections.
3. For any point $`pU_\alpha U_\beta `$ (resp. $`p^{}U_\alpha ^{}U_\beta ^{}`$), there is a $`U_\gamma `$ (resp. $`U_\gamma ^{}`$) such that $`pU_\gamma U_\alpha U_\beta `$ (resp. $`p^{}U_\gamma ^{}U_\alpha ^{}U_\beta ^{}`$).
4. Any inclusion $`U_\alpha U_\beta `$ implies $`U_\alpha ^{}U_\beta ^{}`$.
5. For each $`\alpha \mathrm{\Lambda }`$, $`f(U_\alpha )U_\alpha ^{}`$. Moreover, there is a collection of local smooth liftings of $`f`$, $`\{\stackrel{~}{f}_\alpha :V_\alpha V_\alpha ^{};\alpha \mathrm{\Lambda }\}`$ which defines the given $`C^{\mathrm{}}`$ map $`\stackrel{~}{f}`$, such that any injection $`i_{\beta \alpha }:(V_\alpha ,G_\alpha ,\pi _\alpha )(V_\beta ,G_\beta ,\pi _\beta )`$ is assigned with an injection $`\lambda (i_{\beta \alpha }):(V_\alpha ^{},G_\alpha ^{},\pi _\alpha ^{})(V_\beta ^{},G_\beta ^{},\pi _\beta ^{})`$ satisfying the following compatibility conditions:
$$\stackrel{~}{f}_\beta i_{\beta \alpha }=\lambda (i_{\beta \alpha })\stackrel{~}{f}_\alpha \alpha ,\beta \mathrm{\Lambda }$$
$`(3.1)`$
and
$$\lambda (i_{\gamma \beta }i_{\beta \alpha })=\lambda (i_{\gamma \beta })\lambda (i_{\beta \alpha })\alpha ,\beta ,\gamma \mathrm{\Lambda }.$$
$`(3.2)`$
Definition 3.1: We call such a $`(𝒰,𝒰^{},\{\stackrel{~}{f}_\alpha \},\lambda )`$ a compatible system of the $`C^{\mathrm{}}`$ map $`\stackrel{~}{f}`$. A $`C^{\mathrm{}}`$ map is said to be good if it admits a compatible system.
Lemma 3.2: Let $`pr:EX^{}`$ be an orbibundle over $`X^{}`$. For any $`C^{\mathrm{}}`$ good map $`\stackrel{~}{f}:XX^{}`$ with a compatible system $`\xi =(𝒰,𝒰^{},\{\stackrel{~}{f}_\alpha \},\lambda )`$, there is a canonically defined pull-back orbibundle $`pr:E_{\stackrel{~}{f},\xi }^{}X`$ with an orbibundle morphism $`\overline{f}_{\stackrel{~}{f},\xi }:E_{\stackrel{~}{f},\xi }^{}E`$ which covers $`\stackrel{~}{f}`$.
Definition 3.3: Two compatible systems $`\xi _1,\xi _2`$ of a $`C^{\mathrm{}}`$ good map $`\stackrel{~}{f}:XX^{}`$ are said to be isomorphic if for any orbibundle $`E`$ over $`X^{}`$ there is an orbibundle isomorphism $`\varphi :E_{\stackrel{~}{f},\xi _1}^{}E_{\stackrel{~}{f},\xi _2}^{}`$ such that
$$\overline{f}_{\stackrel{~}{f},\xi _1}=\overline{f}_{\stackrel{~}{f},\xi _2}\varphi .$$
$`(3.3)`$
Example 3.4a: Not every $`C^{\mathrm{}}`$ map is good, as shown in the following example: consider an effective linear representation of a finite group $`(𝐑^n,G)`$. Let $`H^g`$ be the linear subspace of fixed points of an element $`1gG`$. Then the centralizer $`C(g)`$ of $`g`$ in $`G`$ acts on $`H^g`$, and $`(H^g,C(g)/K_g)`$ is an effective linear representation, where $`K_gC(g)`$ is the kernel of the action of $`C(g)`$ on $`H^g`$. Suppose $`H^g\{0\}`$ and there is no homomorphism $`\lambda :C(g)/K_gC(g)`$ such that $`\pi \lambda `$ is the identity homomorphism, where $`\pi :C(g)C(g)/K_g`$ is the projection, then the continuous map $`H^g/C(g)𝐑^n/G`$ induced by inclusion $`H^g𝐑^n`$ is a $`C^{\mathrm{}}`$ map which is not a good one.
Example 3.4b: There could be non-isomorphic compatible systems of the same $`C^{\mathrm{}}`$ map, as shown in the following example: Let $`X=𝐂\times 𝐂/G`$ where $`G=𝐙_2𝐙_2`$ acting on $`𝐂\times 𝐂`$ in the standard way. For the $`C^{\mathrm{}}`$ map $`𝐂/𝐙_2X`$ defined by the inclusion $`𝐂\times \{0\}𝐂\times 𝐂`$, there are two non-isomorphic compatible systems $`(\stackrel{~}{f}_i,\lambda _i):(𝐂,𝐙_2)(𝐂\times 𝐂,G)`$, $`i=1,2`$, where $`\lambda _1(1)=(1,0)`$ and $`\lambda _2(1)=(1,1)`$.
In general it is not only difficult to determine whether a given $`C^{\mathrm{}}`$ map is good or not, but also difficult to classify compatible systems of an arbitrary good map by equivalence up to isomorphism. Nevertheless, it is clear that a good map together with an isomorphism class of compatible systems is the object we ought to deal with in the orbifold quantum cohomology theory.
We end this section with a discussion on the case when the domain of a good map is a 2-dimensional orbifold.
Definition-Construction 3.5:
1. Given a marked Riemann surface $`(\mathrm{\Sigma },𝐳)`$ where $`𝐳=(z_1,\mathrm{},z_k)`$ is the set of marked points, we can give a unique orbifold structure to $`\mathrm{\Sigma }`$ by assigning to each marked point $`z_i`$ an integer $`m_i1`$ (note that $`m_i=1`$ is allowed for convenience). We will call $`(\mathrm{\Sigma },𝐳,𝐦)`$ an orbifold marked Riemann surface, where $`𝐦=(m_1,\mathrm{},m_k)`$ is the set of assigned integers, called multiplicities.
2. An orbifold nodal Riemann surface is a marked nodal Riemann surface with the following data: (i) each irreducible component is an orbifold marked Riemann surface (here a nodal point is considered marked on an irreducible component); (ii) two identified nodal points are assigned with the same multiplicity.
Convention-Definition 3.6:
1. Note that, in the definition of compatible systems, the compatibility conditions $`(3.1),(3.2)`$ give rise to a collection of homomorphisms $`\lambda _\alpha ,\alpha \mathrm{\Lambda }`$, between local groups $`G_\alpha `$ and $`G_\alpha ^{}`$, such that each local smooth lifting $`\stackrel{~}{f}_\alpha :V_\alpha V_\alpha ^{}`$ is $`\lambda _\alpha `$-equivariant. For a good $`C^{\mathrm{}}`$ map whose domain is an orbifold marked Riemann surface, we require that each $`\lambda _\alpha `$ be a monomorphism for any of its compatible systems.
2. A good $`C^{\mathrm{}}`$ map with a compatible system from an orbifold nodal Riemann surface into an orbifold $`X`$ is a collection of good $`C^{\mathrm{}}`$ maps with compatible systems defined on its irreducible components which satisfies the following compatibility condition: for each pair of identified nodal points $`z_\nu `$ and $`z_\omega `$, the homomorphisms $`\lambda _{z_\nu }`$ and $`\lambda _{z_\omega }`$ between local groups, which are determined by the corresponding compatible systems, satisfy the equation
$$\lambda _{z_\nu }(x)\lambda _{z_\omega }(x)=1_{G_p}$$
$`(3.4)`$
in $`G_p`$, where $`pX`$ is the image of the identified nodal points $`z_\nu `$ and $`z_\omega `$ under the good $`C^{\mathrm{}}`$ map, and $`x`$ is a generator of the local cyclic group at $`z_\nu `$ and $`z_\omega `$ ($`z_\nu `$ and $`z_\omega `$ have the same multiplicity, hence the same local cyclic group).
Finally we observe that each good $`C^{\mathrm{}}`$ map with a compatible system from an orbifold nodal Riemann surface with $`k`$ marked points into an orbifold $`X`$ determines a point in the space $`\stackrel{~}{X}^k`$, where $`\stackrel{~}{X}=_{(g)T}X_{(g)}`$ (cf. (2.1)), as follows: let the underlying continuous map be $`f`$ and for each marked point $`z_i`$, $`i=1,\mathrm{},k`$, let $`x_i`$ be the positive generator of the cyclic local group at $`z_i`$, and $`\lambda _{z_i}`$ be the homomorphism determined by the given compatible system, then the determined point in $`\stackrel{~}{X}^k`$ is
$$((f(z_1),(\lambda _{z_1}(x_1))_{G_{f(z_1)}}),\mathrm{},(f(z_k),(\lambda _{z_k}(x_k))_{G_{f(z_k)}})).$$
$`(3.5)`$
Let $`𝐱=(X_{(g_1)},\mathrm{},X_{(g_k)})`$ be a connected component in $`\stackrel{~}{X}^k`$. We say that a good map with a compatible system is of type $`𝐱`$ if the point $`(3.5)`$ it determines in $`\stackrel{~}{X}^k`$ lies in the component $`𝐱`$.
## 4 Orbifold Stable Maps
We start with the definition of pseudo-holomorphic map from a Riemann surface into an almost complex orbifold.
Definition 4.1: A pseudo-holomorphic map from a Riemann surface $`(\mathrm{\Sigma },j)`$ into an almost complex orbifold $`(X,J)`$ is a continuous map $`f:\mathrm{\Sigma }X`$ which satisfies the following conditions:
1. For any point $`z\mathrm{\Sigma }`$, there is a disc neighborhood $`D_z`$ of $`z`$ with a branched covering map $`br_z:\stackrel{~}{D}_zD_z`$ given by $`ww^{m_z}`$ (here $`m_z=1`$ is allowed).
2. Let $`p=f(z)`$. There is a local uniformizing system $`(V_p,G_p,\pi _p)`$ of $`X`$ at $`p`$ and a local smooth lifting $`\stackrel{~}{f}_z:\stackrel{~}{D}_zV_p`$ of $`f`$ in the sense that $`fbr_z=\pi _p\stackrel{~}{f}_z`$.
3. $`\stackrel{~}{f}_z`$ is pseudo-holomorphic, i.e., $`d\stackrel{~}{f}_zj=Jd\stackrel{~}{f}_z`$.
Remarks 4.2:
1. When $`(X,J)`$ is a complex orbifold, i.e., $`J`$ is integrable, a pseudo-holomorphic map $`f:(\mathrm{\Sigma },j)(X,J)`$ is just a holomorphic map from $`(\mathrm{\Sigma },j)`$ into the analytic space $`(X,J)`$.
2. For each pseudo-holomorphic map $`f:(\mathrm{\Sigma },j)(X,J)`$, there is a subset of finitely many points $`\{z_1,z_2,\mathrm{},z_k\}\mathrm{\Sigma }`$ such that for any $`z\mathrm{\Sigma }\{z_1,z_2,\mathrm{},z_k\}`$ the multiplicity $`m_z`$ in Definition 4.1-1 equals one (cf. \[HW\]). We will consider pseudo-holomorphic maps from marked Riemann surfaces into $`(X,J)`$. As a convention we will always mark these points $`\{z_1,z_2,\mathrm{},z_k\}`$ where the multiplicity is greater than one.
Given a pseudo-holomorphic map $`f`$ from a marked Riemann surface $`(\mathrm{\Sigma },𝐳)`$ into $`(X,J)`$, where $`𝐳=(z_1,\mathrm{},z_k)`$ is a set of finitely many distinct marked points on $`\mathrm{\Sigma }`$, there is an orbifold structure on $`\mathrm{\Sigma }`$ with singular set contained in $`𝐳`$ such that $`f`$ can be lifted to a good $`C^{\mathrm{}}`$ map $`\stackrel{~}{f}`$. A crucial technical result is summarized in the following
Lemma 4.3: For any pseudo-holomorphic map $`f`$ from a Riemann surface $`\mathrm{\Sigma }`$ of genus $`g`$ with $`k`$ marked points $`𝐳=(z_1,z_2,\mathrm{},z_k)`$ into $`(X,J)`$, there are finitely many orbifold structures on $`\mathrm{\Sigma }`$ whose singular set is contained in $`𝐳`$, and for each of these orbifold structures there are finitely many pairs $`(\stackrel{~}{f},\xi )`$, where $`\stackrel{~}{f}`$ is a good map whose underlying map is $`f`$, and $`\xi `$ is an isomorphism class of compatible systems of $`\stackrel{~}{f}`$. The total number is bounded from above by a constant $`C(X,g,k)`$ depending only on $`X,g,k`$.
Definition 4.4: An orbifold stable map from a marked nodal Riemann surface into an almost complex orbifold $`(X,J)`$ consists of the following data:
1. A continuous map from the marked nodal Riemann surface into $`(X,J)`$ whose restriction to each irreducible component is pseudo-holomorphic.
2. An orbifold structure on the marked nodal Riemann surface so that it becomes an orbifold nodal Riemann surface, and a good map with a compatible system from the orbifold nodal Riemann surface into $`(X,J)`$ with the given underlying continuous map.
3. Stability condition: on each $`S^2`$ or $`T^2`$ component which is mapped into a point in $`X`$ there are at least three or one special points (marked or nodal).
There is an obvious equivalence relation amongst the set of orbifold stable maps. We denote by $`\overline{}_{g,k}(X,J,A,𝐱)`$ the set of all equivalence classes of orbifold stable maps of genus $`g`$, $`k`$ marked points, type $`𝐱`$ and homology class $`A`$ into $`(X,J)`$.
Remark 4.5: In the algebraic setting of Deligne-Mumford stack, a related notion which is called twisted stable map was discussed in \[AV\]. Their twisted stable map was described in the language of category and functor. Our good map was formulated in elementary differential-geometric language. From the first sight, two notions look quite different. However, D. Abramovich kindly informed us that they are actually equivalent \[A\].
Remark 4.6: If $`f:\mathrm{\Sigma }X`$ is a pseudo-holomorphic map whose image intersects the singular locus of $`X`$ at only finitely many points, then there is a unique choice of orbifold structure on $`\mathrm{\Sigma }`$ together with a unique $`(\stackrel{~}{f},\xi )`$, where $`\stackrel{~}{f}`$ is a good map with an isomorphism class of compatible systems $`\xi `$ whose underlying continuous map is $`f`$. If the image of $`f`$ lies completely inside the singular locus, there could be different choices, and they are regarded as different points in the moduli space.
Definition 4.7:
1. An orbifold $`X`$ is symplectic if there is a closed 2-form $`\omega `$ on $`X`$ whose local liftings are non-degenerate.
2. A projective orbifold is a complex orbifold which is a projective variety as an analytic space.
The usual Gromov Compactness Theorem for pseudo-holomorphic maps combined with Lemma 4.3 gives the following
Proposition 4.8: Suppose that $`X`$ is a symplectic or projective orbifold. The moduli space of orbifold stable maps $`\overline{}_{g,k}(X,J,A,𝐱)`$ is a compact metrizable space under a natural topology, whose “virtual dimension” is $`2d`$, where
$$d=c_1(TX)A+(dim_𝐂X3)(1g)+k\iota (𝐱).$$
Here $`\iota (𝐱):=_{i=1}^k\iota _{(g_i)}`$ for $`𝐱=(X_{(g_1)},\mathrm{},X_{(g_k)})`$.
## 5 Orbifold Quantum Cohomology
For any component $`𝐱=(X_{(g_1)},\mathrm{},X_{(g_k)})`$, there are $`k`$ evaluation maps (cf. (3.5))
$$e_i:\overline{}_{g,k}(X,J,A,𝐱)X_{(g_i)},i=1,\mathrm{},k.$$
$`(5.1)`$
For any set of cohomology classes $`\alpha _iH^{2\iota _{(g_i)}}(X_{(g_i)};𝐐)H_{orb}^{}(X;𝐐)`$, $`i=1,\mathrm{},k`$, the orbifold Gromov-Witten invariant is defined as the virtual integral
$$\mathrm{\Psi }_{(g,k,A,𝐱)}^{X,J}(\alpha _1^{l_1},\mathrm{},\alpha _k^{l_k})=_{\overline{}_{g,k}(X,J,A,𝐱)}^{vir}\underset{i=1}{\overset{k}{}}c_1(L_i)^{l_i}e_i^{}\alpha _i,$$
$`(5.2)`$
where $`L_i`$ is the line bundle generated by cotangent space of the $`i`$-th marked point.
When $`g=0`$ and $`A=0`$, the moduli space $`\overline{}_{g,k}(X,J,A,𝐱)`$ admits a very nice and elementary description, based on which we gave an elementary construction of genus zero, degree zero orbifold Gromov-Witten invariants in \[CR1\]. Even in this case, virtual integration is needed where there is an obstruction bundle. The orbifold cup product (cf. Theorem 2.3) is defined through these orbifold Gromov-Witten invariants. In the general case, we need to use the full scope of the virtual integration machinary developed by \[FO\], \[LT\], \[Ru\] and \[Sie\].
Singularities of an orbifold impose additional difficulties in carrying out virtual integration in the orbifold case. Due to the presence of singularities, even on a closed orbifold, the function of injective radius of the exponential map does not have a positive lower bound. As a consequence, it is not known that a neighborhood of a (good) $`C^{\mathrm{}}`$ map into an orbifold can be completely described by $`C^{\mathrm{}}`$ sections of the pull-back tangent bundle via the exponential map. Our approach is a combination of techniques developed in the smooth case. We first construct a local Kuranishi neighborhood for each stable map in $`\overline{}_{g,k}(X,J,A,𝐱)`$, then find finitely many stable maps whose local Kuranishi neighborhoods (although they may have different dimensions) can be patched together to form a “global virtual neighborhood” of $`\overline{}_{g,k}(X,J,A,𝐱)`$ (cf. \[FO\]). This is similar to the constructions of \[FO\], \[LT\]. We carry out the virtual integration over this “global virtual neighborhood” by constructing a system of compatible “Thom forms” (cf. \[Ru\]).
When $`X`$ has a symplectic torus action, the “global virtual neighborhood” can be constructed so that it respects this torus action. The localization theory can be extended to the case of virtual integration. We leave this to another paper.
Main results of this work are summarized in the following
Theorem 5.1: Let $`X`$ be a closed symplectic or projective orbifold. The orbifold Gromov-Witten invariants defined in (5.2) satisfy the quantum cohomology axioms of Witten-Ruan for ordinary Gromov-Witten invariants (cf. \[Ru1\]) except that in the Divisor Axiom, the divisor class is required to be in the nontwisted sector (i.e. in $`H^2(X;𝐐)`$). In the formulation of axioms, the ordinary cup product is replaced by the orbifold cup product $`_{orb}`$ (cf. Theorem 2.3).
As a consequence, we have
Theorem 5.2: Let $`X`$ be a closed symplectic or projective orbifold. With suitable coefficient ring $`𝒞`$, the small quantum product and the big quantum product are well-defined on $`H_{orb}^{}(X;𝐐)𝒞`$, and have properties similar to those of the ordinary quantum cohomology.
## 6 Closing Remarks
What we have accomplished so far is just a tip of iceberg! For example, it is still a difficult problem to compute orbifold quantum cohomology. This requires developing new machinery such as localization and surgery techniques. There are two topics whose natural home should be orbifold. They are birational geometry and mirror symmetry. For birational geometry, recent results in algebraic geometry show that birational transformation can be decomposed as a sequence of wall-crossings in GIT-quotients \[AKMW\], \[HK\], \[W1\], \[W2\]. The latter is naturally in the orbifold category. For mirror symmetry, the Calabi-Yau 3-folds in most of the known examples are crepant resolutions of Calabi-Yau orbifolds. Therefore, it is more natural to consider mirror symmetry for orbifolds. Moreover, the second author believes that orbifold quantum cohomology is different from the quantum cohomology of crepant resolutions. How to formulate mirror symmetry in the categary of Calabi-Yau orbifolds seems to be an extremely interesting problem. Suppose we can do all of these, we are still working only in the so-called type II string theory. There are orbifold versions for other types of string theory (such as heterotic string theory) as well. The amount of new mathematics we can unearth is unimaginable!
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# The Robustness of Inflation to Changes in Super-Planck-Scale Physics
## I Introduction
Most current models of inflation are based on weakly self-coupled scalar matter fields minimally coupled to gravity. In most of these models, the period of inflation lasts for a number of e-foldings much larger than the number needed to solve the problems of standard cosmology . In these cases, the physical length of perturbations of cosmological interest today (those which today correspond to the observed CMB anisotropies and to the large-scale structure) was much smaller than the Planck length at the beginning of inflation. Hence, the approximations which go into the calculation of the spectrum of cosmological perturbations break down. It is then of interest to investigate whether the predictions are sensitive to the unknown super-Planck-scale physics, or whether the resulting spectrum of perturbations is determined only by infrared physics.
An analogous problem arises for black hole evaporation. The original computations of the thermal spectrum from black holes appear to involve mode matching at super-Planck scales. However, in the case of black holes it can be shown that the predictions are in fact insensitive to modifications of the physics at the ultraviolet end.
Our goal is to explore whether and in which cases the spectrum of fluctuations resulting from inflationary cosmology depends on the unknown ultraviolet physics. We will adapt the method of and consider theories obtained by replacing the linear dispersion relation for the linearized fluctuation equations by classes of nonlinear dispersion relations. We find that for the class of dispersion relations introduced by Unruh one recovers a scale-invariant spectrum of fluctuations in the case of exponential inflation. In contrast, for the class of dispersion relations modelled after the one introduced in , the resulting spectrum may be tilted and may include exponential and oscillatory factors if the dispersion relation becomes complex. Such spectra are inconsistent with observations. We thus conclude that the predictions for observables in weakly coupled scalar field models of inflation depend sensitively on hidden assumptions about super-Planck-scale physics.
## II Framework
We will consider the evolution of linear cosmological fluctuations in a spatially flat homogeneous and isotropic Universe. As is well known (see e.g. for a comprehensive review), the evolution equation of scalar and tensor fluctuations in conformal time $`\eta `$ reduces to harmonic oscillator equations with time dependent masses. In the following, we will therefore simply consider the evolution of a scalar field $`\mathrm{\Phi }(\eta ,𝐱)`$ living on space-time. Introducing $`\mu (n,\eta )`$ via the Fourier transform $`\mathrm{\Phi }(\eta ,𝐱)[1/(2\pi )^{3/2}]d𝐧(\mu /a)e^{i𝐧𝐱}`$ \[where $`a(\eta )`$ is the scale factor\], the evolution equation of the mode with a comoving wavenumber $`n`$ becomes
$$\mu ^{\prime \prime }+\left[n^2\frac{a^{\prime \prime }}{a}\right]\mu =\mathrm{\hspace{0.17em}0}.$$
(1)
The corresponding power spectrum $`P_\mathrm{\Phi }(n)`$ is given by
$$n^3P_\mathrm{\Phi }=n^3\left|\frac{\mu }{a}\right|^2.$$
(2)
In order to study the dependence of the predictions for $`P_\mathrm{\Phi }`$ on super-Planck-scale physics, we will modify the linear dispersion relation $`\omega _\mathrm{p}^2=k^2=(n/a)^2`$ (where $`\omega _\mathrm{p}`$ is the physical frequency) for wavenumbers greater than a critical wavenumber $`k_\mathrm{C}`$ by replacing the $`n^2`$ term in (1) with
$$n_{\mathrm{eff}}^2=a^2(\eta )F^2(k)=a^2(\eta )F^2[n/a(\eta )],$$
(3)
where $`F(k)`$ differs significantly from $`k`$ only for $`k>k_\mathrm{C}`$. We see that, in terms of comoving wavenumbers, we obtain a time dependent dispersion relation.
The two classes of dispersion relations we specifically analyze are the one proposed by Unruh and a generalization of the one studied by Corley and Jacobson . The first class is given by
$$F(k)k_\mathrm{C}\mathrm{tanh}^{1/p}[\left(\frac{k}{k_\mathrm{C}}\right)^p],$$
(4)
where $`p`$ is an arbitrary coefficient. For large values of the wave number, this becomes a constant $`k_\mathrm{C}`$ whereas for small values this is a linear law as expected. The second class of dispersion relations is given by
$$F^2(k)=k^2+k^2b_m\left(\frac{k}{k_\mathrm{C}}\right)^{2m},$$
(5)
where $`m`$ is an integer and the coefficients $`b_m`$ are at this stage arbitrary. Note that for negative $`b_m`$ the dispersion relation becomes complex for $`kk_\mathrm{C}`$. The dispersion relations are shown in Fig. 1.
In order to compute the power spectrum $`P_\mathrm{\Phi }(n)`$ of (2) we need to know both the initial conditions for the mode $`\mu (n,\eta )`$ and the subsequent evolution. We first discuss the initial conditions. We want the state at the initial time $`\eta _\mathrm{i}`$ to correspond as closely as possible to our usual physical intuition of a vacuum state. We here study two prescriptions for this. The first is to canonically quantize the field $`\mu (n,\eta )`$ and to demand that the initial state minimizes the energy . The second is to set the state up in the local Minkowski vacuum . In the case of a linear dispersion relation, both prescriptions give the same result. However, for a nonlinear dispersion relation, we obtain different results (which emphasizes the general point that the predictions of inflationary cosmology depend on the initial state chosen ).
Demanding that the initial state minimize the energy yields
$$\mu (\eta =\eta _\mathrm{i})=\frac{1}{\sqrt{2\omega (\eta _\mathrm{i})}},\mu ^{}(\eta =\eta _\mathrm{i})=\pm i\sqrt{\frac{\omega (\eta _\mathrm{i})}{2}},$$
(6)
where $`\omega `$ is the comoving frequency, whereas prescribing the local Minkowski vacuum state gives
$$\mu (\eta _\mathrm{i})=\frac{1}{\sqrt{2n}},\mu ^{}(\eta _\mathrm{i})=\pm i\sqrt{\frac{n}{2}}.$$
(7)
As is apparent from (1), in the case of a linear dispersion relation $`\omega =n`$ and the two prescriptions for the initial state coincide. It should be clear that the choice (6) is certainly the most physical one. The second choice, as mentioned above, illustrates the fact that the final result does depend on the initial conditions.
## III Calculation of Spectra
We now compute the power spectrum $`P_\mathrm{\Phi }(n)`$ for values of $`n`$ with wavenumbers larger than $`k_\mathrm{C}`$ at the initial time $`\eta _\mathrm{i}`$. On such scales, the time interval can be divided into three regions. The first is $`\eta _\mathrm{i}<\eta <\eta _1(n)`$, during which the physical wavenumber exceeds $`k_\mathrm{C}`$. During this time interval, the mode evolution is non-standard. The second interval lasts from the time $`\eta _1(n)`$ to the time $`\eta _2(n)`$ when the wavelength equals the Hubble radius $`l_\mathrm{H}`$. During this time, the solutions of the mode equation (1) are oscillatory since the $`n^2`$ term in the parentheses in (1) dominates over the $`a^{\prime \prime }/a`$ term. In the third period \[$`\eta >\eta _2(n)`$\], the modes are effectively frozen: the non-decaying solution of the mode equation is $`\mu (\eta )a(\eta )`$.
The values of $`\eta _1`$ and $`\eta _2`$ depend on the background evolution. For a power-law inflation model with $`a(\eta )=l_0|\eta |^{1+\beta }`$, where $`\beta `$ is a number with $`\beta 2`$ <sup>*</sup><sup>*</sup>*The value $`\beta =2`$ corresponds to exponential inflation. and $`l_0`$ has the dimension of a length, the values of $`\eta _1`$ and $`\eta _2`$ are given by
$`|\eta _1(n)|`$ $`=`$ $`\left({\displaystyle \frac{n}{2\pi }}{\displaystyle \frac{l_\mathrm{C}}{l_0}}\right)^{1/(1+\beta )}|b_m|^{1/[2m(1+\beta )]},`$ (8)
$`|\eta _2(n)|`$ $`=`$ $`{\displaystyle \frac{2\pi }{n}}|1+\beta |,`$ (9)
where $`l_\mathrm{C}`$ is the wavelength corresponding to $`k_\mathrm{C}`$. By combining the above formula for $`a(\eta )`$ with (8), it follows that
$$a[\eta _2(n)]^2n^{2+2\beta }.$$
(10)
Assuming that the non-decaying modes mix with coefficients of order unity at the times $`\eta _1`$ and $`\eta _2`$, then based on the above observations about the time dependence of $`\mu (n,\eta )`$ in the various time intervals, we obtain the following ‘master formula’ for the power spectrum at late times:
$$n^3P_\mathrm{\Phi }(n)\frac{n^3}{2\omega (\eta _\mathrm{i})}\left|\frac{\mu [n,\eta _1(n)]}{\mu (n,\eta _\mathrm{i})}\right|^2a[\eta _2(n)]^2.$$
(11)
This result is true if the initial state minimizes the energy density. If the initial state is taken to be the local Minkowski vacuum, then $`\omega (\eta _\mathrm{i})`$ on the r.h.s. of (11) needs to be replaced by $`n`$.
In the case of the linear dispersion relation, $`\mu (n,\eta )`$ also oscillates during the first time interval $`\eta _\mathrm{i}<\eta <\eta _1(n)`$. Since in this case $`\omega (\eta _\mathrm{i})n`$, we immediately obtain
$$n^3P_\mathrm{\Phi }(n)n^{4+2\beta },$$
(12)
the ‘standard’ prediction of inflationary cosmology.
In the case of Unruh’s dispersion relation, the mode equation can be solved exactly during the first time interval in the case of exponential inflation ($`\beta =2`$):
$$\mu (n,\eta )=A_1|\eta |^{x_1}+A_2|\eta |^{x_2},$$
(13)
where $`A_1`$ and $`A_2`$ are two constant determined by the initial conditions and where the exponents $`x_1`$ and $`x_2`$ are given by
$$x_{1,2}\frac{1}{2}\pm \frac{1}{2}\sqrt{916\pi ^2\frac{l_0^2}{l_\mathrm{C}^2}}.$$
(14)
Note that both modes are decaying ($`|\mu |\eta ^{1/2}`$ and $`|\eta |0`$). Since $`\eta _1`$ depends on $`n`$ as given in (8), we have
$$\frac{\mu [n,\eta _1(n)]}{\mu (n,\eta _\mathrm{i})}n^{1/2}.$$
(15)
Since for the minimum energy density initial state $`\mu (n,\eta _\mathrm{i})`$ is independent of $`n`$ in the wavelength interval under consideration, we obtain
$$n^3P_\mathrm{\Phi }(n)n^0,$$
(16)
i.e. the same scale-invariant spectrum as in the case of the linear dispersion relation.
In the case of the Corley/Jacobson dispersion relation the result is quite different. Let us consider the case $`b_m<0`$ (complex dispersion relation). Then, in the wavelength regime $`\lambda (\eta _\mathrm{i})l_\mathrm{C}`$ of interest, the mode equation in the first time interval can be solved exactly in terms of modified Bessel functions
$$\mu (n,\eta )|\eta |^{1/2}I_{1/(2b)}(z)|\eta |^{1/2}[z(\eta )]^{1/2}e^{z(\eta )},$$
(17)
where the argument of the Bessel function is $`z(\eta )\gamma |\eta |^b`$ with
$$b\mathrm{\hspace{0.17em}1}m(1+\beta ),\gamma \frac{\sqrt{|b_m|}}{(2\pi )^mb}\left(\frac{l_\mathrm{C}}{l_0}\right)^mn^{m+1}.$$
(18)
Note, in particular, the exponential factor in (17) which depends on $`n`$. This factor does not cancel with any other $`n`$-dependent term in (11) and is thus the root of the exponential dependence of the final power spectrum on $`n`$. Combining (11), (6), (10) and (17) we obtain
$`n^3P_\mathrm{\Phi }(n)`$ $``$ $`n^3n^{1m}n^me^{2z[\eta _1(n)]}n^{2+2\beta }`$ (19)
$``$ $`n^{4+2\beta }e^{2z[\eta _1(n)]},`$ (20)
where the second factor on the r.h.s. of the first line comes from $`\omega (\eta _\mathrm{i})`$, the third and fourth factors stem from the ratio of $`\mu `$, and the final factor from the $`a(\eta _2)`$ term in (11). The careful matching between growing and decaying modes also reveals the presence of an oscillating factor $`\mathrm{cos}^2(n\eta _2n\eta _1\pi /4)`$ in the final power spectrum. However, it should also be noticed that the initial conditions are fixed in a region where the mode function does not oscillate.
In the case of the Corley-Jacobson dispersion relation with $`b_m>0`$, the modified Bessel functions in the mode equation during the first time interval must be replaced by regular Bessel functions. Hence, the exponential factors in the power spectrum disappear and the final result is unchanged, i.e. we recover a scale invariant spectrum.
## IV Discussion and Conclusions
We have studied the robustness of the predictions for the spectrum of cosmological perturbations of weakly coupled inflationary models. The method used was to replace the usual linear dispersion relation by special classes of nonlinear ones, where the nonlinearity is confined to physical wavelengths $`\lambda `$ smaller than some critical length $`l_\mathrm{C}`$. We found that for the class of dispersion relations first introduced by Unruh , the predictions are unchanged. This is connected with the fact that the initial vacuum state evolves adiabatically up to the time $`\eta _1`$ when $`\lambda =l_\mathrm{C}`$ We thank Bill Unruh for pointing this connection out to us.. However, in the case of the dispersion relation modelled after the one used by Corley and Jacobson in the situation where it becomes complex, the resulting spectrum can have oscillations, non-standard tilts and exponential factors which render the resulting theory in conflict with observations. The specific predictions depend on the sign of $`b_m`$, on the value of $`m`$, and on the initial state chosen. The results are summarized in Table 1.
We thus conclude that the predictions in weakly coupled scalar field-driven inflationary models are not robust to changes in the unknown fundamental physics on sub-Planck lengths. This opens up another potentially very interesting link between fundamental physics and observations. Note, however, that in strongly coupled scalar field models of inflation such as the model discussed in , the spectrum of fluctuations is robust to changes in the underlying sub-Planck-length physics.
Acknowledgements
We are grateful to Lev Kofman, Dominik Schwarz, Carsten Van de Bruck and in particular Bill Unruh for stimulating discussions and useful comments. We acknowledge support from the BROWN-CNRS University Accord which made possible the visit of J. M. to Brown during which most of the work on this project was done, and we are grateful to Herb Fried for his efforts to secure this Accord. One of us (R. B.) wishes to thank Bill Unruh for hospitality at the University of British Columbia during the time when this work was completed. J. M. thanks the High Energy Group of Brown University for warm hospitality. The research was supported in part by the U.S. Department of Energy under Contract DE-FG02-91ER40688, TASK A.
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# Next-to-Leading Order SUSY-QCD Predictions for Associated Production of Gauginos and Gluinos
## I Introduction
Weak scale supersymmetry (SUSY) is a theoretically attractive extension of the Standard Model of particle physics. Supersymmetric theories can solve the Higgs hierarchy puzzle , break electroweak symmetry radiatively at low energies , and explain the unification of the gauge couplings at a high energy scale . SUSY introduces a superpartner for each Standard Model particle with the same quantum numbers, but for a difference in spin of 1/2. If supersymmetry were exact, these superparticles would be degenerate in mass with their Standard Model partners. However, SUSY can be broken softly in such a way that its attractive features survive, while the superpartners become heavy enough to evade current limits from collider searches . A powerful and general parametrization for the soft SUSY-breaking terms is provided by the Minimal Supersymmetric Standard Model (MSSM) . Soft breaking mechanisms require that the superpartners remain lighter than a few TeV, and thus there is reason to expect that high energy investigations, such as those at LEP II, Run II of the Fermilab Tevatron, and the CERN Large Hadron Collider (LHC) will discover supersymmetric particles or place strong constraints on supersymmetric models.
The search for direct experimental evidence of supersymmetry at colliders requires a good understanding of theoretical predictions of the cross sections for production of the superparticles. In the case of hadron colliders, where collisions of strongly interacting hadrons are studied, the large strong coupling strength ($`\alpha _S`$) leads to potentially large contributions beyond the leading order (LO) in a perturbation series expansion of the cross section. To have accurate theoretical estimates of production rates and differential cross sections, it is necessary to include corrections at next-to-leading order (NLO) or beyond. Next-to-leading order calculations of the hadroproduction of gluinos and squarks , top squarks , sleptons , and gauginos have been published, including our brief report on associated production of gauginos and gluinos . In this paper, we provide a detailed exposition of our calculation of associated production, and we present new predictions of total and differential cross sections for a variety of assumptions about the superpartner mass spectrum.
Associated production of a gaugino ($`\stackrel{~}{\chi }`$) with a gluino ($`\stackrel{~}{g}`$) or with a squark ($`\stackrel{~}{q}`$) is potentially a very important production mechanism. Associated production processes are semi-weak in that they involve one somewhat smaller coupling constant than the production of a pair of colored sparticles. However, in popular models of SUSY breaking , the mass spectrum favors much lighter masses for the color-neutral, low-lying neutralinos and charginos than for the colored squarks and gluinos. Their lighter mass means that the phase space for production of neutralinos and charginos, and the relevant partonic luminosities, will be greater than that for gluinos and squarks. This effect is potentially decisive at a collider with limited energy, such as the Tevatron. Indeed, as our numerical results show, the extra phase space may more than offset the smaller coupling, and gauginos may be produced more copiously than squarks at the Tevatron. Another point in favor of associated production is the relative simplicity of the final state. For example, the lowest lying neutralino is the (stable) lightest supersymmetric particle (LSP) in supergravity (SUGRA) models , manifest only as missing energy in the events, and it is the second lightest in gauge-mediated models. The charginos and higher mass neutralinos may decay leptonically leaving a lepton signature plus missing transverse energy; relatively clean events ensue. Furthermore, associated production may be the best channel for measurement of the gluino mass .
In this paper we present our complete NLO calculation in SUSY-QCD of the hadroproduction of a gluino in association with a gaugino, including contributions from virtual loops of colored sparticles and particles, and three-particle final states in which light quarks or gluons are emitted. We extract the ultraviolet, infrared, and collinear divergences by use of dimensional regularization and illustrate how they may be absorbed by the usual renormalization and mass factorization procedures. In computing the virtual contributions, we encountered divergent four-point functions that had not been evaluated previously. We use a combined analytic and numerical phase space slicing method to treat the contributions from real emission of light particles. Associated production was calculated at LO some years ago . Our reason to focus first on the $`\stackrel{~}{\chi }`$ plus $`\stackrel{~}{g}`$ final state, rather than the $`\stackrel{~}{\chi }`$ plus $`\stackrel{~}{q}`$ final state, is that the LO cross sections for $`\stackrel{~}{\chi }+\stackrel{~}{g}`$ are 3 to 6 times greater than those for $`\stackrel{~}{\chi }+\stackrel{~}{q}`$ at the energy of the Tevatron when the mass $`m_{\stackrel{~}{g}}=m_{\stackrel{~}{q}}=300`$ GeV, and 6 to 15 times greater when $`m_{\stackrel{~}{g}}=m_{\stackrel{~}{q}}=600`$ GeV. These comparisons are pertinent for the lighter mass neutralinos $`\stackrel{~}{\chi }_{1,2}^0`$ and chargino $`\stackrel{~}{\chi }_1^\pm `$. In obtaining the $`\stackrel{~}{q}`$ cross sections, we sum over five flavors of degenerate squarks and antisquarks.
Our analysis is general in that it is not tied to a particular SUSY breaking model. We can provide cross sections for arbitrary gluino and gaugino masses, and, indeed, the values of the cross sections at Tevatron and LHC energies depend crucially on the sparticle masses. Mass generation in supersymmetry is accomplished in a hidden sector and transmitted to the MSSM fields. In SUGRA models , transmission is through gravitational interactions, while in gauge-mediated and gaugino-mediated models, it occurs through gauge interactions. Anomaly mediated SUSY breaking, also gravitational in origin, is a fourth possibility . In the presentation of predictions for cross sections, we consider illustrative mass spectra typical of each scenario and consistent with bounds established from current data . We also examine the phenomenologically open case of a gluino with mass light compared to the SUSY scale, $`m_{\stackrel{~}{g}}30`$ GeV. This possibility arises in some models of gauge-mediated SUSY breaking .
As is shown in detail below, the NLO corrections to associated production are generally positive, but they can be modest in size, ranging in the SUGRA model from a few percent at the energy of the Fermilab Tevatron to 100% at the energy of the LHC, depending on the sparticle masses. In the light-gluino case, NLO contributions increase the cross section by factors of 1.3 to 1.4 at the energy of the Tevatron and by factors of 2 to 3.5 at the energy of the LHC. Owing to these enhancements, collider searches for signatures of associated production will generally discover or exclude sparticles with masses larger than one would estimate based on LO production rates alone. More significant from the viewpoint of reliability, the renormalization and factorization scale dependence of the cross sections is reduced by a factor of more than two when NLO contributions are included.
At Run II of the Fermilab Tevatron, for an integrated luminosity of 2 $`\mathrm{fb}^1`$, we expect that 10 or more events could be produced in each of the lighter gaugino channels of the SUGRA model, $`\stackrel{~}{g}\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{g}\stackrel{~}{\chi }_2^0`$, and $`\stackrel{~}{g}\stackrel{~}{\chi }_1^\pm `$, provided that the gluino mass $`m_{\stackrel{~}{g}}`$ is less than 450 GeV. The cross sections for the three heavier gaugino channels, $`\stackrel{~}{g}\stackrel{~}{\chi }_3^0`$, $`\stackrel{~}{g}\stackrel{~}{\chi }_4^0`$, and $`\stackrel{~}{g}\stackrel{~}{\chi }_2^\pm `$, are smaller by an order of magnitude or more than those of the lighter gaugino channels. In the light gluino model that we consider, more than 100 events could be produced in the three lighter gaugino channels provided that the common GUT-scale fermion mass $`m_{1/2}`$ is less than 400 GeV, and as many as 10 events in the three heavier gaugino channels as long as $`m_{1/2}`$ is less than 200 GeV. At the higher energy and luminosity of the LHC, at least a few events should be produced in every channel in the SUGRA model and many more in the light gluino model.
The shapes of the rapidity distributions of the gauginos are not altered appreciably by NLO contributions, but the locations of the maximum cross section in transverse momentum ($`p_T`$) are shifted to smaller values by NLO contributions. At LHC energies where the contribution of the $`qg`$ initial state is important, modifications of the $`p_T`$ spectra can be pronounced.
We begin in Sec. II with a brief review of the LO calculation in order to introduce our notation. In this section, we also introduce the SUSY breaking models we adopt and summarize salient aspects of their predicted mass spectra. This discussion is followed in Sec. III by a detailed presentation of our NLO $`𝒪(\alpha \alpha _S^2)`$ calculation. We present partonic scaling functions in Sec. IV as well as predictions for inclusive and differential cross sections at Tevatron and LHC energies. Our conclusions may be found in Sec. V. In Appendices A – E, we present detailed analytic results related to the NLO calculation.
## II Leading Order Production of Gauginos and Gluinos
We begin with the Born level cross sections for the partonic processes
$`q\overline{q}\stackrel{~}{g}\stackrel{~}{\chi }_k^0`$ , $`q\overline{q}\stackrel{~}{g}\stackrel{~}{\chi }_k^\pm ,`$ (1)
derived first in . In anticipation of the renormalization and mass factorization of the NLO contributions, we proceed in the $`n=42ϵ`$ dimensions of standard dimensional regularization. We assume that there is no mixing between squarks of different generations and that the squark mass eigenstates are aligned with the squark chirality states, equivalent to the assumption that the two squarks of a given flavor are degenerate in mass. We ignore the $`n_f=5`$ light quark masses in all of the kinematics and couplings, and thus study the production of gaugino-like charginos and neutralinos, but not the production of Higgsino-like ones . We assume further that the entries in the chargino and neutralino mass matrices are real, and thus that the unitary transformations from the ($`\stackrel{~}{B}`$, $`\stackrel{~}{W_3}`$, $`\stackrel{~}{H_1}`$, $`\stackrel{~}{H_2}`$) and ($`\stackrel{~}{W}^+`$, $`\stackrel{~}{H}^+`$) bases to the ($`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_2^0`$, $`\stackrel{~}{\chi }_3^0`$, $`\stackrel{~}{\chi }_4^0`$) and ($`\stackrel{~}{\chi }_1^+`$, $`\stackrel{~}{\chi }_2^+`$) bases are given by orthogonal matrices. A result of this convention is that it is possible for the mass of one or more neutralinos to be negative inside a polarization sum. The chargino masses are chosen to be positive as may be done for Dirac fermions.
The Dirac matrix $`\gamma _5`$ (or equivalently the projectors $`P_{L(R)}=(1\pm \gamma _5)/2`$) appearing in the gaugino and gluino couplings is treated in the ‘naive’ scheme in which it anti-commutes with all of the other $`\gamma _\mu `$ matrices. This scheme is acceptable at the one loop level for calculations free from anomaly. In evaluating the Feynman diagrams involving Majorana and explicitly charge-conjugated fermions, we have followed the approach described in Ref. .
We express our leading order results in terms of the Mandelstam variables
$`s`$ $`=(p_a+p_b)^2`$ (2)
$`t`$ $`=(p_bp_2)^2`$ $`t_1`$ $`=(p_bp_2)^2m_{1}^{}{}_{}{}^{2}`$ $`t_2`$ $`=(p_bp_2)^2m_{2}^{}{}_{}{}^{2}`$
$`u`$ $`=(p_ap_2)^2`$ $`u_1`$ $`=(p_ap_2)^2m_{1}^{}{}_{}{}^{2}`$ $`u_2`$ $`=(p_ap_2)^2m_{2}^{}{}_{}{}^{2},`$
where $`p_a`$, $`p_b`$, $`p_1`$, and $`p_2`$ refer to the four-momenta of the incoming quark, the incoming anti-quark, the produced gluino, and the produced gaugino, respectively. Variable $`m_1`$ denotes the mass of the gluino and $`m_2`$ that of the gaugino. The incoming partons are treated as massless. The momenta are on mass-shell, $`p_a^2=p_b^2=0`$, $`p_1^2=m_1^2`$, and $`p_2^2=m_2^2`$. The invariants obey the relation $`s+t+u=m_1^2+m_2^2`$.
After the $`n`$-dimensional phase space integration we obtain the lowest order partonic differential cross section,
$`{\displaystyle \frac{d^2\widehat{\sigma }_{ij}^B}{dt_2du_2}}`$ $`=`$ $`{\displaystyle \frac{\pi S_ϵ}{s^2\mathrm{\Gamma }(1ϵ)}}\left[{\displaystyle \frac{t_2u_2m_2^2s}{\mu ^2s}}\right]^ϵ\mathrm{\Theta }(t_2u_2m_2^2s)`$ (4)
$`\times \mathrm{\Theta }(s(m_1+m_2)^2)\delta (s+t+um_1^2m_2^2)\overline{|_{ij}^B|}^2,`$
where $`S_ϵ=(4\pi )^{2+ϵ}`$. The arbitrary scale $`\mu `$ is introduced, as usual, to provide the correct mass dimension for the coupling in $`n`$ dimensions; $`\overline{|_{ij}^B|}^2`$ is the leading order matrix element summed over the colors and helicities of all of the outgoing particles, and averaged over the colors and helicities of the incoming ones. The indices $`(i,j)`$ label the incident partons. For neutralino production at leading order, the partons are quarks and antiquarks of the same flavor. For chargino production, the incident quarks and antiquarks have different flavor. At next-to-leading order, quark gluon initial states contribute also.
As is shown in Fig. 1, the Born matrix element for associated production of gluinos and gauginos proceeds via $`t`$ or $`u`$ channel exchange of a squark. In the case of charged gauginos, only the left-handed chiral squarks participate, whereas neutral gauginos receive contributions from both left- and right-handed chiral squarks. Furthermore, in the case of charged gaugino production, the squarks exchanged in the $`t`$ and $`u`$ channels correspond to different flavors, while in the neutral case the $`t`$ and $`u`$ channel squarks have the same flavor. Under our assumption that the squark mass eigenstates correspond to squarks of definite chirality, the (massless) incoming quark and anti-quark are forced to have a particular helicity, and thus the sets of graphs in which a right-handed chiral squark is exchanged cannot interfere with those mediated by a left-handed chiral squark. The matrix element has the analytic form
$`\overline{|^B|}^2={\displaystyle \frac{8\pi \widehat{\alpha }_S}{9}}\left[{\displaystyle \frac{X_tt_1t_2}{(tm_{\stackrel{~}{q}_t}^2)^2}}{\displaystyle \frac{2X_{tu}sm_1m_2}{(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}+{\displaystyle \frac{X_uu_1u_2}{(um_{\stackrel{~}{q}_u}^2)^2}}\right],`$ (5)
where $`\widehat{\alpha }_S=\widehat{g}_s^2/4\pi `$ is the coupling between gluinos, squarks, and quarks; $`m_{\stackrel{~}{q}_{t(u)}}`$ is the mass of the squark exchanged in the $`t(u)`$ channel graph; and the $`X`$ represent the gaugino interactions with quark and squark.
For production of a neutralino of type $`\stackrel{~}{\chi }_k^0`$, the $`X`$ are
$`X_t=X_{tu}=X_u=2\left|ee_qN_{k\mathrm{\hspace{0.17em}1}}^{}+{\displaystyle \frac{e}{\mathrm{sin}\theta _W\mathrm{cos}\theta _W}}\left(T_qe_q\mathrm{sin}^2\theta _W\right)N_{k\mathrm{\hspace{0.17em}2}}^{}\right|^2.`$ (6)
In the expressions above, $`e`$ is the electron charge, $`\theta _W`$ the weak mixing angle, $`T_q`$ the third component of the weak isospin for the squark, and $`e_q`$ is the charge of the quark in units of $`e`$. For up-type quarks $`e_q=2/3`$ and for down-type quarks $`e_q=1/3`$. The matrix $`N^{}`$ is the transformation from the interaction to mass eigenbasis defined in Ref. . The expressions for production of a positive chargino of type $`\stackrel{~}{\chi }_k^+`$ are
$`X_t`$ $`={\displaystyle \frac{e^2}{\mathrm{sin}^2\theta _W}}|V_{k\mathrm{\hspace{0.17em}1}}|^2,`$ $`X_{tu}`$ $`={\displaystyle \frac{e^2}{\mathrm{sin}^2\theta _W}}\text{Re}(V_{k\mathrm{\hspace{0.17em}1}}U_{k\mathrm{\hspace{0.17em}1}}^{}),`$ $`X_u`$ $`={\displaystyle \frac{e^2}{\mathrm{sin}^2\theta _W}}|U_{k\mathrm{\hspace{0.17em}1}}|^2.`$ (7)
For the negative chargino $`\stackrel{~}{\chi }_k^{}`$ they have the form,
$`X_t`$ $`={\displaystyle \frac{e^2}{\mathrm{sin}^2\theta _W}}|U_{k\mathrm{\hspace{0.17em}1}}|^2,`$ $`X_{tu}`$ $`={\displaystyle \frac{e^2}{\mathrm{sin}^2\theta _W}}\text{Re}(V_{k\mathrm{\hspace{0.17em}1}}^{}U_{k\mathrm{\hspace{0.17em}1}}),`$ $`X_u`$ $`={\displaystyle \frac{e^2}{\mathrm{sin}^2\theta _W}}|V_{k\mathrm{\hspace{0.17em}1}}|^2.`$ (8)
Matrices $`U`$ and $`V`$ are the chargino transformation matrices from interaction to mass eigenstates defined in Ref. .
To compute cross sections for hadroproduction,
$`h_ah_b\stackrel{~}{g}\stackrel{~}{\chi }_k^0X,`$ $`h_ah_b\stackrel{~}{g}\stackrel{~}{\chi }_k^\pm X,`$ (9)
where $`h_a`$ and $`h_b`$ label the incoming hadrons, one must convolve the partonic cross section with the parton distribution functions. In the high energy scattering limit, one may neglect the mass of the incoming hadrons compared with their momenta, and obtain
$`S`$ $`=(P_a+P_b)^2`$ (10)
$`T`$ $`=(P_bp_2)^2`$ $`T_1`$ $`=(P_bp_2)^2m_{1}^{}{}_{}{}^{2}`$ $`T_2`$ $`=(P_bp_2)^2m_{2}^{}{}_{}{}^{2}`$
$`U`$ $`=(P_ap_2)^2`$ $`U_1`$ $`=(P_ap_2)^2m_{1}^{}{}_{}{}^{2}`$ $`U_2`$ $`=(P_ap_2)^2m_{2}^{}{}_{}{}^{2},`$
in which $`P_i`$ indicates the momentum of hadron $`h_i`$. We define $`x_i`$ by the relations
$`s`$ $`=x_ax_bS,`$ $`t_2`$ $`=x_bT_2,`$ $`u_2`$ $`=x_aU_2.`$ (11)
The convolution with the parton distribution functions may be written as
$`{\displaystyle \frac{d^2\sigma }{dT_2dU_2}}(S,T,U,\mu _F^2)=`$ (12)
$`{\displaystyle \underset{i,j=q,\overline{q}}{}}{\displaystyle _{x_a^{}}^1}𝑑x_a{\displaystyle _{x_b^{}}^1}𝑑x_bx_af_i^{h_a}(x_a,\mu _F^2)x_bf_j^{h_b}(x_b,\mu _F^2){\displaystyle \frac{d^2\widehat{\sigma }_{ij}(s,t,u,\mu _F^2)}{dt_2du_2}}.`$ (13)
In this equation $`\mu _F`$ refers to the factorization scale, and $`d^2\widehat{\sigma }_{ij}/dt_2du_2`$ is the hard cross section, equal to the Born cross section at leading order. The lower limits of integration on the convolution are
$`x_a^{}`$ $`=`$ $`{\displaystyle \frac{T_1}{S+U_2}},`$ (14)
$`x_b^{}`$ $`=`$ $`{\displaystyle \frac{x_aU_2m_2^2+m_1^2}{x_aS+T_2}}.`$ (15)
The differential cross section in the transverse momentum ($`p_T`$) and rapidity ($`y`$) of the gaugino is related to the differential cross section in $`U_2`$ and $`T_2`$ by
$`{\displaystyle \frac{d^2\sigma (S,p_T,y,\mu _F^2)}{dp_Tdy}}=2p_TS{\displaystyle \frac{d^2\sigma (S,T,U,\mu _F^2)}{dT_2dU_2}},`$ (16)
with
$`p_T^2`$ $`=`$ $`{\displaystyle \frac{T_2U_2}{S}}m_2^2={\displaystyle \frac{t_2u_2}{s}}m_2^2,`$ (17)
$`y`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{log}\left({\displaystyle \frac{T_2}{U_2}}\right).`$ (18)
The total cross section is obtained by integrating over the full range of transverse momentum and rapidity,
$`\sigma (S,\mu _F^2)`$ $`=`$ $`{\displaystyle _0^{p_T^{max}(0)}}𝑑p_T{\displaystyle _{y^{max}(p_T)}^{y^{max}(p_T)}}𝑑y{\displaystyle \frac{d^2\sigma (S,p_T,y,\mu _F^2)}{dp_Tdy}}`$ (19)
$`=`$ $`{\displaystyle _{y^{max}(0)}^{y^{max}(0)}}𝑑y{\displaystyle _0^{p_T^{max}(y)}}𝑑p_T{\displaystyle \frac{d^2\sigma (S,p_T,y,\mu _F^2)}{dp_Tdy}}.`$ (20)
The limits of integration are
$`p_T^{max}(y)`$ $`=`$ $`{\displaystyle \frac{1}{2\sqrt{S}\text{Cosh}(y)}}\sqrt{\left(S+m_2^2m_1^2\right)^24m_2^2S\text{Cosh}^2(y)},`$ (21)
and
$`y^{max}(p_T)`$ $`=`$ $`\text{ArcCosh}\left({\displaystyle \frac{S+m_2^2m_1^2}{2\sqrt{S(p_T^2+m_2^2)}}}\right).`$ (22)
### A Supersymmetry Breaking Models
The physical gluino and gaugino masses that we use, as well as the gaugino mixing matrices, are based on four popular SUSY breaking models plus a fifth scenario in which the gluino mass is relatively light.
For our default minimal SUGRA scenario , we select the common scalar and fermion masses at the GUT scale to be $`m_0=100`$ GeV and $`m_{1/2}=150`$ GeV. The trilinear coupling $`A_0=300`$ GeV, and the ratio of the Higgs vacuum expectation values, $`\mathrm{tan}\beta =4`$. The absolute value of the Higgs mass parameter $`\mu `$ is fixed by electroweak symmetry breaking, and we choose $`\mu >0`$. (Our sign convention for $`A_0`$ is opposite to that in the ISASUGRA code .) For this scenario, the neutralino masses $`m_{\stackrel{~}{\chi }_{14}^0}`$ are 55, 104, 283, and 309 GeV with $`m_{\stackrel{~}{\chi }_3^0}<0`$ inside a polarization sum. The chargino masses $`m_{\stackrel{~}{\chi }_{1,2}^\pm }`$ are 101 and 308 GeV and therefore almost degenerate with the masses of $`\stackrel{~}{\chi }_{2,4}^0`$. The gluino mass $`m_{\stackrel{~}{g}}`$ is 410 GeV, and the squark mass is 359 GeV. All of these masses are above the exclusion limits established from LEP and Tevatron collider data . Since the gluino and gaugino masses vary principally with $`m_{1/2}`$, we freeze the values of the other four parameters, and we vary $`m_{1/2}`$ over the range 100 to 400 GeV. The squark, gluino, and gaugino masses all increase as $`m_{1/2}`$ increases.
In considering gauge-mediated SUSY breaking (GMSB), we adopt the parameters of model I studied for the SUSY/Higgs Run II workshop , with $`\mathrm{tan}\beta =2.5`$, $`\mu >0`$, one messenger SU(5) generation, and a messenger scale $`M=2\mathrm{\Lambda }`$. Parameter $`\mathrm{\Lambda }`$ is the scale of SUSY breaking. We examine six cases in which $`\mathrm{\Lambda }`$ varies from 40 to 150 TeV. GMSB does not favor associated production at Tevatron energies because it results in a pattern of gaugino masses in which $`M_3(M)/M_2(M)=\alpha _3(M)/\alpha _2(M)`$, where $`M_3`$ and $`M_2`$ are the masses of the gluino and weak gaugino, and $`\alpha _3`$ and $`\alpha _2`$ are the SU(3) and SU(2) gauge couplings. Since $`M`$ is a low scale, $`\alpha _3`$ is still quite strong. The gluino is generally heavy compared to the other gauginos. Selecting $`\mathrm{\Lambda }=40`$ TeV, $`M_1,M_2,\mu =`$ 56.47, 112.8, 241.7 GeV, and we obtain neutralino masses $`m_{\stackrel{~}{\chi }_{14}^0}=`$ 45, 88, 245, and 281 GeV with $`m_{\stackrel{~}{\chi }_3^0}<0`$ inside a polarization sum. The chargino masses $`m_{\stackrel{~}{\chi }_{1,2}^\pm }`$ are 82 and 277 GeV and again almost degenerate with the masses of $`\stackrel{~}{\chi }_{2,4}^0`$. The gluino mass $`m_{\stackrel{~}{g}}=367`$ GeV, and the squark mass is 471 GeV. The spectrum of masses is similar to that of our default SUGRA scenario, and the masses grow as we increase $`\mathrm{\Lambda }`$. For comparable gluino and gaugino masses, we find that LO cross sections are roughly a factor of 5 (3) smaller at Tevatron (LHC) energies than in the SUGRA model, related to the larger squark mass.
Our anomaly mediated (AMSB) scenario is based on the work in Ref. . It is less well-defined in the sense that scalar masses are not understood, and thus the value of $`\mu `$ is not determined through radiative electroweak symmetry breaking. However, $`M_1,M_2`$, and $`M_3`$ are well specified. We fix the squark masses to be 350 GeV and choose $`\mathrm{tan}\beta =2`$. The gaugino masses are controlled by the gravitino mass. The gluino tends to be heavy in this scenario, disfavoring associated production. However, the gluino mass has phase $`\pi `$ relative to $`M_1`$ and $`M_2`$, resulting in constructive interference at LO in the production of $`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_2^0`$, and $`\stackrel{~}{\chi }_4^0`$, and negative interference in production of $`\stackrel{~}{\chi }_3^0`$, in contrast to the SUGRA scenario. In AMSB the lightest neutralino is always a $`\stackrel{~}{W}`$ and has a large coupling to (s)quarks, in contrast to the $`\stackrel{~}{B}`$-like lightest neutralino of the SUGRA model. We vary the gravitino mass parameter $`m_{3/2}`$ from 30 to 60 TeV. For $`M_1,M_2,\mathrm{and}\mu =`$ 272, 80, and -300 GeV, we obtain neutralino masses $`m_{\stackrel{~}{\chi }_{14}^0}=`$ 91, 269, 309, and 371 GeV with $`m_{\stackrel{~}{\chi }_4^0}<0`$ inside a polarization sum. The chargino masses are $`m_{\stackrel{~}{\chi }_{1,2}^\pm }=`$ 91 and 318 GeV. The gluino mass $`m_{\stackrel{~}{g}}=672`$ GeV. The masses grow as we increase $`m_{3/2}`$. For the $`\stackrel{~}{g}\stackrel{~}{\chi }_1^0`$ channel the LO cross section is roughly a factor of 15 larger at Tevatron energies than in the SUGRA model, for comparable masses. However, the fact that the combination of the gluino and neutralino masses exceeds 750 GeV makes this model an unlikely candidate for discovery at the Tevatron.
Gaugino dominated boundary conditions offer another interesting possibility, exemplified by gaugino-mediated SUSY breaking (g̃MSB) . In this class of models, gauge fields propagate freely in the five-dimensional bulk, whereas fermions are confined to one or more four-dimensional hyper-surfaces. Supersymmetry is broken at a distant point in the extra dimension, giving mass to the gauginos, and the gauge interactions communicate this breaking to the scalar fermions as well. We present results based on the model of Ref. , which combines this simple mechanism of SUSY breaking with a model of quark masses and mixings. This model has two input parameters, $`m_{1/2}`$ and the mass of the down-type Higgs at the GUT scale, $`m_{H_d}`$. Fixing, for example, $`m_{H_d}=200`$ GeV and $`m_{1/2}=150`$ GeV (which determines $`\mathrm{tan}\beta `$ 15 ), we find a spectrum with gluino mass 379 GeV, neutralino masses 57, 103, -224, and 249 GeV, chargino masses 101 and 251 GeV, and squark masses of about 330 GeV. This spectrum is similar to the default SUGRA scenario, with the principal differences that the Higgsino masses are somewhat lighter because the non-universal boundary conditions at the GUT scale result (through the requirement of radiative EWSB) in a non-SUGRA $`\mu `$-term at the weak scale. For the light gaugino-like states the cross sections are virtually unchanged with respect to the SUGRA scenario, whereas the LO cross sections for $`\stackrel{~}{\chi }_3^0`$, $`\stackrel{~}{\chi }_4^0`$, and $`\stackrel{~}{\chi }_2^\pm `$ are 7, 5, and 5 times larger than those in the default SUGRA scenario, because of the increased phase space for these states.
An intriguing scenario is that of a light gluino LSP . Gluino masses in the range of 25 to 35 GeV may still be allowed . Since the work of Ref. treats only the strong SUSY sector, we make some assumptions about the weak parameters, respecting LEP limits on neutralino and chargino masses. We choose $`m_{\stackrel{~}{g}}=`$ 30 GeV and $`m_{\stackrel{~}{q}}=`$ 450 GeV. For the weak sector, we adopt masses typical of SUGRA models, discussed above. Since the gluino mass is light, there is much more phase space available, and cross sections for associated production are substantial at Tevatron energies, reaching $`1\mathrm{p}\mathrm{b}^1`$ for a wide range of values of $`m_{1/2}`$.
Because the GMSB, $`\stackrel{~}{g}`$MSB and AMSB cross sections at LO are not too dissimilar from those of the SUGRA case at Tevatron energies, we focus our NLO work on the SUGRA and light gluino models.
## III Next-to-Leading Order Contributions
The next-to-leading order contributions to the associated production of gluinos and gauginos can be separated into virtual corrections that contain internal loops of colored particles, and 2-to-3 parton real emission contributions in which a light gluon or quark is emitted. The kinematics of the virtual contributions are identical to the Born case described in Sec. II whereas the presence of an additional out-going particle in the emission contributions requires integration over a three-body (rather than two-body) phase space. It is useful to separate the real emission contributions into parts in which the additional parton’s energy approaches zero (and thus the three body final state effectively becomes a two-body one) and parts in which the additional parton is hard (energetic). We refer to these two parts as soft emission and hard emission contributions, respectively.
All of these NLO contributions contain singularities. The virtual corrections contain ultraviolet (UV) singularities that may be absorbed into the definitions of the couplings and operators in the usual renormalization procedure. Both virtual and emission contributions contain infrared (IR) singularities when the energy of the produced or exchanged particle approaches zero. These singularities cancel when the virtual and emission contributions are combined. Finally, there are collinear singularities when the produced particle is emitted collinearly with another massless colored object. These singularities are absorbed into the (universal) NLO definition of the parton distribution functions.
### A Virtual Corrections
In this subsection we present the virtual corrections to the associated production of gauginos and gluinos in hadron collisions. They arise from the interference of the Born matrix elements presented in Sec. II with the one-loop amplitudes shown generically in Fig. 2. In these diagrams, the crossed regions indicate contributions from self-energy corrections (Fig. 3) and vertex corrections (Fig. 4) that are present one at a time at next-to-leading order. Additional contributions arise from the box diagrams in Fig. 5. We include the full supersymmetric spectrum of strongly interacting particles in the virtual loops, i.e. squarks and gluinos as well as quarks and gluons.
Since these virtual loop contributions contain ultraviolet and infrared singularities, we regularize the cross sections by computing the phase space and matrix elements in $`n=42ϵ`$ dimensions. We then obtain the virtual differential cross section from
$`{\displaystyle \frac{d^2\widehat{\sigma }_{ij}^V}{dt_2du_2}}`$ $`=`$ $`{\displaystyle \frac{\pi S_ϵ}{s^2\mathrm{\Gamma }(1ϵ)}}\left[{\displaystyle \frac{t_2u_2m_2^2s}{\mu ^2s}}\right]^ϵ\mathrm{\Theta }(t_2u_2m_2^2s)\mathrm{\Theta }(s(m_1+m_2)^2)`$ (24)
$`\times \delta (s+t+um_1^2m_2^2)\overline{(^B^V+^V^B)}.`$
As in the Born case, the matrix elements are summed (averaged) over the colors and spins of the outgoing (incoming) particles.
We calculate the traces of Dirac matrices with the help of the computer algebra program FORM using the so-called “naive” $`\gamma _5`$ scheme. In this scheme, $`\gamma _5`$ anticommutes with all other $`\gamma _\mu `$ matrices, which is justified for anomaly-free one-loop amplitudes . The $`\gamma _5`$ matrix enters the calculation through both the quark-squark-gluino and quark-squark-gaugino Yukawa couplings. The integration over the internal loop momenta is simplified by reducing all tensorial integration kernels to expressions that are only scalar functions of the loop momentum . The resulting one-, two-, three-, and some four-point functions were computed in the context of other physical processes . However, two previously unknown divergent four-point functions are computed here for the first time due to the fact that the final state particles, i.e. the gluino and the gaugino, have different masses in general. The absorptive parts are obtained with Cutkosky cutting rules and the real parts with dispersion techniques. The results are collected in Appendix B.
The virtual one-loop corrections contain ultraviolet divergences that appear as poles in $`1/ϵ`$ in the one- and two-point functions. They are removed by renormalization of the coupling constants in the modified-minimal-subtraction scheme ($`\overline{\mathrm{MS}}`$) scheme at the renormalization scale $`\mu `$ , and of the masses of the heavy particles (squarks and gluinos) in the on-shell scheme. A difficulty arises from the fact that gluons have $`n2`$ possible polarizations, whereas gluinos have 2, leading to violation of supersymmetry in the $`\overline{\mathrm{MS}}`$ scheme. The simplest procedure to restore supersymmetry, which we adopt here, is through a finite shift in the quark-squark-gluino Yukawa coupling:
$$\widehat{g}_s=g_s\left[1+\frac{\alpha _s}{4\pi }\left(\frac{2}{3}N\frac{1}{2}C_F\right)\right]=g_s\left[1+\frac{\alpha _s}{3\pi }\right].$$
(25)
This shift was discussed first in Ref. .
The virtual corrections can be classified into a $`C_F`$ and a $`N_C`$ color class depending on the color flow or the Abelian or non-Abelian nature of the correction vertices. In addition to UV singularities they have collinear and infrared singularities that appear as $`1/ϵ`$ or $`1/ϵ^2`$ poles in the derivatives of the two-point- and in the three- and four-point functions. The generally divergent scalar integrals are always multiplied by finite coefficient functions proportional to parts of the Born matrix elements. The full result is given in Appendix A.
### B Real Emission Contributions
At NLO, the production of gluinos and gauginos receives contributions from real emission of gluons or massless quarks and anti-quarks. In the following sub-sections both of these types of two-to-three partonic contributions are dealt with separately.
Following the notation developed in Ref. , we express our results in terms of the following sets of invariants,
$`s`$ $`=(p_a+p_b)^2`$ $`s_5`$ $`=(p_1+p_2)^2`$ (26)
$`s_3`$ $`=(p_3+p_2)^2m_1^2`$ $`s_4`$ $`=(p_3+p_1)^2m_1^2`$ (27)
$`t`$ $`=(p_bp_2)^2`$ $`t^{}`$ $`=(p_bp_3)^2`$ (28)
$`u`$ $`=(p_ap_2)^2`$ $`u^{}`$ $`=(p_ap_3)^2`$ (29)
$`u_6`$ $`=(p_bp_1)^2m_1^2`$ $`u_7`$ $`=(p_ap_1)^2m_1^2,`$ (30)
where $`p_a`$ and $`p_b`$ are the four-momenta of the incoming (massless) partons, $`p_1`$ and $`p_2`$ are the $`\stackrel{~}{g}`$ and $`\stackrel{~}{\chi }`$ momenta, and $`p_3`$ is the momentum of the additional massless parton. We also find it useful to define the following derived quantities:
$`t_1`$ $`=tm_1^2`$ $`t_2`$ $`=tm_2^2`$ (31)
$`u_1`$ $`=um_1^2`$ $`u_2`$ $`=um_2^2`$ (32)
$`\mathrm{\Delta }_u`$ $`=m_1^2m_{\stackrel{~}{q}_u}^2`$ $`\mathrm{\Delta }_t`$ $`=m_1^2m_{\stackrel{~}{q}_t}^2`$ (33)
$`u_{6\mathrm{\Delta }}`$ $`=u_6+\mathrm{\Delta }_u`$ $`u_{7\mathrm{\Delta }}`$ $`=u_7+\mathrm{\Delta }_t`$ (34)
$`s_{4\mathrm{\Delta }}`$ $`=s_4+\mathrm{\Delta }_u`$ $`s_{3\mathrm{\Delta }}`$ $`=s_3+\mathrm{\Delta }_t.`$ (35)
Energy and momentum conservation provide relations among these quantities:
$`s_4`$ $`=s+t_2+u_1`$ $`s_3`$ $`=s+u_6+u_7`$ (36)
$`s_5`$ $`=s+t^{}+u^{}`$ $`u_6`$ $`=st_2t^{}`$ (37)
$`u_7`$ $`=su_2u^{}.`$ (38)
The $`n`$-dimensional three-body phase space may be derived conveniently if we integrate the general fully differential cross section in the 1-3 rest frame . In this frame the $`4`$-dimensional components of the $`n`$-dimensional momenta are expressed as:
$`p_a`$ $`=`$ $`(\omega _a,0,\omega _a\mathrm{sin}\psi ,\omega _a\mathrm{cos}\psi )`$ (39)
$`p_b`$ $`=`$ $`(\omega _b,,0,0,\omega _b)`$ (40)
$`p_1`$ $`=`$ $`(E_1,\omega _3\mathrm{sin}\theta _1\mathrm{sin}\theta _2,\omega _3\mathrm{sin}\theta _1\mathrm{cos}\theta _2,\omega _3\mathrm{cos}\theta _1)`$ (41)
$`p_2`$ $`=`$ $`(E_2,\omega _a\mathrm{sin}\psi ,\omega _a\mathrm{cos}\psi +\omega _b)`$ (42)
$`p_3`$ $`=`$ $`(\omega _3,\omega _3\mathrm{sin}\theta _1\mathrm{sin}\theta _2,\omega _3\mathrm{sin}\theta _1\mathrm{cos}\theta _2,\omega _3\mathrm{cos}\theta _1),`$ (43)
with
$`\omega _a`$ $`={\displaystyle \frac{s+u_2}{2\sqrt{s_4+m_1^2}}}`$ $`\omega _b`$ $`={\displaystyle \frac{s+t_2}{2\sqrt{s_4+m_1^2}}}`$ (44)
$`\omega _3`$ $`={\displaystyle \frac{s_4}{2\sqrt{s_4+m_1^2}}}`$ $`E_1`$ $`={\displaystyle \frac{s_4+2m_1^2}{2\sqrt{s_4+m_1^2}}}`$ (45)
$`E_2`$ $`={\displaystyle \frac{t_2+u_2+2m_1^2}{2\sqrt{s_4+m_1^2}}}`$ $`\mathrm{cos}\psi `$ $`={\displaystyle \frac{s(s_4+m_1^2+m_2^2)+t_2u_2}{(s+t_2)(s+u_2)}}.`$ (46)
Using this parameterization, we may express the invariants defined in Eq. (26) in terms of $`\theta _1`$, $`\theta _2`$, and the $`\theta `$-independent variables $`\omega _{(a,b,3)}`$, $`E_{(1,2)}`$, and $`\psi `$. For the real emission contributions, it is sometimes convenient to parameterize these momenta with the $`\widehat{z}`$ axis aligned along $`p_a`$ or $`p_2`$. As these alternate frames are related by a simple spatial rotation, the expressions for $`E_{(1,2)}`$, and $`\omega _{(a,b,3)}`$ remain unchanged. The general three-body cross section may be expressed in this frame as
$`{\displaystyle \frac{d^3\widehat{\sigma }_{ij}^R}{ds_4du_2dt_2}}`$ $`=`$ $`{\displaystyle \frac{S_{ϵ}^{}{}_{}{}^{2}\mu ^{2ϵ}}{2s^2\mathrm{\Gamma }(12ϵ)}}\left[{\displaystyle \frac{t_2u_2sm_2^2}{s\mu ^2}}\right]^ϵ\mathrm{\Theta }(t_2u_2sm_2^2)\mathrm{\Theta }(s_4)`$ (47)
$`\times `$ $`\mathrm{\Theta }(s(m_1+m_2)^2){\displaystyle \frac{s_4^{12ϵ}}{(s_4+m_1^2)^{1ϵ}}}\delta (s+t_2+u_1s_4){\displaystyle 𝑑\mathrm{\Omega }_n\overline{|^R|}^2},`$ (48)
in which $`\overline{|^R|}^2`$ is the real emission matrix element squared, summed over final spins and colors and averaged over initial spins and colors, and the $`n`$-dimensional angular integration is $`d\mathrm{\Omega }_n=\mathrm{sin}^{12ϵ}(\theta _1)d\theta _1\mathrm{sin}^{2ϵ}(\theta _2)d\theta _2`$.
In evaluating the integration over the angular variables in Eq. (47), we follow the procedure outlined in Ref. , in which we use the relations among the invariants, Eq. (36), to reduce all of the angular integrals to the form,
$`I_n^{(k,l)}`$ $`=`$ $`{\displaystyle _0^\pi }\mathrm{sin}^{12ϵ}(\theta _1)𝑑\theta _1{\displaystyle _0^\pi }\mathrm{sin}^{2ϵ}(\theta _2)𝑑\theta _2`$ (50)
$`\times (a+b\mathrm{cos}\theta _1)^k(A+B\mathrm{cos}\theta _1+C\mathrm{sin}\theta _1\mathrm{cos}\theta _2)^l,`$
the necessary expressions for which may be found in Ref. . The angular integrations involving negative powers of $`t^{}`$ and $`u^{}`$ produce poles in $`ϵ`$ which correspond to collinear singularities in which particle 3 is collinear with particle $`a`$ or $`b`$ (c.f. Figs. (6) and (7)). Because these singularities follow a universal structure, they may be removed from the cross section and absorbed into the parton distribution functions according to the usual mass factorization procedure . The non-zero mass of the gluino kinematically forbids collinear emission, and thus the gluino has no associated collinear singularities.
The collinear singular pieces have the factorized form
$`{\displaystyle \frac{d^2\widehat{\sigma }_{ij}^R(s,t_2,u_2,\mu ^2)}{dt_2du_2}}`$ $`=`$ $`{\displaystyle _0^1}𝑑x_1x_1{\displaystyle _0^1}𝑑x_2x_2{\displaystyle \underset{k,l}{}}\mathrm{\Gamma }_{ki}(x_1,\mu _F^2,\mu ^2,ϵ)`$ (52)
$`\times \mathrm{\Gamma }_{lj}(x_2,\mu _F^2,\mu ^2,ϵ){\displaystyle \frac{d^2\widehat{\sigma }_{kl}^R(\widehat{s},\widehat{t}_2,\widehat{u}_2,\mu _F^2)}{d\widehat{t}_2d\widehat{u}_2}},`$
in which $`\widehat{s}=x_1x_2s`$, $`\widehat{u}_2=x_1u_2`$, $`\widehat{t}_2=x_2t_2`$. The universal splitting functions $`\mathrm{\Gamma }_{ij}`$ contain the collinear divergences associated with incoming parton $`j`$ splitting into parton $`i`$, and the hard scattering cross section, $`d^2\widehat{\sigma }_{kl}^R/d\widehat{t}_2d\widehat{u}_2`$, is free from singularities. The splitting functions may be redefined by an arbitrary finite term, and thus one must choose a factorization scheme. In order to use recent sets of parton distributions extracted from data we adopt the $`\overline{\mathrm{MS}}`$ scheme, in which the splitting functions at $`𝒪(\alpha _S)`$ are
$`\mathrm{\Gamma }_{ij}(x,\mu _F^2,\mu ^2,ϵ)=\delta _{ij}\delta (x1)+{\displaystyle \frac{\alpha _S}{2\pi }}\left[{\displaystyle \frac{1}{ϵ}}+\gamma _E\mathrm{log}(4\pi )+\mathrm{log}\left({\displaystyle \frac{\mu _F^2}{\mu ^2}}\right)\right]P_{ij}(x).`$ (53)
In the above expression, the $`P_{ij}(x)`$ are the Altarelli-Parisi evolution kernels ,
$`P_{qq}(x_i)=P_{\overline{q}\overline{q}}(x_i)`$ $`=`$ $`C_F\left[{\displaystyle \frac{1+x_i^2}{1x_i}}\mathrm{\Theta }(1\delta _ix_i)+\left(2\mathrm{log}\delta _i+{\displaystyle \frac{3}{2}}\right)\delta (1x_i)\right],`$ (54)
$`P_{gq}(x_i)=P_{g\overline{q}}(x_i)`$ $`=`$ $`C_F{\displaystyle \frac{1+(1x_i)^2}{x_i}},`$ (55)
$`P_{qg}(x_i)=P_{\overline{q}g}(x_i)`$ $`=`$ $`T_F\left[x_i^2+(1x_i)^2\right],`$ (56)
$`P_{gg}(x_i)`$ $`=`$ $`2N_C\left[{\displaystyle \frac{1}{x_i(1x_i)}}+x_i(1x_i)2\right]\mathrm{\Theta }(1x_i\delta _i)`$ (58)
$`+\left[2N_C\mathrm{log}\delta _i+{\displaystyle \frac{1}{2}}\beta _0^L\right]\delta (1x_i),`$
with $`C_F=4/3`$ and $`T_F=1/2`$ for $`N_C=3`$ colors of quarks; $`\beta _0^L=11N_C/62n_FT_F/3`$. The quantities $`\delta _i`$ express the slicing of $`s_4`$ into hard and soft regimes in terms of the $`x_a`$ and $`x_b`$ variables. They can be related to the $`\mathrm{\Delta }`$ of Section III B 1 by $`\delta _a=\mathrm{\Delta }/(s+u_2)`$ and $`\delta _b=\mathrm{\Delta }/(s+t_2)`$.
We set the renormalization and factorization scales equal to each other, $`\mu _F=\mu `$, and expand Eq. (52) to $`𝒪(\alpha _S)`$ to derive the expression for the reduced cross section,
$`{\displaystyle \frac{d^2\widehat{\sigma }_{ij}^R(s,t_2,u_2,\mu ^2)}{dt_2du_2}}`$ $`=`$ $`{\displaystyle \frac{d^2\widehat{\sigma }_{ij}^R(s,t_2,u_2,\mu ^2)}{dt_2du_2}}`$ (61)
$`+{\displaystyle \frac{\alpha _S}{2\pi }}{\displaystyle \frac{1}{\overline{ϵ}}}{\displaystyle _0^1}𝑑x_1x_1P_{li}(x_1){\displaystyle \frac{d^2\widehat{\sigma }_{lj}^0(x_1s,t_2,x_1u_2,\mu ^2)}{dt_2d\widehat{u}_2}}`$
$`+{\displaystyle \frac{\alpha _S}{2\pi }}{\displaystyle \frac{1}{\overline{ϵ}}}{\displaystyle _0^1}𝑑x_2x_2P_{kj}(x_2){\displaystyle \frac{d^2\widehat{\sigma }_{ik}^0(x_2s,x_2t_2,u_2,\mu ^2)}{d\widehat{t}_2du_2}}.`$
We employ the compact notation $`\overline{ϵ}^1=ϵ^1\gamma _E+\mathrm{log}(4\pi )`$; $`d^2\widehat{\sigma }_{ik}^0/dt_2du_2`$ is the leading order cross section for $`ik\stackrel{~}{g}\stackrel{~}{\chi }`$. The resulting hard scattering cross section is free from collinear singularities, as the implicit $`ϵ`$-dependence of $`d^2\widehat{\sigma }_{ij}/dt_2du_2`$ cancels with the explicit $`ϵ`$-dependence of the second and third terms.
#### 1 Gluon Emission
The NLO real contributions with an additional gluon in the final state,
$`q\overline{q}g\stackrel{~}{g}\stackrel{~}{\chi },`$ (62)
proceed from the Feynman diagrams shown in Fig. 6. As was the case for the leading order cross section for production of neutralinos with gluinos, the set of graphs in which a right-handed squark is exchanged cannot interfere with the graphs in which a left-handed squark is exchanged because the incoming quark and anti-quark must have definite helicity. Production of charginos in association with gluinos involves only left-handed squarks.
In addition to the collinear singularities described above, this set of corrections also has infrared singularities that arise when the energy of the emitted gluon approaches zero. These singularities appear as poles in $`s_4`$ in the reduced cross section, and must also be extracted so that they can be combined with corresponding terms in the virtual corrections and shown to cancel.
To make this cancellation conveniently, we slice the gluon emission corrections into hard and soft pieces,
$$\frac{d^2\widehat{\sigma }_{ij}^R}{dt_2du_2}=_0^\mathrm{\Delta }𝑑s_4\frac{d^3\widehat{\sigma }_{ij}^S}{dt_2du_2ds_4}+_\mathrm{\Delta }^{s_{4}^{}{}_{}{}^{max}}𝑑s_4\frac{d^3\widehat{\sigma }_{ij}^H}{dt_2du_2ds_4},$$
(63)
where $`\mathrm{\Delta }`$ is an arbitrary cut-off between what we call soft gluon radiation and hard gluon radiation. When the cut-off is much smaller than the other invariants, $`\mathrm{\Delta }s,t,u,m_i^2`$, the $`s_4`$ integration for the soft term becomes very simple and can be evaluated analytically. This operation results in singular terms
$`{\displaystyle \frac{d^2\widehat{\sigma }_{ij}^B}{dt_2du_2}}`$ $`[\left({\displaystyle \frac{C_F\alpha _S}{\pi }}\right)\{{\displaystyle \frac{1}{\overline{ϵ}^2}}+({\displaystyle \frac{3}{2}}+\mathrm{log}{\displaystyle \frac{\mu ^2}{s}}){\displaystyle \frac{1}{\overline{ϵ}}}\}`$ (65)
$`\left({\displaystyle \frac{N_C\alpha _S}{2\pi }}\right)\{\mathrm{log}\left({\displaystyle \frac{(s+t_2)(s+u_2)}{sm_1^2}}\right)1\}{\displaystyle \frac{1}{\overline{ϵ}}}].`$
In Eq. (65) we distinguish the contributions from the $`N_C`$ and $`C_F`$ color classes.
The singular expression may then be combined with the virtual corrections discussed in Sec. III A to yield the combined “soft and virtual” contribution free from infrared singularities. The residual finite soft contributions are presented in Appendix C.
In the hard gluon regime, there are collinear singularities, but no IR singularities, and thus the most singular terms are proportional to $`ϵ^1`$. After the mass subtraction described above is performed, the results are singularity-free, and they can be presented as minimally subtracted, singularity-free integrals,
$`\widehat{I}\left(f(\theta _1,\theta _2)\right)`$ $`=`$ $`{\displaystyle _0^\pi }𝑑\theta _1{\displaystyle _0^\pi }𝑑\theta _2\mathrm{sin}^{12ϵ}(\theta _1)\mathrm{sin}^{2ϵ}(\theta _2)f(\theta _1,\theta _2)`$ (67)
$`{\displaystyle \frac{1}{ϵ}}\underset{ϵ0}{lim}\left(ϵ{\displaystyle _0^\pi }𝑑\theta _1{\displaystyle _0^\pi }𝑑\theta _2\mathrm{sin}^{12ϵ}(\theta _1)\mathrm{sin}^{2ϵ}(\theta _2)f(\theta _1,\theta _2)\right),`$
which contain only the finite terms with the $`ϵ^1`$ poles subtracted. The resulting expression consists of a simple power series in $`ϵ`$, which may then be evaluated in 4 dimensions by setting $`ϵ0`$. Note that the function $`f(\theta _1,\theta _2)`$ can involve coefficients for angular expressions that have mass dimension, and thus the mass dimension of $`\widehat{I}(f(\theta _1,\theta _2))`$ will depend on $`f(\theta _1,\theta _2)`$. The gluon emission matrix elements are presented in Appendix D.
The cutoff on the $`s_4`$ integration introduces an implicit logarithmic dependence on $`\mathrm{\Delta }`$ that is matched by the explicit logarithms of $`\mathrm{\Delta }`$ which appear in the combined soft and virtual term. The total correction is independent of the value of $`\mathrm{\Delta }`$. Choosing for illustration $`m_{1/2}=400`$ GeV, we display in Fig. 8 the dependence of various contributions on our cutoff $`\delta `$. The Born contribution is obviously independent of $`\delta `$, but its contribution helps to show the relative magnitude of different terms. The combined soft and virtual contribution is positive but falls as an explicit analytic function of $`\mathrm{log}\delta `$. The hard part of the gluon emission contribution is negative, but its numerical value grows more positive as a implicit function of $`\mathrm{log}\delta `$. The figure shows that two contributions balance each other well, and the combined soft/virtual plus hard contribution is independent of the cutoff for $`\delta <210^3`$. The figure also shows that the net small next-to-leading order contribution is obtained after large cancellations take place. The case chosen for display in Fig. 8 is a worst case. With $`m_{1/2}=400`$ GeV, $`m_{\stackrel{~}{\chi }_4^0}=679`$ GeV, and $`m_{\stackrel{~}{g}}=1012`$ GeV. The energy chosen is that of the Tevatron, $`\sqrt{S}=2`$ TeV, so phase space limitations are relatively severe for this set of masses. For all other cases, the cross sections are independent of $`\delta `$ for a larger range of $`\delta `$.
#### 2 Light Quark (Anti-Quark) Emission
A second set of real emission corrections involves an additional light quark (or anti-quark) in the final state, through the partonic reactions,
$`qgq\stackrel{~}{g}\stackrel{~}{\chi }`$ , $`\overline{q}g\overline{q}\stackrel{~}{g}\stackrel{~}{\chi }.`$ (68)
The set of Feynman diagrams contributing to emission of a quark is shown in Fig. 7. They include diagrams in which an incoming gluon splits into a $`q\overline{q}`$ pair as well as diagrams in which an intermediate squark splits into a quark and either a gluino or gaugino. This set of corrections does not have an IR divergence, and thus it is not necessary to slice it into hard and soft regimes. However, after all of the initial state collinear singularities are removed by the mass factorization procedure described above, the matrix elements may still contain integrable singularities if the mass of the squark is larger than the mass of the gluino or gaugino. In these cases, the intermediate squark state can be on its mass-shell, and the variables $`s_{4\mathrm{\Delta }}`$ and $`s_{3\mathrm{\Delta }}`$ go to zero inside the region of integration. This problem was encountered previously , and we follow the same procedure. These singularities represent the LO production of a squark and a gluino or gaugino, followed by the LO decay of the squark. They may be removed if one includes the full Breit-Wigner form for the squark propagator, which regulates the squark resonance by the squark width. This procedure amounts to the replacements,
$`{\displaystyle \frac{1}{s_{4\mathrm{\Delta }}}}`$ $``$ $`{\displaystyle \frac{1}{s_{4\mathrm{\Delta }}+im_{\stackrel{~}{q}_u}\mathrm{\Gamma }_{\stackrel{~}{q}_u}}}𝒫\left({\displaystyle \frac{1}{s_{4\mathrm{\Delta }}}}\right)i\pi \delta (s_{4\mathrm{\Delta }})`$ (69)
$`{\displaystyle \frac{1}{s_{3\mathrm{\Delta }}}}`$ $``$ $`{\displaystyle \frac{1}{s_{3\mathrm{\Delta }}+im_{\stackrel{~}{q}_t}\mathrm{\Gamma }_{\stackrel{~}{q}_t}}}𝒫\left({\displaystyle \frac{1}{s_{3\mathrm{\Delta }}}}\right)i\pi \delta (s_{3\mathrm{\Delta }}),`$ (70)
where $`𝒫`$ indicates the principal value function, and the final distribution identity holds in the limit of small squark widths, $`\mathrm{\Gamma }_{\stackrel{~}{q}}m_{\stackrel{~}{q}}`$. The replacement removes the singularities, and when both $`s_{4\mathrm{\Delta }}`$ and $`s_{3\mathrm{\Delta }}`$ are zero it generates an additional real term from the product of $`\delta `$ functions.
There is a further subtlety associated with the requirement that we not double-count the region of phase space in which the squark is on-shell. Properly, the kinematic configuration with an on-shell squark is included in the LO production of a squark and a gluino or a squark and a gaugino, and thus should not be considered as a genuine higher order correction to the production of gluinos with gauginos. To avoid double counting, we thus subtract the on-shell squark contribution by defining the total cross section (for illustration, we deal with the $`s_{4\mathrm{\Delta }}`$ singular case),
$`\widehat{\sigma }={\displaystyle _0^{s_4^{max}}}𝑑s_4{\displaystyle _{t_2^{}(s_4)}^{t_2^+(s_4)}}𝑑t_2{\displaystyle \frac{d^2\widehat{\sigma }}{dt_2ds_4}}={\displaystyle _0^{s_4^{max}}}𝑑s_4{\displaystyle \frac{f(s_{4\mathrm{\Delta }})}{s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}}^2\mathrm{\Gamma }_{\stackrel{~}{q}}^2}}.`$ (71)
The on-shell contribution then corresponds to $`f(0)/(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}}^2\mathrm{\Gamma }_{\stackrel{~}{q}}^2)`$, with
$`f(0)=\widehat{\sigma }_{\stackrel{~}{\chi }\stackrel{~}{q}}^B{\displaystyle \frac{m_{\stackrel{~}{q}}\mathrm{\Gamma }_{\stackrel{~}{q}}}{\pi }}{\displaystyle \frac{\mathrm{\Gamma }_{\stackrel{~}{q}q\stackrel{~}{g}}^B}{\mathrm{\Gamma }_{\stackrel{~}{q}}}}\widehat{\sigma }_{\stackrel{~}{\chi }\stackrel{~}{q}}^BBR(\stackrel{~}{q}q\stackrel{~}{g})(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}}^2\mathrm{\Gamma }_{\stackrel{~}{q}}^2)\delta (s_{4\mathrm{\Delta }}).`$ (72)
It can be subtracted leaving a genuine NLO contribution
$`\widehat{\sigma }^{\mathrm{NLO}}`$ $`=`$ $`{\displaystyle _0^{s_4^{max}}}𝑑s_4{\displaystyle \frac{f(s_{4\mathrm{\Delta }})f(0)}{s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}}^2\mathrm{\Gamma }_{\stackrel{~}{q}}^2}}`$ (73)
which may once again be expressed as a principal value function, since $`s_{4\mathrm{\Delta }}/(s_{4\mathrm{\Delta }}+m_{\stackrel{~}{q}}^2\mathrm{\Gamma }_{\stackrel{~}{q}}^2)𝒫(1/s_{4\mathrm{\Delta }})`$ in the limit of small squark width. The $`s_{3\mathrm{\Delta }}`$ singular terms may be treated in a similar way, with the added complication that the integration over $`s_{3\mathrm{\Delta }}`$ is hidden in the angular integrations .
The quark emission matrix elements are presented in Appendix E.
For the parameters of the SUGRA model that we adopt, the gluino mass remains greater than the squark mass for all values of $`m_{1/2}`$, and there is never an intermediate squark-to-gluino-plus-quark final state singularity. However, the two chargino masses and all four neutralino masses are always less than the squark mass, and the final-state on-shell squark-to-gaugino-plus-quark singularity comes into play in all cases. In the light gluino model, with $`m_{\stackrel{~}{q}}=450`$ GeV, the gluino mass of 30 GeV is light enough that on-shell intermediate squark decay into a gluino is always active. In this light gluino model, as $`m_{1/2}`$ is varied from 100 to 400 GeV, the masses of the lighter gauginos ($`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_2^0`$, and $`\stackrel{~}{\chi }_1^\pm `$) remain less than the squark mass so that the on-shell intermediate squark decay into a gaugino is active over the whole range of $`m_{1/2}`$ for these light channels. The situation changes for the heavier gauginos ($`\stackrel{~}{\chi }_3^0`$, $`\stackrel{~}{\chi }_4^0`$, and $`\stackrel{~}{\chi }_2^\pm `$). For small $`m_{1/2}`$ their masses are below the squark mass, and the on-shell decay is active. However, above roughly $`m_{1/2}=250`$ GeV, the masses of the heavier gauginos exceed the squark mass and the on-shell possibility closes.
## IV Quantitative Results
In this section we collect our main results on total and differential cross sections for the associated production of gauginos and gluinos at Tevatron and LHC energies.
### A Scaling Functions
We begin with the cross section at the parton level expressed as
$`\widehat{\sigma }_{ij}`$ $`=`$ $`{\displaystyle \frac{\alpha \alpha _S(\mu )}{m^2}}\left\{f_{ij}^B(\eta )+4\pi \alpha _S(\mu )\left[f_{ij}^{V+S}(\eta ,\mu )+f_{ij}^H(\eta ,\mu )\right]\right\}.`$ (74)
It has been integrated over the Mandelstam invariants $`t`$ and $`s_4`$ and depends on the partonic center-of-mass energy $`s`$ through the scaling variable
$`\eta `$ $`=`$ $`{\displaystyle \frac{s}{4m^2}}1,`$ (75)
where $`m`$ is the average mass of the produced sparticles,
$`m`$ $`=`$ $`{\displaystyle \frac{m_1+m_2}{2}}.`$ (76)
It also depends on the produced masses $`m_1`$ and $`m_2`$ and on the squark mass $`m_{\stackrel{~}{q}}`$ (through the internal squark propagator). The common renormalization and factorization scale is denoted by $`\mu `$. The partonic initial state is labeled $`i,j=g,q,\overline{q}`$.
Equation (74) defines the dimensionless scaling functions $`f_{ij}`$, studied in Ref. . These functions are independent of the coupling constants $`\alpha `$ and $`\alpha _S`$, of parton densities, and of the collider type and energy. They permit precise checks of individual contributions and of the threshold, resonance, and high energy behaviors of the production process.
The Born $`f_{ij}^B`$ and the summed virtual and soft scaling functions $`f_{ij}^{V+S}`$ receive contributions only from $`q\overline{q}`$ initial states, where $`q=u\mathrm{or}d`$, with the possible emission of a soft and/or collinear gluon. The hard scaling function $`f_{ij}^H`$ has contributions from $`qg`$ initial states when an additional quark or antiquark is emitted together with the gluino and the gaugino. We eliminate the explicit dependence of the soft contributions on the technical cut-off $`\delta =\mathrm{\Delta }/m^2`$ by subtracting the $`\mathrm{log}^{(1,2)}\delta `$ terms. These terms are then added to the hard contribution such that this contribution is also independent of $`\delta `$. In Sec. III we show that our results are independent of $`\delta `$ at least in the range $`\delta [10^5;10^3]`$, and we use the value $`\delta =10^4`$ in the following.
The scaling functions for the production of a $`\stackrel{~}{g}`$ and a $`\stackrel{~}{\chi }_2^0`$ are presented in Fig. 9 and those for a $`\stackrel{~}{g}`$ and a $`\stackrel{~}{\chi }_1^\pm `$ in Fig. 10. Here we set the scale $`\mu `$ equal to the average particle mass $`m`$. The masses are those of our default SUGRA scenario. As discussed in Sec. III B 2, the emission of an additional quark or antiquark can lead to intermediate on-shell squarks and therefore to a singular squark propagator in Feynman diagrams. After the LO two-body $`q+g\stackrel{~}{g}+\stackrel{~}{q}`$ contribution is removed, the remaining integrable singularities can be identified as spikes in the $`gu`$ and $`gd`$ scaling functions in the two figures. They yield finite contributions after integration over the momentum fractions $`x_a`$ and $`x_b`$ of the incoming partons or, equivalently, over the partonic center-of-mass energy $`s=x_ax_bS`$.
Evident from Figs. 9 and 10 is that next-to-leading order contributions do not alter either the threshold or high energy asymptotic behaviors in $`\eta `$, unlike, e.g., the situation for pair production of heavy quarks . The combined virtual and soft scaling functions $`f_{ij}^{V+S}`$ contribute negatively but are small in magnitude when compared with the hard scaling functions $`f_{ij}^H`$. The figures show that one should expect only modest enhancements in predicted rates when the $`u\overline{u}`$, $`d\overline{d}`$, and $`u\overline{d}`$ production channels are dominant, as is true at Tevatron energies, for which the quark/antiquark parton luminosity is large and the range in $`\eta `$ is limited ( $`\eta <15`$ for these two channels in our default scenario). The contribution of the $`qg`$ channel can become important if phase space is open to large values of $`\eta `$. At LHC energies, $`\eta `$ extends to nearly 800. This range in $`\eta `$, along with the large $`qg`$ luminosity, suggests that the $`qg`$ channel will supply significant enhancements in the predicted rates at LHC energies, as is demonstrated below.
Scaling functions for the $`\stackrel{~}{g}\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{g}\stackrel{~}{\chi }_4^0`$, and $`\stackrel{~}{g}\stackrel{~}{\chi }_2^\pm `$ channels show behavior similar to that seen in Figs. 9 and 10, with the notable exception that the positive excursion at large $`\eta `$ in the $`qg`$ channels is relatively more prominent for $`\stackrel{~}{g}\stackrel{~}{\chi }_4^0`$ and $`\stackrel{~}{g}\stackrel{~}{\chi }_2^\pm `$ than for $`\stackrel{~}{g}\stackrel{~}{\chi }_2^0`$ and $`\stackrel{~}{g}\stackrel{~}{\chi }_1^\pm `$. The $`\stackrel{~}{g}\stackrel{~}{\chi }_3^0`$ channel is distinguished primarily by the fact that the peak in the $`\eta `$ distribution at the Born level occurs near $`\eta =0.4`$ whereas the peaks occur at larger $`\eta `$, in the range of $`\eta =`$ 1 to 2, in all other cases. There is also a noticeable difference in threshold behaviors of the Born and NLO hard-gluon emission contributions for this channel. We relate these differences to the fact that only the $`\stackrel{~}{g}\stackrel{~}{\chi }_3^0`$ channel exhibits positive interference at the Born level between the t- and u-channel contributions, c.f., Eq. (5). Because $`m_{\stackrel{~}{\chi }_3^0}`$ is negative in the SUGRA model, the term proportional to $`X_{tu}`$ in Eq. (5) is positive for $`\stackrel{~}{g}\stackrel{~}{\chi }_3^0`$ production but negative in all other cases.
### B Hadronic Total Cross Sections
The hadronic total cross section is obtained from the partonic cross section through
$`\sigma ^{h_1h_2}(S,\mu )`$ $`=`$ $`{\displaystyle \underset{i,j=g,q,\overline{q}}{}}{\displaystyle _\tau ^1}dx_a{\displaystyle _{\tau /x_a}^1}dx_bf_i^{h_1}(x_a,\mu )f_j^{h_2}(x_b,\mu )\widehat{\sigma }_{ij}(x_ax_bS,\mu ),`$ (77)
where
$`\tau `$ $`=`$ $`{\displaystyle \frac{4m^2}{S}},`$ (79)
and $`\sqrt{S}`$ is the hadronic center-of-mass energy (2 TeV for Run II at the Fermilab $`p\overline{p}`$ collider Tevatron and 14 TeV at the CERN $`pp`$ collider LHC). Our NLO predictions are calculated in the $`\overline{\mathrm{MS}}`$ scheme with the CTEQ5M parametrization for the parton densities $`f(x,\mu )`$ in the proton and antiproton and a two-loop approximation for the strong coupling constant $`\alpha _S`$ with $`\mathrm{\Lambda }^{(5)}=226`$ MeV. This value of $`\mathrm{\Lambda }^{(5)}`$ is used also in the renormalization group evolution equations for our SUSY scenarios. In obtaining LO cross sections, we use the CTEQ5L LO parton densities and the one-loop approximation for $`\alpha _S`$, with $`\mathrm{\Lambda }^{(5)}=146`$ MeV.
#### 1 SUGRA model results
For the SUGRA scenario, we present the total hadronic cross sections for the associated production of gluinos and gauginos at Run II of the Tevatron in Fig. 11 and for the LHC in Fig. 12. We vary the SUGRA parameter $`m_{1/2}`$ from 100 to 400 GeV and keep the other SUGRA parameters fixed at the values described in Sec. II A. The squark mass runs from 250 GeV to 890 GeV in this region. The cross sections are presented as a function of the physical gluino mass $`m_{\stackrel{~}{g}}`$. The corresponding gaugino mass ranges are 31 to 163 GeV for $`\stackrel{~}{\chi }_1^0`$, 62 to 317 GeV for $`\stackrel{~}{\chi }_2^0`$ and $`\stackrel{~}{\chi }_1^\pm `$, 211 to 666 GeV for $`\stackrel{~}{\chi }_3^0`$, and 240 to 679 GeV for $`\stackrel{~}{\chi }_4^0`$ and $`\stackrel{~}{\chi }_2^\pm `$. The chargino cross sections are summed over positive and negative charges. The renormalization and factorization scale $`\mu `$ is set equal to the average particle mass $`m`$. We truncate Fig. 11 at a cross section of $`10^5\mathrm{pb}`$ since the anticipated integrated luminosity at Run II is at most 30 $`\mathrm{fb}^1`$. For the convenience of the reader, we provide numerical values of the cross sections in Table I for a few selected points.
For small $`m_{\stackrel{~}{g}}`$ one might expect the largest cross section for the lightest gaugino, $`\stackrel{~}{\chi }_1^0`$. However, its coupling is dominantly of type $`\stackrel{~}{B}`$ and therefore smaller than the $`\stackrel{~}{W}_3`$-type coupling of $`\stackrel{~}{\chi }_2^0`$ which, in turn, has a larger cross section despite its larger mass. The heavier gauginos $`\stackrel{~}{\chi }_{3,4}^0`$ and $`\stackrel{~}{\chi }_2^\pm `$ are dominantly Higgsino and are therefore suppressed by several orders of magnitude with respect to the lighter gauginos because of the light quark Yukawa couplings. At the LHC, the $`\stackrel{~}{\chi }_1^\pm `$ cross section is dominant. At small values of $`m_{\stackrel{~}{g}}`$, the LHC cross sections are a factor of about 30 greater than at the Tevatron, and at large $`m_{\stackrel{~}{g}}`$, the factor is about $`10^4`$.
Comparing the NLO predictions in Figs. 11 and 12 (solid curves) with the LO predictions (dashed curves) we observe that the NLO corrections are all positive and substantially larger at the LHC than at the Tevatron. At the Tevatron, some of the NLO predictions fall below the LO predictions at large mass, a point to which we return below. The NLO enhancements are more evident in the ratio of the NLO over the LO cross section ($`K`$ factors) shown in Figs. 13 and 14. The $`K`$ factors are computed at the scale $`\mu =m`$. At the Tevatron, the NLO corrections amount to at most a 10% increase in cross section. At the LHC, they appear generally in the range of 20 to 40% but can amount to a factor of 2 for $`\stackrel{~}{\chi }_4^0`$ and $`\stackrel{~}{\chi }_2^\pm `$.
The very modest size of the NLO enhancement at the Tevatron is somewhat expected from the behavior of the scaling functions, but it is also attributable partly to differences in the NLO and LO parton densities. Recalculating the $`K`$ factors with the CTEQ4 parametrization , we find increases in $`K`$ by as much as 0.1 at the Tevatron energy. The change from CTEQ4 to CTEQ5 is interesting. The $`u`$ quark density at NLO decreased by 1 to 5% and the $`d`$ quark density at NLO increased by up to 10% over the range 0.1 $`<x<`$ 0.6. On the other hand, the LO $`u`$ quark density increased by about 1% and the LO $`d`$ quark density by up to 20%. These purely parton density effects result in a net increase in the LO cross sections and decrease in the NLO cross sections, a drop in the calculated $`K`$ factors from CTEQ4 to CTEQ5.
The contribution from the $`qg`$ initial state at the energy of the Tevatron is insignificant (less than $`10^3`$ of the total) for all six gaugino channels, but it is considerable at the energy of the LHC. In Fig. 15 we show the fraction of the NLO cross section at the LHC attributed to the $`qg`$ initial state. At the energy of the LHC, the contribution from the $`q\overline{q}`$ channel is less dominant, and the $`qg`$ contribution becomes significant owing to the large gluon density. A similar effect is seen in single top quark production for comparable values of the average produced mass $`m`$.
The large $`K`$ factor for $`\stackrel{~}{\chi }_4^0`$ and $`\stackrel{~}{\chi }_2^\pm `$ production at the LHC is associated with a large contribution from the $`gq`$ channel at this energy. For these two channels, there is strong interference at LO between the t- and u-channel exchange diagrams that does not occur in the NLO quark emission graphs.
The dependence of the predicted cross sections on the renormalization and factorization scale is reduced considerably at next-to-leading order. As an example, in Fig. 16 we show the scale dependences for the $`\stackrel{~}{\chi }_2^0`$ channel at the Tevatron. All six channels show similar behavior at the Tevatron. Cross sections vary by $`\pm 23\%`$ at LO but only by $`\pm 8\%`$ at NLO as the scale ratio $`\mu /m`$ is varied over the range 0.5 to 2.0, a substantial improvement in reliability. Here $`m`$ is the average mass for the default scenario. At the Tevatron, the NLO and LO cross sections intersect at scale ratio near unity. In Fig. 17 we show the $`\mu `$ dependences for the $`\stackrel{~}{\chi }_2^0`$ channel at the LHC. Again, all six channels show similar behavior. Cross sections vary by $`\pm 12\%`$ at LO but only by $`\pm 4.5\%`$ at NLO in the region $`0.5<\mu /m<2.0`$. However, unlike the situation at the Tevatron, owing to the important contribution of the $`qg`$ channel at the energy of the LHC, the NLO and LO cross sections do not intersect at scale ratio near unity. Instead, if they intersect at all, the crossing point is at a very low scale.
Uncertainties in the cross section from parton density variation may be estimated roughly if we compare NLO results obtained with CTEQ4M and CTEQ5M. For $`\stackrel{~}{\chi }_2^0`$ production at the Tevatron in our default SUGRA scenario, we compute cross sections of 0.0244 pb (CTEQ4M) and 0.0219 pb (CTEQ5M), a $`12\%`$ difference. At the LHC, the cross sections are 1.138 pb (CTEQ4M), and 1.095 pb (CTEQ5M), a $`4\%`$ difference.
#### 2 Light gluino model
For the light gluino model, we present the total hadronic cross sections for the associated production of gluinos and gauginos at Run II of the Tevatron in Fig. 18 and for the LHC in Fig. 19. As mentioned in Sec. II A, we fix $`m_{\stackrel{~}{g}}=`$ 30 GeV and $`m_{\stackrel{~}{q}}=`$ 450 GeV. We display the cross sections as functions of the mass of the common GUT-scale fermion mass $`m_{1/2}`$. As we vary $`m_{1/2}`$ from 100 to 400 GeV, the gaugino masses range from 31 to 163 GeV for $`\stackrel{~}{\chi }_1^0`$, 62 to 317 GeV for $`\stackrel{~}{\chi }_2^0`$ and $`\stackrel{~}{\chi }_1^\pm `$, 211 to 666 GeV for $`\stackrel{~}{\chi }_3^0`$, and 240 to 679 GeV for $`\stackrel{~}{\chi }_4^0`$ and $`\stackrel{~}{\chi }_2^\pm `$. It is worth noting that the coupling strengths also vary with $`m_{1/2}`$. The chargino cross sections are summed over positive and negative charges. The renormalization and factorization scale $`\mu `$ is set equal to the average of the masses of the gaugino and gluino in each of the channels.
Evident in Figs. 18 and 19 is that the cross sections do not depend strongly on the gaugino masses. The values of $`m_{1/2}`$ in these figures extend below the value $`150`$ GeV believed excluded since LEP data set a lower bound on the mass of $`\stackrel{~}{\chi }_1^\pm `$ of about 100 GeV. Nevertheless, even above $`m_{1/2}=150`$ GeV, the Tevatron cross sections for the three lighter gluino channels are predicted to be in the range of 0.1 to 0.5 $`\mathrm{pb}`$. The cross sections would be increased if $`m_{\stackrel{~}{q}}`$ were reduced from the value 450 GeV that we use.
The relatively large cross sections suggest that associated production is a good channel for discovery of a light gluino at the Tevatron, for closing the window on this possibility, and/or for setting limits on light gaugino masses. We remark in this connection that the usual searches for a light gluino LSP begin with the assumption of pair production of gluinos. In this situation, the dominant background is QCD pair production of hadronic jets. Hard cuts on transverse momentum must be made to reduce this background to tolerable levels. The cuts, in turn, mitigate against gluinos of modest mass. By contrast, if light gluinos are produced in association with gauginos, one can search for light gluino monojets accompanied by leptons and/or missing transverse energy from gaugino decays.
The $`K`$ factors for the light gluino case are shown in Figs. 20 and 21. At the Tevatron, the NLO corrections amount to a $`30\%`$ to $`40\%`$ increase in cross section. At the LHC, they are generally in the range of a factor of 2 to 3.5. The large $`K`$ factors owe their origins to the important role of the $`gq`$ channel. The gluon parton density is very large at small values of $`x`$. Contributions from the $`gq`$ production channel are more intense in the light gluino case than in the SUGRA case where average produced masses and, thus, typical values of $`x`$ are larger. At the energy of the Tevatron, the $`gq`$ channel accounts for more than 20% of the NLO cross section in the $`\stackrel{~}{\chi }_1^0`$ and $`\stackrel{~}{\chi }_1^\pm `$ channels at small values of $`m_{1/2}`$ and more than 10% at large $`m_{1/2}`$. At the energy of the LHC, the fraction exceeds 60% in the $`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_2^0`$, and $`\stackrel{~}{\chi }_1^\pm `$ channels for the entire range of $`m_{1/2}`$. It hovers above 50% in the $`\stackrel{~}{\chi }_3^0`$, $`\stackrel{~}{\chi }_4^0`$, and $`\stackrel{~}{\chi }_2^\pm `$ channels at small values of $`m_{1/2}`$ and near 20% at large $`m_{1/2}`$.
Motivated by the curious behavior of the $`K`$ factors for the three heavier gaugino channels at the LHC energy, we examined the change in the scaling functions for the $`\stackrel{~}{\chi }_4^0`$ case as $`m_{1/2}`$ is varied. In the $`q\overline{q}`$ channel, the net (virtual plus soft plus hard) NLO contribution is positive for all $`\eta `$, and, relative to the Born contribution, its magnitude grows gradually with $`m_{1/2}`$. The $`q\overline{q}`$ channel accounts for a slightly increasing component of $`K`$, hovering about 1.5. On the other hand, the contribution from the $`gq`$ channel changes markedly as $`m_{1/2}`$ increases. Below about $`m_{1/2}=200`$ GeV, its scaling function is large, with significant support below $`\eta =1`$, where the gluon parton density is large. As $`m_{1/2}`$ increases above 200 GeV, the $`gq`$ scaling function decreases in magnitude. The $`gq`$ channel supplies a component of $`K`$ that increases slightly from about 2 to 2.5 as $`m_{1/2}`$ increases to 200 GeV and then falls gradually to below 0.5 at $`m_{1/2}=400`$ GeV.
We attribute the sharp decrease of the $`K`$ factor for the three heavier gaugino channnels to the role of on-shell intermediate squark decay into a gaugino, discussed in Sec. III B 2. In the light gluino model, as $`m_{1/2}`$ is varied from 100 to 400 GeV, the masses of the lighter gauginos ($`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_2^0`$, and $`\stackrel{~}{\chi }_1^\pm `$) remain less than the squark mass so that the on-shell intermediate squark decay into a gaugino is active over the whole range of $`m_{1/2}`$ for these light channels. The situation changes for the heavier gauginos ($`\stackrel{~}{\chi }_3^0`$, $`\stackrel{~}{\chi }_4^0`$, and $`\stackrel{~}{\chi }_2^\pm `$). For small $`m_{1/2}`$ their masses are below the squark mass, and the on-shell decay is active. However, above roughly $`m_{1/2}=250`$ GeV, the masses of the heavier gauginos exceed the squark mass and the on-shell possibility closes.
### C Differential Cross Sections
In Figs. 22 and 23, we display the differential cross sections in the rapidity $`y`$ and in the transverse momentum $`p_T`$ of $`\stackrel{~}{\chi }_1^\pm `$ at the energy of the Tevatron collider. Here the $`\stackrel{~}{\chi }_1^\pm `$ and $`\stackrel{~}{g}`$ masses are set at their default values in the SUGRA scenario, 101 and 410 GeV, respectively. The rapidity distribution is obtained after integration over all $`p_T`$, and the $`p_T`$ distribution after integration over all $`y`$. More restrictive selections could be made, but until experimental conditions are better known, any such restrictions (cuts) would be unmotivated. The NLO (solid) curve in Fig. 22 shows a modest enhancement in the rapidity distribution in the central region with respect to the LO (dashed) curve, but the shape of the distribution is unchanged qualitatively. The $`p_T`$ distributions in Fig. 23 show that the NLO contribution tends to shift the distribution to somewhat smaller values of $`p_T`$. Since the contribution of the $`gq`$ initial state is very small at the energy of the Tevatron, the shift in the $`p_T`$ distribution is associated with next-to-leading order corrections in the dominant $`q\overline{q}`$ initial state.
The features of the $`y`$ distributions are qualitatively similar for all gauginos except for the expected and systematic narrowing of the $`y`$ distribution with increasing gaugino mass. We show one example representative of the full set. The $`p_T`$ distributions for the different gauginos are also qualitatively similar except that the maximum in the distribution moves to larger $`p_T`$ as the gaugino mass is increased. The location of the peak is specified roughly by $`m/2`$, where $`m`$ is the average mass of the produced sparticles. The one exception is $`\stackrel{~}{\chi }_3^0`$ production. In this case, the peak occurs at a smaller value (about 100 GeV for the default masses), an effect correlated with the fact that the location of the peak in the scaling function occurs at a smaller value of $`\eta `$. Interference effects enhance the cross section at small $`p_T`$ for $`\stackrel{~}{\chi }_3^0`$ production.
We show differential cross sections in $`y`$ and $`p_T`$ for $`\stackrel{~}{\chi }_1^\pm `$ at the energy of the LHC collider in Figs. 24 and 25, and in Fig. 26 we present the $`p_T`$ distribution for $`\stackrel{~}{\chi }_4^0`$. The rapidity and transverse momentum distributions are much broader than at the Tevatron. As at the energy of the Tevatron, the features of the $`y`$ distributions are qualitatively similar for all gauginos except for the expected and systematic narrowing of the $`y`$ distribution with increasing gaugino mass. The NLO contributions enhance the $`y`$ distributions at all $`y`$.
At the energy of the LHC, the $`p_T`$ spectra are qualitatively similar for the relatively light $`\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{\chi }_2^0`$, and $`\stackrel{~}{\chi }_1^\pm `$ states, illustrated by the $`\stackrel{~}{\chi }_1^\pm `$ case in Fig. 25. The enhancement factor $`K`$ is near unity at small $`p_T`$ but becomes sizeable at larger $`p_T`$. For $`\stackrel{~}{\chi }_4^0`$ and $`\stackrel{~}{\chi }_2^\pm `$, the $`p_T`$ spectra are altered significantly by NLO contributions. As illustrated in Fig. 26, the NLO contribution associated with the important $`qg`$ channel fills in the distribution at small $`p_T`$ and softens the overall $`p_T`$ distribution in these two cases. At LHC energies, it is therefore not a good approximation to assume that the enhancement factor $`K`$ is roughly independent of $`p_T`$.
## V Conclusions
In this paper we report a complete next-to-leading order analysis of the associated production of gauginos and gluinos at hadron colliders. If supersymmetry exists at the electroweak scale, the cross section for this process is expected to be observable at the Fermilab Tevatron and/or the CERN LHC. It is enhanced by the large color charge of the gluino and the relatively small mass of the light gauginos in many SUSY models. Associated production represents a chance to study in detail the parameters of the soft SUSY-breaking Lagrangian. The rates are proportional to the phases of the gaugino and gluino masses, and to the mixings in the squark and chargino/neutralino sectors. Thus, in combination with other channels, associated production could allow one to measure some or all of these quantities.
The physical gluino and gaugino masses that we use, as well as the gaugino mixing matrices, are based on four popular SUSY breaking models plus a fifth scenario in which the gluino mass is relatively light. Because the LO cross sections in gauge-mediated, gaugino-mediated, and anomaly-mediated supersymmetry breaking models are not too dissimilar from those of the SUGRA case at Tevatron energies, we focus our NLO work on the SUGRA model and on a model with a light gluino LSP, with $`m_{\stackrel{~}{g}}=`$ 30 GeV.
In the SUGRA model, the largest cross sections at the Fermilab Tevatron energy are those for $`\stackrel{~}{\chi }_2^0`$, enhanced by its $`\stackrel{~}{W}_3`$-like coupling with respect to the $`\stackrel{~}{B}`$-like $`\stackrel{~}{\chi }_1^0`$, and the $`\stackrel{~}{\chi }_1^\pm `$, which is about equal in mass with the $`\stackrel{~}{\chi }_2^0`$. The NLO corrections to associated production are generally positive, but they can be modest in size, ranging in the SUGRA model from a few percent at the energy of the Tevatron to 100% at the energy of the LHC, depending on the sparticle masses. In the light-gluino case, NLO contributions increase the cross section by factors of 1.3 to 1.4 at the energy of the Tevatron and by factors of 2 to 3.5 at the energy of the LHC. The large $`K`$ factors owe their origins to the important role of the $`gq`$ channel that enters first at NLO.
Owing to the NLO enhancements, collider searches for signatures of associated production will generally discover or exclude sparticles with masses larger than one would estimate based on LO production rates alone. More significant from the viewpoint of reliability, the renormalization and factorization scale dependence of the cross sections is reduced by a factor of more than two when NLO contributions are included.
At Run II of the Fermilab Tevatron, for an integrated luminosity of 2 $`\mathrm{fb}^1`$, we expect that 10 or more events could be produced in each of the lighter gaugino channels of the SUGRA model, $`\stackrel{~}{g}\stackrel{~}{\chi }_1^0`$, $`\stackrel{~}{g}\stackrel{~}{\chi }_2^0`$, and $`\stackrel{~}{g}\stackrel{~}{\chi }_1^\pm `$, provided that the gluino mass $`m_{\stackrel{~}{g}}`$ is less than 450 GeV. The cross sections for the three heavier gaugino channels, $`\stackrel{~}{g}\stackrel{~}{\chi }_3^0`$, $`\stackrel{~}{g}\stackrel{~}{\chi }_4^0`$, and $`\stackrel{~}{g}\stackrel{~}{\chi }_2^\pm `$, are smaller by an order of magnitude or more than those of the lighter gaugino channels. In the light gluino model, more than 100 events could be produced in the three lighter gaugino channels provided that the common GUT-scale fermion mass $`m_{1/2}`$ is less than 400 GeV, and as many as 10 events in the three heavier gaugino channels as long as $`m_{1/2}`$ is less than 200 GeV. At the higher energy and luminosity of the LHC, at least a few events should be produced in every channel in the SUGRA model and many more in the light gluino model.
The shapes of the rapidity distributions of the gauginos are not altered appreciably by NLO contributions, but the locations of the maximum cross section in transverse momentum ($`p_T`$) are shifted to smaller values by NLO contributions. At LHC energies where the contribution of the $`qg`$ initial state is important, modifications of the $`p_T`$ spectra can be pronounced.
The relatively large cross sections suggest that associated production is a good channel for discovery of a light gluino at the Tevatron, for closing the window on this possibility, and/or for setting limits on light gaugino masses. The usual searches for a light gluino LSP are based on the assumption that gluinos are produced in pairs. In this situation, the dominant background is QCD production of hadronic jets. Hard cuts on transverse momentum must be made to reduce this background to tolerable levels. The cuts, in turn, mitigate against gluinos of modest mass. By contrast, if light gluinos are produced in association with gauginos, one can search for light gluino monojets accompanied by leptons and/or missing transverse energy from gaugino decays.
## Acknowledgments
Work in the High Energy Physics Division at Argonne National Laboratory is supported by the U.S. Department of Energy, Division of High Energy Physics, under Contract W-31-109-ENG-38. M. Klasen is supported by the Bundesministerium für Bildung und Forschung under Contract 05 HT9GUA 3, by the Deutsche Forschungsgemeinschaft under Contract KL 1266/1-1, and by the European Commission under Contract ERBFMRXCT980194. The authors are grateful for correspondence with W. Beenakker and T. Plehn and conversations with S. Mrenna. T. Tait benefitted from discussions with D. Kaplan, G. Kribs, and C.–P. Yuan.
## A Virtual Loop Contributions
At the one-loop level in SUSY-QCD, virtual corrections contribute to the hadroproduction of supersymmetric particles through the interference of self-energy corrections, vertex corrections, and box diagrams with the tree-level diagrams. For the associated production of gluinos and gauginos one has to calculate the self-energy corrections in Fig. 3, the vertex corrections in Fig. 4, and the box diagrams in Fig. 5.
The self-energy corrections for the external quark and antiquark factorize the complete Born matrix element and are independent of the underlying scattering process. Since the quark and antiquark are treated as massless, only the squark-gluino loop contributes to
$`\overline{|^q|}^2`$ $`=`$ $`i[B_0(0;m_{\stackrel{~}{g}},m_{\stackrel{~}{q}})+B_0^{}(0;m_{\stackrel{~}{g}},m_{\stackrel{~}{q}})(m_{\stackrel{~}{g}}^2m_{\stackrel{~}{q}}^2)]2C_F\widehat{g}_s^2\overline{|^B|}^2,`$ (A1)
where the factor of two accounts for the sum of the quark and antiquark contributions. Functions $`B_0(p^2;m_1,m_2)`$ and $`B_0^{}(p^2;m_1,m_2)`$ stand for the scalar two-point integral and its derivative. They are defined in Appendix B. The quark-gluon loop contains an infrared singularity
$`\overline{|_{\mathrm{IR}}^q|}^2`$ $`=`$ $`C_ϵ{\displaystyle \frac{1}{ϵ}}2C_Fg_s^2\overline{|^B|}^2,`$ (A2)
where $`C_ϵ=(4\pi )^ϵ/(16\pi ^2)e^{ϵ\gamma _E}`$. This infrared singularity is not shown in the result above since it is canceled by an ultraviolet singularity when infrared and ultraviolet singularities are not distinguished in dimensional regularization.
The external gluino self-energy also factorizes the complete Born matrix element and is independent of the underlying scattering process:
$`\overline{|^{\stackrel{~}{g}}|}^2`$ $`=`$ $`i\left[{\displaystyle \frac{A_0(m_{\stackrel{~}{g}})}{m_{\stackrel{~}{g}}^2}}(1ϵ)4B_0^{}(m_{\stackrel{~}{g}}^2;0,m_{\stackrel{~}{g}})m_{\stackrel{~}{g}}^2\right]N_Cg_s^2\overline{|^B|}^2`$ (A3)
$`+`$ $`i\left[{\displaystyle \frac{A_0(m_{\stackrel{~}{q}})}{2m_{\stackrel{~}{g}}^2}}+B_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0){\displaystyle \frac{m_{\stackrel{~}{g}}^2+m_{\stackrel{~}{q}}^2}{2m_{\stackrel{~}{g}}^2}}+B_0^{}(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0)(m_{\stackrel{~}{g}}^2m_{\stackrel{~}{q}}^2)\right]n_f4{\displaystyle \frac{1}{2}}\widehat{g}_s^2\overline{|^B|}^2`$ (A4)
$`+`$ $`i[{\displaystyle \frac{A_0(m_t)}{2m_{\stackrel{~}{g}}^2}}{\displaystyle \frac{A_0(m_{\stackrel{~}{q}})}{2m_{\stackrel{~}{g}}^2}}+B_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},m_t){\displaystyle \frac{m_{\stackrel{~}{g}}^2m_t^2+m_{\stackrel{~}{q}}^2}{2m_{\stackrel{~}{g}}^2}}+B_0^{}(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},m_t)`$ (A6)
$`\times (m_{\stackrel{~}{g}}^2+m_t^2m_{\stackrel{~}{q}}^2)]4{\displaystyle \frac{1}{2}}\widehat{g}_s^2\overline{|^B|}^2.`$
The scalar one-point integral $`A_0(m)`$ is defined in Appendix B. Except for the gluino mass and the $`N_C`$ color factor, the gluon-gluino loop is identical to the heavy-quark self-energy. It contains an infrared singularity in the derivative of the scalar two-point integral:
$`\overline{|_{\mathrm{IR}}^{\stackrel{~}{g}}|}^2`$ $`=`$ $`C_ϵ{\displaystyle \frac{1}{ϵ}}4{\displaystyle \frac{N_C}{2}}g_s^2\overline{|^B|}^2.`$ (A7)
The quark-squark loop with a color factor of 1/2 contributes through two different fermion number flows due to the Majorana nature of gluinos. We take into account $`n_f=5`$ light (s)quark flavors and a heavy top (s)quark. We do not include mixing in the top squark sector and take the top squark mass equal to the light squark masses. We set $`m_t=175`$ GeV.
The gaugino couples only electroweakly to the quarks and squarks and thus does not give rise to strong self-energy corrections. All external particle self-energies have been renormalized on-shell and multiplied by a factor of 1/2 for proper wave function renormalization.
The self-energy correction of the internal squark propagator depends on the off-shell squark four-momentum squared. Therefore, it factorizes only the corresponding $`t`$\- or $`u`$-channel interference piece of the Born matrix element:
$`\overline{|^{\stackrel{~}{q_t}}|}^2`$ $`=`$ $`i\left[B_0(t;m_{\stackrel{~}{g}},0){\displaystyle \frac{tm_{\stackrel{~}{g}}^2}{tm_{\stackrel{~}{q}}^2}}B_0(m_{\stackrel{~}{q}}^2;m_{\stackrel{~}{g}},0){\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2}{tm_{\stackrel{~}{q}}^2}}\right]4C_F\widehat{g}_s^2[\overline{|^t^t|}+\overline{|^t^u|}]`$ (A8)
$``$ $`i\left[B_0(t;m_{\stackrel{~}{q}},0){\displaystyle \frac{t+m_{\stackrel{~}{q}}^2}{tm_{\stackrel{~}{q}}^2}}B_0(m_{\stackrel{~}{q}}^2;m_{\stackrel{~}{q}},0){\displaystyle \frac{2m_{\stackrel{~}{q}}^2}{tm_{\stackrel{~}{q}}^2}}\right]4C_Fg_s^2[\overline{|^t^t|}+\overline{|^t^u|}].`$ (A9)
The ultraviolet divergences cancel between the quark-gluino loop contribution and its supersymmetric counterpart, the squark-gluon loop contribution. Since the gluon tadpole contribution is quadratic in the loop momentum it vanishes in dimensional regularization. The squark tadpole contribution vanishes after renormalization. The $`u`$-channel result can be obtained from the $`t`$-channel result given above through the exchange $`tu`$.
Like the self-energy correction of the squark propagator, the corrections to the quark-squark-gluino and quark-squark-gaugino vertices depend on the four-momentum squared of the squark and factorize only the $`t`$\- or $`u`$-channel interference pieces of the Born matrix element. The quark-squark-gaugino vertex receives corrections through a gluon and a gluino exchange between the initial state quark or antiquark and the squark that are proportional to the gauge and Yukawa coupling, respectively. For the $`t`$-channel we find
$`\overline{|^{q\stackrel{~}{q}\stackrel{~}{\chi }}|}^2`$ $`=`$ $`i\left[B_0(t;m_{\stackrel{~}{g}},0)B_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{q}},0)+C_0(m_{\stackrel{~}{\chi }}^2,0,t;0,m_{\stackrel{~}{q}},m_{\stackrel{~}{g}})(t+m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2m_{\stackrel{~}{\chi }}^2)\right]`$ (A11)
$`\times {\displaystyle \frac{m_{\stackrel{~}{g}}m_{\stackrel{~}{\chi }}}{tm_{\stackrel{~}{\chi }}^2}}4C_F\widehat{g}_s^2{\displaystyle \frac{X_t^{q^{}\stackrel{~}{q}\stackrel{~}{\chi }}}{X_t^{q\stackrel{~}{q}^{}\stackrel{~}{\chi }}}}[\overline{|^t^t|}+\overline{|^t^u|}]`$
$`+`$ $`i[B_0(t;m_{\stackrel{~}{q}},0){\displaystyle \frac{t+m_{\stackrel{~}{\chi }}^2}{2(tm_{\stackrel{~}{\chi }}^2)}}B_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{q}},0){\displaystyle \frac{m_{\stackrel{~}{\chi }}^2}{tm_{\stackrel{~}{\chi }}^2}}+C_0(m_{\stackrel{~}{\chi }}^2,0,t;m_{\stackrel{~}{q}},0,0)`$ (A13)
$`\times (m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{q}}^2)]4C_Fg_s^2[\overline{|^t^t|}+\overline{|^t^u|}].`$
Function $`C_0(p_1^2,p_2^2,(p_1+p_2)^2;m_1,m_2,m_3)`$ stands for the scalar three-point integral defined in Appendix B. The ratio of quark-squark-gaugino couplings $`X_t^{q^{}\stackrel{~}{q}\stackrel{~}{\chi }}/X_t^{q\stackrel{~}{q}^{}\stackrel{~}{\chi }}`$ accounts for reversed flavor flow in the vertex correction with respect to the underlying Born matrix element in the case of the exchange of an additional gluino. For neutralinos with real couplings it reduces to unity, whereas for charginos with real couplings it is given by a ratio of chargino mixing matrix elements. The infrared singularities
$`\overline{|_{\mathrm{IR}}^{q\stackrel{~}{q}\stackrel{~}{\chi }}|}^2`$ $`=`$ $`C_ϵ\left[{\displaystyle \frac{1}{ϵ}}+{\displaystyle \frac{1}{ϵ}}\mathrm{log}\left({\displaystyle \frac{m_{\stackrel{~}{q}}^2t}{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}}\right){\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}{tm_{\stackrel{~}{\chi }}^2}}\right]4C_Fg_s^2[\overline{|^t^t|}+\overline{|^t^u|}]`$ (A14)
arise from the gluon exchange correction. In dimensional regularization, the gauge bosons have $`(n2)`$ degrees of freedom whereas their supersymmetric counterparts, the gauginos, have two. This difference leads to a mismatch between the quark-quark-gauge boson gauge couplings and the quark-squark-gaugino Yukawa couplings through finite next-to-leading order terms. The (super-)symmetry between the gauge and Yukawa couplings can be restored through a finite renormalization contribution
$`\overline{|_{\mathrm{finite}}^{q\stackrel{~}{q}\stackrel{~}{\chi }}|}^2`$ $`=`$ $`C_ϵC_Fg_s^2[\overline{|^t^t|}+\overline{|^t^u|}],`$ (A15)
that can be found by comparing the quark-squark-gaugino vertex correction given above with the corresponding quark-quark-gauge boson vertex correction in exact supersymmetry. All $`u`$-channel results can again be obtained through the exchange $`tu`$.
The quark-squark-gluino vertex correction can be obtained from the quark-squark-gaugino vertex correction if $`m_{\stackrel{~}{\chi }}`$ is replaced with $`m_{\stackrel{~}{g}}`$ and the color factor $`C_F`$ with $`C_FN_C/2`$. There are, however, two additional contributions with a color factor of $`N_C/2`$, due to the non-Abelian gauge coupling of the gluino to the gluon, when a gluon is exchanged between the final state gluino and the initial state quark or antiquark and the squark. For the total quark-squark-gluino vertex correction in the $`t`$-channel we find
$`\overline{|^{q\stackrel{~}{q}\stackrel{~}{g}}|}^2`$ $`=`$ $`\overline{|^{q\stackrel{~}{q}\stackrel{~}{\chi }}|}^2(m_{\stackrel{~}{\chi }}m_{\stackrel{~}{g}},C_FC_FN_C/2)`$ (A16)
$`+`$ $`i[B_0(t;m_{\stackrel{~}{g}},0)({\displaystyle \frac{m_{\stackrel{~}{g}}^2}{tm_{\stackrel{~}{g}}^2}}+ϵ)B_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{g}},0){\displaystyle \frac{t}{tm_{\stackrel{~}{g}}^2}}+C_0(m_{\stackrel{~}{g}}^2,0,t;m_{\stackrel{~}{g}},0,0)(tm_{\stackrel{~}{g}}^2)`$ (A17)
$`+`$ $`B_0(t;0,m_{\stackrel{~}{q}}){\displaystyle \frac{t+m_{\stackrel{~}{g}}^2}{2(tm_{\stackrel{~}{g}}^2)}}B_0(m_{\stackrel{~}{g}}^2;0,m_{\stackrel{~}{g}}){\displaystyle \frac{t}{tm_{\stackrel{~}{g}}^2}}+C_0(m_{\stackrel{~}{g}}^2,0,t;0,m_{\stackrel{~}{g}},m_{\stackrel{~}{q}}){\displaystyle \frac{m_{\stackrel{~}{g}}^4tm_{\stackrel{~}{q}}^2}{tm_{\stackrel{~}{g}}^2}}]`$ (A19)
$`\times 4{\displaystyle \frac{N_C}{2}}g_s^2[\overline{|^t^t|}+\overline{|^t^u|}].`$
It contains the following infrared singularities:
$`\overline{|_{\mathrm{IR}}^{q\stackrel{~}{q}\stackrel{~}{g}}|}^2`$ $`=`$ $`\overline{|_{\mathrm{IR}}^{q\stackrel{~}{q}\stackrel{~}{\chi }}|}^2(m_{\stackrel{~}{\chi }}m_{\stackrel{~}{g}},C_FC_FN_C/2)`$ (A20)
$``$ $`C_ϵ\left[{\displaystyle \frac{1}{2ϵ^2}}+{\displaystyle \frac{1}{ϵ}}+{\displaystyle \frac{1}{2ϵ}}\mathrm{log}\left({\displaystyle \frac{Q^2}{m_{\stackrel{~}{g}}^2}}\right){\displaystyle \frac{1}{ϵ}}\mathrm{log}\left({\displaystyle \frac{m_{\stackrel{~}{g}}^2t}{m_{\stackrel{~}{g}}^2}}\right)\right]`$ (A22)
$`\times 4{\displaystyle \frac{N_C}{2}}g_s^2[\overline{|^t^t|}+\overline{|^t^u|}].`$
The finite renormalization contribution for the gluino vertex correction is
$`\overline{|_{\mathrm{finite}}^{q\stackrel{~}{q}\stackrel{~}{g}}|}^2`$ $`=`$ $`C_ϵ\left({\displaystyle \frac{4}{3}}N_CC_F\right)g_s^2[\overline{|^t^t|}+\overline{|^t^u|}].`$ (A23)
For the $`u`$-channel contribution, $`t`$ and $`u`$ have to be exchanged as before.
Turning to the box diagrams that contribute to the associated production of gluinos and gauginos, we notice that they depend naturallly on the four-momentum squared of the exchanged squark. Therefore, we do not expect the full Born matrix element to factorize. In addition, the traces of the Dirac matrices project out terms that depend on the final state masses separately from those that do not, so that the squared $`t`$-channel and $`u`$-channel diagrams and the interference term can only be factorized individually.
For the first box diagram in Fig. 5 we find
$`\overline{|^{\mathrm{Box1}}|}^2`$ $`=`$ $`iB_0(t;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)4\overline{|^t^u|}/s`$ (A24)
$`+`$ $`iB_0(t;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A26)
$`(2m_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)2/(tm_{\stackrel{~}{\chi }}^2))`$
$`+`$ $`iB_0(s;0,0)(tm_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A28)
$`(2s+4(tm_{\stackrel{~}{g}}^2)2m_{\stackrel{~}{\chi }}^2+2m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iB_0(s;0,0)(tm_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A32)
$`(8sm_{\stackrel{~}{g}}^4/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)+2sm_{\stackrel{~}{\chi }}^2/(tm_{\stackrel{~}{\chi }}^2)+8sm_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{\chi }}^2)`$
$`6sm_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{g}}^2)+2m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{\chi }}^2)2m_{\stackrel{~}{\chi }}^4/(tm_{\stackrel{~}{\chi }}^2)+2m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{g}}^2)`$
$`2m_{\stackrel{~}{g}}^4/(tm_{\stackrel{~}{g}}^2))`$
$`+`$ $`iB_0(s;0,0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A34)
$`(2m_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)+2/(tm_{\stackrel{~}{\chi }}^2))`$
$`+`$ $`iB_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A36)
$`(2(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2/s+2(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2/s2s2(tm_{\stackrel{~}{g}}^2)+2m_{\stackrel{~}{\chi }}^2+2m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iB_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A39)
$`(4sm_{\stackrel{~}{g}}^4/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)2sm_{\stackrel{~}{\chi }}^2/(tm_{\stackrel{~}{\chi }}^2)4sm_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{\chi }}^2)`$
$`+4sm_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{g}}^2)2m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{\chi }}^2)+2m_{\stackrel{~}{\chi }}^4/(tm_{\stackrel{~}{\chi }}^2))`$
$`+`$ $`iB_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A41)
$`(2(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2/s2(tm_{\stackrel{~}{g}}^2)/sm_{\stackrel{~}{g}}^22(tm_{\stackrel{~}{g}}^2)4m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iB_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A44)
$`(4s/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^44s/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2+2s/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2`$
$`2/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2+2/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4)`$
$``$ $`iB_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)4\overline{|^t^u|}/s`$ (A45)
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{\chi }}^2,s;0,m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A48)
$`(2m_{\stackrel{~}{q}}^2s+4m_{\stackrel{~}{q}}^2(tm_{\stackrel{~}{g}}^2)2m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2+2m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2+2s(tm_{\stackrel{~}{g}}^2)+2sm_{\stackrel{~}{g}}^2`$
$`2(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^22(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+2m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^22m_{\stackrel{~}{g}}^4)`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{\chi }}^2,s;0,m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A54)
$`(8m_{\stackrel{~}{q}}^2s/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4+2m_{\stackrel{~}{q}}^2s/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2+8m_{\stackrel{~}{q}}^2s/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2`$
$`6m_{\stackrel{~}{q}}^2s/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+2m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^22m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^4`$
$`+2m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^22m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4+8s/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^6`$
$`+6s/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2+8s/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^42s/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^28s/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4`$
$`+2/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^42/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^4m_{\stackrel{~}{g}}^22/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^4+2/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^4m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{\chi }}^2,s;0,m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)8\overline{|^t^u|}`$ (A55)
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{\chi }}^2,s;0,m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A59)
$`(m_{\stackrel{~}{q}}^2s/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)+2m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{\chi }}^2)`$
$`m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{g}}^2)+s/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2s/(tm_{\stackrel{~}{g}}^2)`$
$`+2/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4+3/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2+m_{\stackrel{~}{\chi }}^2/(tm_{\stackrel{~}{g}}^2)2/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2)`$
$``$ $`iC_0(m_{\stackrel{~}{g}}^2,0,t;m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)4\overline{|^t^u|}`$ (A60)
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,0,t;m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A62)
$`(1+m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2/(tm_{\stackrel{~}{g}}^2))`$
$``$ $`iC_0(m_{\stackrel{~}{\chi }}^2,0,t;m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)4\overline{|^t^u|}`$ (A63)
$`+`$ $`iC_0(m_{\stackrel{~}{\chi }}^2,0,t;m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A65)
$`(1+m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2/(tm_{\stackrel{~}{\chi }}^2))`$
$`+`$ $`iC_0(0,0,s;0,0,0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A67)
$`(m_{\stackrel{~}{q}}^2s/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)s/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+s/(tm_{\stackrel{~}{g}}^2))`$
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2,0,0;0,m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)4\overline{|^t^u|}s`$ (A68)
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2,0,0;0,m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)4\overline{|^t^t|}s`$ (A69)
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2,0,0;0,m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)^2g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A71)
$`(s/(tm_{\stackrel{~}{\chi }}^2)s/(tm_{\stackrel{~}{g}}^2))`$
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2,0,0;0,m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)^3g_s^2(C_FN_C/2)2\overline{|^t^t|}`$ (A73)
$`(s/(tm_{\stackrel{~}{\chi }}^2)/(tm_{\stackrel{~}{g}}^2)).`$
Function $`K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)=s^22s(m_{\stackrel{~}{\chi }}^2+m_{\stackrel{~}{g}}^2)+(m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2)^2`$ is the triangle (Källén) function of the partonic center-of-mass energy and the masses of the produced particles. The corresponding $`u`$-channel contribution is obtained by exchanging $`t`$ and $`u`$.
For the second box diagram in Fig. 5 we find
$`\overline{|^{\mathrm{Box2}}|}^2`$ $`=`$ $`iB_0(t;m_{\stackrel{~}{g}},0)(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)4\overline{|^t^u|}/s`$ (A74)
$`+`$ $`iB_0(t;m_{\stackrel{~}{g}},0)(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A77)
$`(2/(tm_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2+2/(tm_{\stackrel{~}{g}}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2`$
$`+2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+2/(um_{\stackrel{~}{\chi }}^2))`$
$`+`$ $`iB_0(s;m_{\stackrel{~}{q}},m_{\stackrel{~}{q}})(um_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A79)
$`(2s4(um_{\stackrel{~}{g}}^2)+2m_{\stackrel{~}{\chi }}^22m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iB_0(s;m_{\stackrel{~}{q}},m_{\stackrel{~}{q}})(um_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A83)
$`(8s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^42s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^28s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2`$
$`+6s/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^22/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2+2/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^42/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2`$
$`+2/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4)`$
$`+`$ $`iB_0(s;m_{\stackrel{~}{q}},m_{\stackrel{~}{q}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A85)
$`(2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^22/(um_{\stackrel{~}{\chi }}^2))`$
$`+`$ $`iB_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0)(um_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A87)
$`(2(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2/s2(um_{\stackrel{~}{g}}^2)/sm_{\stackrel{~}{g}}^2+2s+2(um_{\stackrel{~}{g}}^2)2m_{\stackrel{~}{\chi }}^22m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iB_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0)(um_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A90)
$`(4s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4+2s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2+4s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2`$
$`4s/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+2/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^22/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^4)`$
$``$ $`iB_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0)(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)4\overline{|^t^u|}/s`$ (A91)
$`+`$ $`iB_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0)(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A93)
$`(2/(tm_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2)`$
$`+`$ $`iB_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{q}},0)(um_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A95)
$`(2(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2/s+2(um_{\stackrel{~}{g}}^2)/sm_{\stackrel{~}{g}}^2+2(um_{\stackrel{~}{g}}^2)+4m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iB_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{q}},0)(um_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A98)
$`(4s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4+4s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^22s/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2`$
$`+2/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^22/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4)`$
$`+`$ $`iB_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{q}},0)(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A100)
$`(2/(tm_{\stackrel{~}{g}}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{\chi }}^2,s;m_{\stackrel{~}{q}},0,m_{\stackrel{~}{q}})(um_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A103)
$`(2m_{\stackrel{~}{q}}^2s+4m_{\stackrel{~}{q}}^2(um_{\stackrel{~}{g}}^2)2m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2+2m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^22s(um_{\stackrel{~}{g}}^2)2sm_{\stackrel{~}{g}}^2`$
$`+2(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2+2(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^22m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2+2m_{\stackrel{~}{g}}^4)`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{\chi }}^2,s;m_{\stackrel{~}{q}},0,m_{\stackrel{~}{q}})(um_{\stackrel{~}{q}}^2)/K(s,m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A109)
$`(8m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4+2m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2+8m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2`$
$`6m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+2m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^22m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^4`$
$`+2m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^22m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^48s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^6`$
$`6s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^28s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^4+2s/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2+8s/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4`$
$`2/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^4+2/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^4m_{\stackrel{~}{g}}^2+2/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^42/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^4m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{\chi }}^2,s;m_{\stackrel{~}{q}},0,m_{\stackrel{~}{q}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A111)
$`(12m_{\stackrel{~}{q}}^2/s+2(um_{\stackrel{~}{g}}^2)/sm_{\stackrel{~}{\chi }}^2/s+3m_{\stackrel{~}{g}}^2/s)`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{\chi }}^2,s;m_{\stackrel{~}{q}},0,m_{\stackrel{~}{q}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A116)
$`(2m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)+4m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+2m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{\chi }}^2)`$
$`2m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{g}}^2)s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2s/(um_{\stackrel{~}{\chi }}^2)+s/(um_{\stackrel{~}{g}}^2)`$
$`+s^2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)4/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^43/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2`$
$`+3/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,0,t;0,m_{\stackrel{~}{q}},m_{\stackrel{~}{g}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A118)
$`(1+2m_{\stackrel{~}{q}}^2/s(um_{\stackrel{~}{g}}^2)/s2m_{\stackrel{~}{g}}^2/s)`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,0,t;0,m_{\stackrel{~}{q}},m_{\stackrel{~}{g}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A122)
$`(2m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)+2m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2+2m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{\chi }}^2)`$
$`s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2s/(um_{\stackrel{~}{\chi }}^2)s^2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)`$
$`2/(tm_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2m_{\stackrel{~}{g}}^2/(um_{\stackrel{~}{\chi }}^2))`$
$`+`$ $`iC_0(m_{\stackrel{~}{\chi }}^2,0,t;0,m_{\stackrel{~}{q}},m_{\stackrel{~}{g}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A124)
$`(1+2m_{\stackrel{~}{q}}^2/s(um_{\stackrel{~}{g}}^2)/s+m_{\stackrel{~}{\chi }}^2/s3m_{\stackrel{~}{g}}^2/s)`$
$`+`$ $`iC_0(m_{\stackrel{~}{\chi }}^2,0,t;0,m_{\stackrel{~}{q}},m_{\stackrel{~}{g}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A128)
$`(2m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)+2m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{g}}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+2m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{g}}^2)`$
$`s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2s/(um_{\stackrel{~}{g}}^2)s^2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)`$
$`2/(tm_{\stackrel{~}{g}}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4m_{\stackrel{~}{g}}^2/(um_{\stackrel{~}{g}}^2))`$
$`+`$ $`iC_0(0,0,s;m_{\stackrel{~}{q}},m_{\stackrel{~}{g}},m_{\stackrel{~}{q}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A129)
$`+`$ $`iC_0(0,0,s;m_{\stackrel{~}{q}},m_{\stackrel{~}{g}},m_{\stackrel{~}{q}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A131)
$`(2m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)+s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+s^2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2))`$
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2,0,0;m_{\stackrel{~}{q}},0,m_{\stackrel{~}{q}},m_{\stackrel{~}{g}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^t^u|}`$ (A134)
$`(2m_{\stackrel{~}{q}}^2(um_{\stackrel{~}{g}}^2)/s+m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2/s5m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2/s2m_{\stackrel{~}{q}}^2+2m_{\stackrel{~}{q}}^4/s+2(um_{\stackrel{~}{g}}^2)/sm_{\stackrel{~}{g}}^2`$
$`m_{\stackrel{~}{\chi }}^2/sm_{\stackrel{~}{g}}^2+3/sm_{\stackrel{~}{g}}^4+s+(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{\chi }}^2+3m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,m_{\stackrel{~}{g}}^2,0,0;m_{\stackrel{~}{q}},0,m_{\stackrel{~}{q}},m_{\stackrel{~}{g}})(um_{\stackrel{~}{q}}^2)g_s^2(C_FN_C/2)2\overline{|^u^u|}`$ (A141)
$`(2m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{\chi }}^2)3m_{\stackrel{~}{q}}^2s/(um_{\stackrel{~}{g}}^2)`$
$`4m_{\stackrel{~}{q}}^2s^2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)+4m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4+m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2`$
$`3m_{\stackrel{~}{q}}^2/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+4m_{\stackrel{~}{q}}^4s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)2m_{\stackrel{~}{q}}^4/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2`$
$`+2m_{\stackrel{~}{q}}^4/(um_{\stackrel{~}{g}}^2)2s/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^4s/(um_{\stackrel{~}{\chi }}^2)m_{\stackrel{~}{g}}^2+2s/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2`$
$`+s^2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^2+s^2/(um_{\stackrel{~}{g}}^2)+s^3/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)`$
$`2/(um_{\stackrel{~}{\chi }}^2)/(um_{\stackrel{~}{g}}^2)m_{\stackrel{~}{g}}^6m_{\stackrel{~}{g}}^4/(um_{\stackrel{~}{\chi }}^2)+m_{\stackrel{~}{g}}^4/(um_{\stackrel{~}{g}}^2)).`$
The corresponding $`u`$-channel contribution is obtained by exchanging $`t`$ and $`u`$.
For the third box diagram in Fig.5 we find
$`\overline{|^{\mathrm{Box3}}|}^2`$ $`=`$ $`iB_0(t;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^u|}/s`$ (A142)
$`+`$ $`iB_0(t;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^t|}`$ (A144)
$`(m_{\stackrel{~}{g}}^2/((tm_{\stackrel{~}{\chi }}^2)(tm_{\stackrel{~}{g}}^2))1/(tm_{\stackrel{~}{\chi }}^2))`$
$`+`$ $`iB_0(u;m_{\stackrel{~}{g}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^u|}/s`$ (A145)
$`+`$ $`iB_0(u;m_{\stackrel{~}{g}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^t|}`$ (A148)
$`(m_{\stackrel{~}{g}}^2/((tm_{\stackrel{~}{\chi }}^2)(tm_{\stackrel{~}{g}}^2))+m_{\stackrel{~}{\chi }}^2/((tm_{\stackrel{~}{\chi }}^2)(um_{\stackrel{~}{\chi }}^2))+1/(tm_{\stackrel{~}{\chi }}^2)`$
$`+m_{\stackrel{~}{g}}^2/((tm_{\stackrel{~}{g}}^2)(um_{\stackrel{~}{g}}^2)))`$
$``$ $`iB_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^u|}/s`$ (A149)
$`+`$ $`iB_0(m_{\stackrel{~}{g}}^2;m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^t|}`$ (A151)
$`(m_{\stackrel{~}{\chi }}^2/((tm_{\stackrel{~}{\chi }}^2)(um_{\stackrel{~}{\chi }}^2)))`$
$``$ $`iB_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{g}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^u|}/s`$ (A152)
$`+`$ $`iB_0(m_{\stackrel{~}{\chi }}^2;m_{\stackrel{~}{g}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^t|}(m_{\stackrel{~}{g}}^2/((tm_{\stackrel{~}{g}}^2)(um_{\stackrel{~}{g}}^2)))`$ (A153)
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,0,t;m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^u|}`$ (A155)
$`(1(tm_{\stackrel{~}{g}}^2)/(2s)+m_{\stackrel{~}{\chi }}^2/(2s)m_{\stackrel{~}{g}}^2/(2s))`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,0,t;m_{\stackrel{~}{q}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^t|}`$ (A157)
$`(1/2+m_{\stackrel{~}{q}}^2/(2(tm_{\stackrel{~}{g}}^2))m_{\stackrel{~}{\chi }}^2/(2(tm_{\stackrel{~}{g}}^2)))`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,0,u;0,m_{\stackrel{~}{q}},m_{\stackrel{~}{g}})(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^u|}`$ (A159)
$`(1/2+m_{\stackrel{~}{q}}^2/s(tm_{\stackrel{~}{g}}^2)/(2s)m_{\stackrel{~}{g}}^2/s)`$
$`+`$ $`iC_0(m_{\stackrel{~}{g}}^2,0,u;0,m_{\stackrel{~}{q}},m_{\stackrel{~}{g}})(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^t|}`$ (A163)
$`(1/2+m_{\stackrel{~}{q}}^2s/(2(tm_{\stackrel{~}{\chi }}^2)(tm_{\stackrel{~}{g}}^2))+m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2/((tm_{\stackrel{~}{\chi }}^2)(um_{\stackrel{~}{\chi }}^2))`$
$`+1/2m_{\stackrel{~}{q}}^2/(tm_{\stackrel{~}{\chi }}^2)sm_{\stackrel{~}{g}}^2/(2(tm_{\stackrel{~}{\chi }}^2)(tm_{\stackrel{~}{g}}^2))+s/(2(tm_{\stackrel{~}{g}}^2))`$
$`m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2/((tm_{\stackrel{~}{\chi }}^2)(um_{\stackrel{~}{\chi }}^2))m_{\stackrel{~}{g}}^2/(2(tm_{\stackrel{~}{\chi }}^2)))`$
$`+`$ $`iC_0(m_{\stackrel{~}{\chi }}^2,0,u;m_{\stackrel{~}{g}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^u|}`$ (A165)
$`(1/2+(tm_{\stackrel{~}{g}}^2)/(2s)m_{\stackrel{~}{\chi }}^2/(2s)+m_{\stackrel{~}{g}}^2/(2s))`$
$`+`$ $`iC_0(m_{\stackrel{~}{\chi }}^2,0,u;m_{\stackrel{~}{g}},0,0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^t|}`$ (A168)
$`(1/2m_{\stackrel{~}{q}}^2s/(2(tm_{\stackrel{~}{\chi }}^2)(tm_{\stackrel{~}{g}}^2))m_{\stackrel{~}{q}}^2/(2(tm_{\stackrel{~}{g}}^2))`$
$`+sm_{\stackrel{~}{g}}^2/(2(tm_{\stackrel{~}{\chi }}^2)(tm_{\stackrel{~}{g}}^2))s/(2(tm_{\stackrel{~}{g}}^2))+m_{\stackrel{~}{\chi }}^2/(2(tm_{\stackrel{~}{g}}^2)))`$
$`+`$ $`iC_0(m_{\stackrel{~}{\chi }}^2,0,t;0,m_{\stackrel{~}{g}},m_{\stackrel{~}{q}})(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^u|}`$ (A170)
$`(m_{\stackrel{~}{q}}^2/s+(tm_{\stackrel{~}{g}}^2)/(2s)+m_{\stackrel{~}{g}}^2/s)`$
$`+`$ $`iC_0(m_{\stackrel{~}{\chi }}^2,0,t;0,m_{\stackrel{~}{g}},m_{\stackrel{~}{q}})(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^t|}`$ (A173)
$`(1/2+m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2/((tm_{\stackrel{~}{\chi }}^2)(tm_{\stackrel{~}{g}}^2))+m_{\stackrel{~}{q}}^2/(2(tm_{\stackrel{~}{\chi }}^2))m_{\stackrel{~}{g}}^4/((tm_{\stackrel{~}{\chi }}^2)(tm_{\stackrel{~}{g}}^2))`$
$`m_{\stackrel{~}{g}}^2/(2(tm_{\stackrel{~}{\chi }}^2)))`$
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,0,m_{\stackrel{~}{g}}^2,0;0,m_{\stackrel{~}{g}},m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^u|}`$ (A175)
$`(s(tm_{\stackrel{~}{g}}^2)+m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,0,m_{\stackrel{~}{g}}^2,0;0,m_{\stackrel{~}{g}},m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)g_s^2N_C2\overline{|^t^t|}`$ (A177)
$`(s(tm_{\stackrel{~}{g}}^2)+m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{g}}^2)`$
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,0,m_{\stackrel{~}{g}}^2,0;0,m_{\stackrel{~}{g}},m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)^2g_s^2N_C2\overline{|^t^u|}`$ (A179)
$`(1/2(tm_{\stackrel{~}{g}}^2)/(2s)+m_{\stackrel{~}{\chi }}^2/(2s)m_{\stackrel{~}{g}}^2/(2s))`$
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,0,m_{\stackrel{~}{g}}^2,0;0,m_{\stackrel{~}{g}},m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)^2g_s^2N_C2\overline{|^t^t|}`$ (A181)
$`(1+s/(2(tm_{\stackrel{~}{\chi }}^2))+s/(2(tm_{\stackrel{~}{g}}^2))m_{\stackrel{~}{\chi }}^2/(2(tm_{\stackrel{~}{g}}^2))+m_{\stackrel{~}{g}}^2/(2(tm_{\stackrel{~}{g}}^2)))`$
$`+`$ $`iD_0(m_{\stackrel{~}{\chi }}^2,0,m_{\stackrel{~}{g}}^2,0;0,m_{\stackrel{~}{g}},m_{\stackrel{~}{q}},0)(tm_{\stackrel{~}{q}}^2)^3g_s^2N_C2\overline{|^t^t|}`$ (A183)
$`(s/(2(tm_{\stackrel{~}{\chi }}^2)(tm_{\stackrel{~}{g}}^2))1/(2(tm_{\stackrel{~}{g}}^2))).`$
The corresponding $`u`$-channel contribution is obtained by exchanging $`t`$ and $`u`$.
The coefficients of the infrared singularities in the gluon exchange box diagrams can be simplified considerably. The gluino-exchange box diagram is infrared finite. For the first box diagram in Fig. 5 in the $`t`$-channel we find
$`\overline{|_{\mathrm{IR}}^{\mathrm{Box1}}|}^2`$ $`=`$ $`C_ϵ[{\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{1}{ϵ}}\mathrm{log}\left({\displaystyle \frac{(tm_{\stackrel{~}{q}}^2)^2s}{(m_{\stackrel{~}{g}}^2m_{\stackrel{~}{q}}^2)(m_{\stackrel{~}{\chi }}^2m_{\stackrel{~}{q}}^2)Q^2}}\right)+{\displaystyle \frac{1}{ϵ}}\mathrm{log}\left({\displaystyle \frac{m_{\stackrel{~}{q}}^2t}{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2}}\right){\displaystyle \frac{tm_{\stackrel{~}{q}}^2}{tm_{\stackrel{~}{g}}^2}}`$ (A185)
$`+{\displaystyle \frac{1}{ϵ}}\mathrm{log}\left({\displaystyle \frac{m_{\stackrel{~}{q}}^2t}{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}}\right){\displaystyle \frac{tm_{\stackrel{~}{q}}^2}{tm_{\stackrel{~}{\chi }}^2}}\left]4\right(C_F{\displaystyle \frac{N_C}{2}})g_s^2[\overline{|^t^t|}+\overline{|^t^u|}],`$
and the $`u`$-channel contribution can be obtained by exchanging $`t`$ and $`u`$. The result is completely symmetric under the exchange $`m_{\stackrel{~}{g}}m_{\stackrel{~}{\chi }}`$.
The infrared singular pieces of the third $`t`$-channel box diagram in Fig. 5 are
$`\overline{|_{\mathrm{IR}}^{\mathrm{Box3}}|}^2`$ $`=`$ $`C_ϵ\left[{\displaystyle \frac{1}{2ϵ^2}}+{\displaystyle \frac{1}{2ϵ}}\mathrm{log}\left({\displaystyle \frac{Q^2}{m_{\stackrel{~}{g}}^2}}\right){\displaystyle \frac{1}{ϵ}}\mathrm{log}\left({\displaystyle \frac{m_{\stackrel{~}{q}}^2t}{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}}\right){\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}{tm_{\stackrel{~}{\chi }}^2}}{\displaystyle \frac{1}{ϵ}}\mathrm{log}\left({\displaystyle \frac{m_{\stackrel{~}{g}}^2u}{m_{\stackrel{~}{g}}^2}}\right)\right]`$ (A187)
$`\times 4{\displaystyle \frac{N_C}{2}}g_s^2[\overline{|^t^t|}+\overline{|^t^u|}].`$
For the $`u`$-channel result $`t`$ and $`u`$ have to be exchanged.
Finally we sum the infrared singularities encountered in the virtual corrections and separate them into $`C_F`$ and $`N_C`$ color classes:
$`\overline{|_{\mathrm{IR}}^V|}^2`$ $`=`$ $`C_ϵ\left[{\displaystyle \frac{1}{ϵ^2}}+{\displaystyle \frac{3}{2ϵ}}+{\displaystyle \frac{1}{ϵ}}\mathrm{log}\left({\displaystyle \frac{Q^2}{s}}\right)\right]4C_Fg_s^2\overline{|^B|}^2`$ (A189)
$`+C_ϵ\left[{\displaystyle \frac{1}{ϵ}}+{\displaystyle \frac{1}{ϵ}}\mathrm{log}\left({\displaystyle \frac{(tm_{\stackrel{~}{g}}^2)(um_{\stackrel{~}{g}}^2)}{sm_{\stackrel{~}{g}}^2}}\right)\right]4{\displaystyle \frac{N_C}{2}}g_s^2\overline{|^B|}^2.`$
## B Scalar Integrals
The tensor integrals that occur in virtual loop diagrams can be reduced to one-, two-, three-, and four-point integrals that are scalar functions of the loop momentum $`q`$ . The general scalar one-point integral is defined as
$`A_0(m)`$ $`=`$ $`Q^{2ϵ}{\displaystyle \frac{\text{d}^nq}{(2\pi )^n}\frac{1}{q^2m^2}}.`$ (B1)
The general scalar two-point integral is
$`B_0(p^2;m_1,m_2)`$ $`=`$ $`Q^{2ϵ}{\displaystyle \frac{\text{d}^nq}{(2\pi )^n}\frac{1}{[q^2m_1^2][(q+p)^2m_2^2]}},`$ (B2)
and its derivative is
$`B_0^{}(p^2;m_1,m_2)`$ $`=`$ $`{\displaystyle \frac{}{q^2}}B_0(q^2;m_1,m_2)|_{q^2=p^2}.`$ (B3)
The general scalar three-point integral is
$`C_0(p_1^2,p_2^2,(p_1+p_2)^2;m_1,m_2,m_3)`$ $`=`$ $`Q^{2ϵ}`$ (B5)
$`\times {\displaystyle }{\displaystyle \frac{\text{d}^nq}{(2\pi )^n}}{\displaystyle \frac{1}{[q^2m_1^2][(q+p_1)^2m_2^2][(q+p_1+p_2)^2m_3^2]}},`$
and the general scalar four-point integral is
$`D_0(p_1^2,p_2^2,p_3^2,(p_1+p_2+p_3)^2,(p_1+p_2)^2,(p_2+p_3)^2;m_1,m_2,m_3,m_4)`$ $`=`$ $`Q^{2ϵ}`$ (B7)
$`\times {\displaystyle }{\displaystyle \frac{\text{d}^nq}{(2\pi )^n}}{\displaystyle \frac{1}{[q^2m_1^2][(q+p_1)^2m_2^2][(q+p_1+p_2)^2m_3^2][(q+p_1+p_2+p_3)^2m_4^2]}}.`$
The $`p_i,i=1\mathrm{}4,`$ are the four-momenta of the external particles, and the $`m_i,i=1\mathrm{}4,`$ are the masses of the adjacent internal particles.
The scalar two- and three-point integrals relevant for the associated production of gluinos and gauginos were calculated previously in a different physical context . We recalculated them and checked that the results in Ref. are correct. The divergent four-point integrals that contribute to the gluon exchange box diagrams in Fig. 5 were unknown, and they are presented here for the first time. We calculate the absorptive parts with Cutkosky cutting rules and the real parts with dispersion techniques.
The exchange of a massless gluon between the initial state quark and antiquark in the first box diagram of Fig. 5 leads to the following divergent four-point integral:
$`D_0(0,0,m_{\stackrel{~}{g}}^2,m_{\stackrel{~}{\chi }}^2,s,t;0,0,0,m_{\stackrel{~}{q}})`$ $`=`$ $`iC_ϵ{\displaystyle \frac{1}{s(tm_{\stackrel{~}{q}}^2)}}[{\displaystyle \frac{1}{ϵ^2}}{\displaystyle \frac{1}{ϵ}}\mathrm{ln}{\displaystyle \frac{s(tm_{\stackrel{~}{q}}^2)^2}{(m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2)(m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2)Q^2}}`$ (B11)
$`2\text{Li}_2\left(1+{\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2}{tm_{\stackrel{~}{q}}^2}}\right)2\text{Li}_2\left(1+{\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}{tm_{\stackrel{~}{q}}^2}}\right)\text{Li}_2\left(1+{\displaystyle \frac{(m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2)(m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2)}{sm_{\stackrel{~}{q}}^2}}\right)`$
$`{\displaystyle \frac{\pi ^2}{4}}+{\displaystyle \frac{1}{2}}\mathrm{ln}^2\left({\displaystyle \frac{s}{Q^2}}\right){\displaystyle \frac{1}{2}}\mathrm{ln}^2\left({\displaystyle \frac{s}{m_{\stackrel{~}{q}}^2}}\right)+2\mathrm{ln}\left({\displaystyle \frac{s}{Q^2}}\right)\mathrm{ln}\left({\displaystyle \frac{tm_{\stackrel{~}{q}}^2}{m_{\stackrel{~}{q}}^2}}\right)`$
$`\mathrm{ln}\left({\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2}{Q^2}}\right)\mathrm{ln}\left({\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{g}}^2}{m_{\stackrel{~}{q}}^2}}\right)\mathrm{ln}\left({\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}{Q^2}}\right)\mathrm{ln}\left({\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}{m_{\stackrel{~}{q}}^2}}\right)].`$
This integral also contributes to the pair production of gauginos of unequal mass . In the limit of two final state particles of equal mass our result agrees with Ref. .
In the third box diagram of Fig. 5, the exchange of a massless gluon between the the final state gluino and the initial state antiquark gives rise to a second divergent four-point integral:
$`D_0(0,m_{\stackrel{~}{g}}^2,0,m_{\stackrel{~}{\chi }}^2,t,u;m_{\stackrel{~}{q}},m_{\stackrel{~}{g}},0,0)`$ $`=`$ $`iC_ϵ{\displaystyle \frac{1}{(tm_{\stackrel{~}{q}}^2)(um_{\stackrel{~}{g}}^2)}}[{\displaystyle \frac{1}{2ϵ^2}}{\displaystyle \frac{1}{ϵ}}(\mathrm{ln}\left({\displaystyle \frac{t+m_{\stackrel{~}{q}}^2}{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}}\right)`$ (B14)
$`+\mathrm{ln}\left({\displaystyle \frac{u+m_{\stackrel{~}{g}}^2}{m_{\stackrel{~}{g}}Q}}\right)){\displaystyle \frac{\pi ^2}{8}}+\mathrm{ln}^2({\displaystyle \frac{u+m_{\stackrel{~}{g}}^2}{m_{\stackrel{~}{g}}Q}})+2\mathrm{ln}({\displaystyle \frac{u+m_{\stackrel{~}{g}}^2}{m_{\stackrel{~}{g}}Q}}\left)\mathrm{ln}\right({\displaystyle \frac{t+m_{\stackrel{~}{q}}^2}{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}})`$
$`+\text{Li}_2(1+{\displaystyle \frac{um_{\stackrel{~}{g}}^2}{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}})+\text{Li}_2(1+{\displaystyle \frac{m_{\stackrel{~}{q}}^2(um_{\stackrel{~}{g}}^2)}{m_{\stackrel{~}{g}}^2(m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2)}})2\text{Li}_2(1+{\displaystyle \frac{m_{\stackrel{~}{q}}^2m_{\stackrel{~}{\chi }}^2}{tm_{\stackrel{~}{q}}^2}})].`$
Again, our result agrees with Ref. for the case of two final state particles of equal mass.
The supersymmetric counterpart of the first box diagram is the second box diagram in Fig. 5. Since supersymmetry is broken and the gluino propagator and the two squark propagators are massive, this diagram does not have infrared divergences. The corresponding four-point integral can be expressed in terms of 16 dilogarithms .
## C Soft Gluon Emission Contribution
In this Appendix, we collect the expressions for the finite pieces of the soft gluon emission contribution. The finite pieces remain after mass-factorization and cancellation of soft poles between soft gluon emission and virtual contributions, as described in the text. We begin with the $`C_F`$ color class,
$`{\displaystyle \frac{d^2\widehat{\sigma }^S}{dt_2du_2}}`$ $`=`$ $`{\displaystyle \frac{d^2\widehat{\sigma }^B}{dt_2du_2}}\left({\displaystyle \frac{C_F\alpha _S}{\pi }}\right)\{\mathrm{Li}_2\left({\displaystyle \frac{u_2t_2sm_2^2}{(s+t_2)(s+u_2)}}\right)`$ (C3)
$`+{\displaystyle \frac{1}{2}}\mathrm{log}^2\left({\displaystyle \frac{\mu ^2}{m_1^2\delta ^2}}\right)+\mathrm{log}\left({\displaystyle \frac{(s+t_2)(s+u_2)}{sm_1^2}}\right)\mathrm{log}\left({\displaystyle \frac{\mu ^2}{m_1^2\delta ^2}}\right)`$
$`+{\displaystyle \frac{1}{2}}\mathrm{log}^2\left({\displaystyle \frac{(s+t_2)(s+u_2)}{sm_1^2}}\right)\},`$
where $`\delta =\mathrm{\Delta }/m_1^2`$ is the cut-off between hard and soft emission mentioned in the text. Its appearance in these terms is the explicit cut-off dependence that matches the implicit logarithmic behavior of the hard real emission contributions.
The soft terms associated with the $`N_C`$ color class are
$`{\displaystyle \frac{d^2\widehat{\sigma }^S}{dt_2du_2}}`$ $`=`$ $`{\displaystyle \frac{d^2\widehat{\sigma }^B}{dt_2du_2}}\left({\displaystyle \frac{N_C\alpha _S}{2\pi }}\right)\{2+\mathrm{Li}_2\left({\displaystyle \frac{u_2t_2sm_2^2}{(s+t_2)(s+u_2)}}\right)`$ (C6)
$`+{\displaystyle \frac{1}{2}}\mathrm{log}^2\left({\displaystyle \frac{\mu ^2}{m_1^2\delta ^2}}\right)`$
$`+[\mathrm{log}\left({\displaystyle \frac{(s+t_2)(s+u_2)}{sm_1^2}}\right)1]\mathrm{log}\left({\displaystyle \frac{(s+t_2)(s+u_2)}{sm_1^2}}\right)\}.`$
Again logarithmic dependence on the hard/soft cut-off $`\delta `$ is apparent.
## D Hard Gluon Emission Contribution
In this Appendix, we collect explicit expressions for the thirty-six matrix elements that contribute to the real emission of a gluon in the 2-to-3 partonic subprocess $`q\overline{q}g\stackrel{~}{g}\stackrel{~}{\chi }`$.
The hard gluon emission cross section (in four dimensions) is
$`{\displaystyle \frac{d^3\widehat{\sigma }^h}{ds_4dt_2du_2}}`$ $`=`$ $`{\displaystyle \frac{d^3\widehat{\sigma }_1^g}{ds_4dt_2du_2}}+{\displaystyle \frac{\alpha _S\widehat{\alpha }_S}{16\pi ^2}}{\displaystyle \frac{s_4\delta (s+t_2+u_1s_4)}{36s^2(s_4+m_1^2)}}{\displaystyle \underset{i=\mathrm{1..8}}{}}{\displaystyle \underset{j=i\mathrm{..8}}{}}\widehat{M}_{ij}^g,`$ (D1)
where $`i`$ and $`j`$ label the diagrams in Fig. 6. The remainder of the factorization process is
$`{\displaystyle \frac{d^3\widehat{\sigma }_1^g}{ds_4dt_2du_2}}`$ $`=`$ $`{\displaystyle \frac{C_F\alpha _S\widehat{\alpha }_S\delta (s+t_2+u_1s_4)}{36\pi s^2}}\mathrm{log}\left({\displaystyle \frac{\mu ^2(s_4+m_1^2)}{s_4^2}}\right)`$ (D6)
$`\times \{\left({\displaystyle \frac{s_4^22s_4(s+u_2)+2(s+u_2)^2}{s_4(s+u_2)}}\right)({\displaystyle \frac{X_tt_2}{(tm_{\stackrel{~}{q}_t}^2)^2}}`$
$`+{\displaystyle \frac{2X_{tu}sm_1m_2}{(tm_{\stackrel{~}{q}_t}^2)[(\mathrm{\Delta }_ust_2)(s+u_2)+ss_4]}}+{\displaystyle \frac{X_uu_2[s_4u_2u_1(s+u_2)]}{[(\mathrm{\Delta }_ust_2)(s+u_2)+ss_4]^2}})`$
$`+\left({\displaystyle \frac{s_4^22s_4(s+t_2)+2(s+t_2)^2}{s_4(s+t_2)}}\right)({\displaystyle \frac{X_tt_2[s_4t_2t_1(s+t_2)]}{[(\mathrm{\Delta }_tsu_2)(s+t_2)+ss_4]^2}}`$
$`+{\displaystyle \frac{2X_{tu}sm_1m_2}{(um_{\stackrel{~}{q}_u}^2)[(\mathrm{\Delta }_tsu_2)(s+t_2)+ss_4]}}+{\displaystyle \frac{X_uu_2}{(um_{\stackrel{~}{q}_u}^2)^2}})\}.`$
The elements $`\widehat{M}_{ij}^g`$ are
$`\widehat{M}_{11}^g`$ $`=`$ $`{\displaystyle \frac{16C_F\pi X_tt_2(s_4+m_1^2)}{(tm_{\stackrel{~}{q}_t}^2)^2(s+u_2)}},`$ (D7)
$`\widehat{M}_{12}^g`$ $`=`$ $`\left({\displaystyle \frac{8(C_FN_C/2)X_t}{(tm_{\stackrel{~}{q}_t}^2)(\mathrm{\Delta }_tsu_2)}}\right)`$ (D11)
$`\times \{s[2t_2(s+u_2)+s_4(s+t_2+u_2)]\widehat{I}\left({\displaystyle \frac{1}{t^{}u^{}}}\right)+[st_2]\widehat{I}\left({\displaystyle \frac{u^{}}{t^{}}}\right)`$
$`+[\mathrm{\Delta }_ts(\mathrm{\Delta }_t+t_1)+t_2(s(m_2^2m_1^2)+\mathrm{\Delta }_t(t_2+u_2\mathrm{\Delta }_t)u_2t_2)]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}}}\right)`$
$`+[\mathrm{\Delta }_t(s+u_2t_2)+t_2(\mathrm{\Delta }_tsu_2)]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)+[st_2]\widehat{I}\left({\displaystyle \frac{u_{7\mathrm{\Delta }}}{t^{}}}\right)\},`$
$`\widehat{M}_{13}^g`$ $`=`$ $`{\displaystyle \frac{8N_CX_tt_2(s+u_2m_1^2)}{s_4(tm_{\stackrel{~}{q}_t}^2)^2}}\widehat{I}(1),`$ (D12)
$`\widehat{M}_{14}^g`$ $`=`$ $`\left({\displaystyle \frac{4(C_FN_C/2)X_tt_2}{(tm_{\stackrel{~}{q}_t}^2)^2}}\right)\{[4t_1(s+u_2)]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{7\mathrm{\Delta }}}}\right)`$ (D14)
$`+[2(2t_2+2m_2^2su_2)]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)2\widehat{I}\left({\displaystyle \frac{u^{}}{u_{7\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{15}^g`$ $`=`$ $`{\displaystyle \frac{16(C_FN_C/2)X_{tu}s^2m_1m_2}{(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}\widehat{I}\left({\displaystyle \frac{1}{t^{}u^{}}}\right),`$ (D15)
$`\widehat{M}_{16}^g`$ $`=`$ $`\left({\displaystyle \frac{16C_FX_{tu}m_1m_2}{(tm_{\stackrel{~}{q}_t}^2)}}\right)\{[\mathrm{\Delta }_ust_2]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)`$ (D17)
$`{\displaystyle \frac{2\pi (s_4+m_1^2)(\mathrm{\Delta }_ust_2)}{s_4[(\mathrm{\Delta }_ust_2)(s+u_2)+ss_4]}}+{\displaystyle \frac{2\pi (s_4+m_1^2)}{s_4(s+u_2)}}\},`$
$`\widehat{M}_{17}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2}{s_4(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}\right)\left\{[s+u_2]\widehat{I}\left({\displaystyle \frac{t^{}}{u^{}}}\right)+[st_2]\widehat{I}(1)\right\},`$ (D18)
$`\widehat{M}_{18}^g`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_{tu}m_1m_2}{(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}\right)\{2[s(\mathrm{\Delta }_u+m_2^2m_1^2)u_2(\mathrm{\Delta }_ut_2)]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)`$ (D20)
$`+[2t_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{22}^g`$ $`=`$ $`\left(8C_FX_t\right)\{[\mathrm{\Delta }_t(\mathrm{\Delta }_t+t_1)]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}^2}}\right)[2\mathrm{\Delta }_t+t_1]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}}}\right)+\mathrm{\Delta }_t\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}^2}}\right)`$ (D23)
$`\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)\left({\displaystyle \frac{2\pi (s_4+m_1^2)}{s_4}}\right)({\displaystyle \frac{(s+t_2)\mathrm{\Delta }_t(\mathrm{\Delta }_t+t_1)}{[(\mathrm{\Delta }_ust_2)(s+u_2)+ss_4]^2}}`$
$`{\displaystyle \frac{2\mathrm{\Delta }_t+t_1}{[(\mathrm{\Delta }_ust_2)(s+u_2)+ss_4]}}+{\displaystyle \frac{1}{(s+t_2)}})\},`$
$`\widehat{M}_{23}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_t}{s_4(tm_{\stackrel{~}{q}_t}^2)}}\right)\times `$ (D27)
$`\{[\mathrm{\Delta }_t((s+t_2)(\mathrm{\Delta }_t+m_2^2m_1^2)+s_4t_2)+t_2(ss_4s^2st_2su_2t_2u_2)]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}}}\right)`$
$`+[s+t_2]\widehat{I}\left({\displaystyle \frac{u_{7\mathrm{\Delta }}}{t^{}}}\right)+[\mathrm{\Delta }_t(m_1^2+m_2^2+s_4)t_2(s+u_2\mathrm{\Delta }_t)]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)`$
$`[t+m_1^2+s_4]\widehat{I}(1)\},`$
$`\widehat{M}_{24}^g`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_t}{(tm_{\stackrel{~}{q}_t}^2)}}\right)\{[4t_2\mathrm{\Delta }_t(\mathrm{\Delta }_t+m_2^2m_1^2)]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}^2}}\right)+[2\mathrm{\Delta }_t(t_2+2m_2^2)]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}^2}}\right)`$ (D29)
$`[4t_2(2\mathrm{\Delta }_t+m_2^2m_1^2)]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}}}\right)2[t_2+2m_2^2]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{25}^g`$ $`=`$ $`\left({\displaystyle \frac{16C_FX_{tu}m_1m_2}{(um_{\stackrel{~}{q}_u}^2)}}\right)\{[\mathrm{\Delta }_tsu_2]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}}}\right)`$ (D31)
$`{\displaystyle \frac{2\pi (\mathrm{\Delta }_tsu_2)(s_4+m_1^2)}{s_4[(\mathrm{\Delta }_tsu_2)(s+t_2)+ss_4]}}+{\displaystyle \frac{2\pi (s_4+m_1^2)}{s_4(s+t_2)}}\},`$
$`\widehat{M}_{26}^g`$ $`=`$ $`\left({\displaystyle \frac{16(C_FN_C/2)X_{tu}sm_1m_2}{(\mathrm{\Delta }_tsu_2)(\mathrm{\Delta }_ust_2)}}\right)\{[s]\widehat{I}\left({\displaystyle \frac{1}{t^{}u^{}}}\right)+[\mathrm{\Delta }_tu_2]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}}}\right)`$ (D33)
$`+[\mathrm{\Delta }_ut_2]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)+[\mathrm{\Delta }_u+\mathrm{\Delta }_tst_2u_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}u_{7\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{27}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2}{s_4(um_{\stackrel{~}{q}_u}^2)}}\right)\{[\mathrm{\Delta }_ts+s^2ss_4+\mathrm{\Delta }_tt_2+st_2su_2t_2u_2]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}}}\right)`$ (D35)
$`+[s+u_22\mathrm{\Delta }_t]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)+2\widehat{I}(1)\},`$
$`\widehat{M}_{28}^g`$ $`=`$ $`\left({\displaystyle \frac{4(C_FN_C/2)X_{tu}m_1m_2}{(um_{\stackrel{~}{q}_u}^2)(\mathrm{\Delta }_ust_2}}\right)\{2[su_1(s+t_2)(\mathrm{\Delta }_tu_2)]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}}}\right)`$ (D37)
$`+2[s+t_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)+[2\mathrm{\Delta }_t(s+t_2)+2\mathrm{\Delta }_uu_2+2su_1]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}u_{7\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{33}^g`$ $`=`$ $`\left({\displaystyle \frac{8N_CX_tt_2}{s_4^2(tm_{\stackrel{~}{q}_t}^2)^2}}\right)\left\{2m_1^2(s+u_2)\widehat{I}(1)+s_4\widehat{I}(u^{})\right\},`$ (D38)
$`\widehat{M}_{34}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_tt_2}{s_4(tm_{\stackrel{~}{q}_t}^2)^2}}\right)\{[\mathrm{\Delta }_t(s_42t)2m_1^2(u_2+2s)]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)`$ (D40)
$`+[m_2^2+m_1^2s+t_2u_2]\widehat{I}(1)\},`$
$`\widehat{M}_{35}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2}{s_4(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}\right)\left\{[s+t_2]\widehat{I}\left({\displaystyle \frac{u^{}}{t^{}}}\right)+[su_2]\widehat{I}(1)\right\},`$ (D41)
$`\widehat{M}_{36}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2}{s_4(tm_{\stackrel{~}{q}_t}^2)}}\right)\{[\mathrm{\Delta }_u(s+u_2)+s^2ss_4st_2+su_2t_2u_2]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)`$ (D43)
$`+[s+t_22\mathrm{\Delta }_u]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)+2\widehat{I}(1)\},`$
$`\widehat{M}_{37}^g`$ $`=`$ $`{\displaystyle \frac{16N_CX_{tu}sm_1m_2(s_4+2m_1^2)}{s_4^2(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}\widehat{I}(1),`$ (D44)
$`\widehat{M}_{38}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2}{s_4(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}\right)\{[u_2(t_2\mathrm{\Delta }_u)+s(\mathrm{\Delta }_us_4m_2^23m_1^2)]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)`$ (D46)
$`+[s+t_2]\widehat{I}\left({\displaystyle \frac{u^{}}{u_{6\mathrm{\Delta }}}}\right)+[u_2s]\widehat{I}(1)\},`$
$`\widehat{M}_{44}^g`$ $`=`$ $`\left({\displaystyle \frac{8C_FX_tt_2}{(tm_{\stackrel{~}{q}_t}^2)^2}}\right)\left\{[\mathrm{\Delta }_t(\mathrm{\Delta }_ttm_1^2)]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}^2}}\right)+[t+m_1^22\mathrm{\Delta }_t]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)+\widehat{I}(1)\right\},`$ (D47)
$`\widehat{M}_{45}^g`$ $`=`$ $`\left({\displaystyle \frac{8C_FX_{tu}m_1m_2}{(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}\right)\left\{[t_2(u_2\mathrm{\Delta }_t)+s(\mathrm{\Delta }_t+m_2^2m_1^2)]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{7\mathrm{\Delta }}}}\right)+u_2\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)\right\},`$ (D48)
$`\widehat{M}_{46}^g`$ $`=`$ $`\left({\displaystyle \frac{4(C_FN_C/2)X_{tu}m_1m_2}{(tm_{\stackrel{~}{q}_t}^2)(\mathrm{\Delta }_tsu_2)}}\right)\{2[(t_2\mathrm{\Delta }_u)(s+u_2)+st_1]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)`$ (D50)
$`+2[s+u_2]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)+[2\mathrm{\Delta }_u(s+u_2)+2st_1+2\mathrm{\Delta }_tt_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}u_{7\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{47}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2}{s_4(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}\right)\{[t_2(u_2\mathrm{\Delta }_t)+\mathrm{\Delta }_tss(s_4+m_2^2+3m_1^2)]\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)`$ (D52)
$`+[s+u_2]\widehat{I}\left({\displaystyle \frac{t^{}}{u_{7\mathrm{\Delta }}}}\right)+[t_2s]\widehat{I}(1)\},`$
$`\widehat{M}_{48}^g`$ $`=`$ $`\left({\displaystyle \frac{8(C_FN_C/2)X_{tu}sm_1m_2}{(tm_{\stackrel{~}{q}_t}^2)(um_{\stackrel{~}{q}_u}^2)}}\right)\{[2s+2m_1^2\mathrm{\Delta }_t\mathrm{\Delta }_u+t+u]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}u_{7\mathrm{\Delta }}}}\right)`$ (D54)
$`+\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)+\widehat{I}\left({\displaystyle \frac{1}{u_{7\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{55}^g`$ $`=`$ $`{\displaystyle \frac{16C_F\pi X_uu_2(s_4+m_1^2)}{(um_{\stackrel{~}{q}_u}^2)^2(s+t_2)}},`$ (D55)
$`\widehat{M}_{56}^g`$ $`=`$ $`\left({\displaystyle \frac{8(C_FN_C/2)X_u}{(um_{\stackrel{~}{q}_u}^2)(\mathrm{\Delta }_ust_2)}}\right)`$ (D59)
$`\times \{[\mathrm{\Delta }_u^2(su_2)+\mathrm{\Delta }_u(su_1+u_2(t_2+u_2))+su_2(m_2^2m_1^2)t_2u_2^2]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)`$
$`+[su_2]\widehat{I}\left({\displaystyle \frac{u_{6\mathrm{\Delta }}}{u^{}}}\right)+[\mathrm{\Delta }_u(s+t_2)u_2(s+t_2)]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)`$
$`+s[(s+t_2)(s+t_1)+u_2u_1]\widehat{I}\left({\displaystyle \frac{1}{u^{}t^{}}}\right)+[su_2]\widehat{I}\left({\displaystyle \frac{t^{}}{u^{}}}\right)\},`$
$`\widehat{M}_{57}^g`$ $`=`$ $`\left({\displaystyle \frac{8N_CX_uu_2[s+t_2m_1^2]}{s_4(um_{\stackrel{~}{q}_u}^2)^2}}\right)\widehat{I}(1),`$ (D60)
$`\widehat{M}_{58}^g`$ $`=`$ $`\left({\displaystyle \frac{4(C_FN_C/2)X_uu_2}{(um_{\stackrel{~}{q}_u}^2)^2}}\right)\{4u_1[s+t_2]\widehat{I}\left({\displaystyle \frac{1}{t^{}u_{6\mathrm{\Delta }}}}\right)+2[3u_1+2m_1^2s_4]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)`$ (D62)
$`2\widehat{I}\left({\displaystyle \frac{t^{}}{u_{6\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{66}^g`$ $`=`$ $`\left(8C_FX_u\right)\{\mathrm{\Delta }_u[\mathrm{\Delta }_u+u_1]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}^2}}\right)[2\mathrm{\Delta }_u+u_1]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)`$ (D66)
$`+\mathrm{\Delta }_u\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}^2}}\right)\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)`$
$`\left({\displaystyle \frac{2\pi (s_4+m_1^2)}{s_4}}\right)({\displaystyle \frac{\mathrm{\Delta }_u(\mathrm{\Delta }_u+u_1)(s+u_2)}{[(\mathrm{\Delta }_ust_2)(s+u_2)+ss_4]^2}}`$
$`{\displaystyle \frac{2\mathrm{\Delta }_u+u_1}{[(\mathrm{\Delta }_ust_2)(s+u_2)+ss_4]}}+{\displaystyle \frac{1}{s+u_2}})\},`$
$`\widehat{M}_{67}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_u}{s_4(um_{\stackrel{~}{q}_u}^2)}}\right)\{[s+u_2]\widehat{I}\left({\displaystyle \frac{u_{6\mathrm{\Delta }}}{u^{}}}\right)[s_4+m_1^2+u]\widehat{I}(1)`$ (D70)
$`+[\mathrm{\Delta }_u^2(s+u_2)+\mathrm{\Delta }_u(s(m_2^2m_1^2)`$
$`+u_2(m_2^2m_1^2))t_2u_2)]\widehat{I}({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}})`$
$`+[\mathrm{\Delta }_u(s+t+u)+u_2(\mathrm{\Delta }_ust_2)]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{68}^g`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_u}{(um_{\stackrel{~}{q}_u}^2)}}\right)\{4\mathrm{\Delta }_uu_2[\mathrm{\Delta }_u+m_2^2m_1^2]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}^2}}\right)+2\mathrm{\Delta }_u[u_2+2m_2^2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}^2}}\right)`$ (D72)
$`+4u_2[2\mathrm{\Delta }_um_2^2+m_1^2]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)2[u_2+2m_2^2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{77}^g`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_uu_2}{s_4^2(um_{\stackrel{~}{q}_u}^2)^2}}\right)\left\{4m_1^2[s+t_2]\widehat{I}(1)+2s_4\widehat{I}(t^{})\right\},`$ (D73)
$`\widehat{M}_{78}^g`$ $`=`$ $`\left({\displaystyle \frac{2N_CX_uu_2}{s_4(um_{\stackrel{~}{q}_u}^2)^2}}\right)\{2[m_2^2+m_1^2st_2+u_2]\widehat{I}(1)`$ (D76)
$`+[s_4(s+t_2)+\mathrm{\Delta }_u(\mathrm{\Delta }_u3m_2^25m_1^23u_2)`$
$`+(\mathrm{\Delta }_ust_2)(\mathrm{\Delta }_u+m_2^2+3m_1^2+s+t_2+u_2)]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{88}^g`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_uu_2}{(um_{\stackrel{~}{q}_u}^2)^2}}\right)\{2\mathrm{\Delta }_u[\mathrm{\Delta }_u+m_2^2+m_1^2+u_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}^2}}\right)`$ (D78)
$`+2[u_2+m_2^2+m_1^22\mathrm{\Delta }_u]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)+2\widehat{I}(1)\}.`$
## E (Anti-) Quark Emission Contributions
In this Appendix we present the matrix elements for the contributions from emission of an additional light quark or anti-quark in the final state. These matrix elements may be obtained from those for gluon emission upon crossing one of the initial state partons and the final state gluon. However, we compute them ab initio in order to provide a check on the validity of the results. In all cases, the result of the crossing agrees with the explicit computation. We limit our presentation to the quark emission contributions since the anti-quark emission contributions may be obtained from these expressions by replacements of $`tu`$, $`t^{}u^{}`$, and $`u_6u_7`$ everywhere.
The quark emission contribution has two parts, consisting of the remainder of terms that are collinear singular, and thus have the angular integrations done analytically, as well as a large class of terms which are collinear finite and for which the angular integrations are done numerically. The matrix elements with collinear singularities may be written in a form very similar to the gluon emission cross section, with
$`{\displaystyle \frac{d^3\widehat{\sigma }^h}{ds_4dt_2du_2}}`$ $`=`$ $`{\displaystyle \frac{d^3\widehat{\sigma }_1^q}{ds_4dt_2du_2}}+{\displaystyle \frac{\alpha _S\widehat{\alpha }_S}{16\pi ^2}}{\displaystyle \frac{s_4\delta (s+t_2+u_1s_4)}{96s^2(s_4+m_1^2)}}{\displaystyle \underset{i=\mathrm{1..2}}{}}{\displaystyle \underset{j=i\mathrm{..8}}{}}\widehat{M}_{ij}^q.`$ (E1)
The remainder of the factorization is
$`{\displaystyle \frac{d^3\widehat{\sigma }_1^q}{ds_4dt_2du_2}}`$ $`=`$ $`{\displaystyle \frac{\alpha _S\widehat{\alpha }_SC_F}{48\pi s^2}}\left(1+\mathrm{log}\left[{\displaystyle \frac{\mu ^2(s_4+m_1^2)}{s_4^2}}\right]\right){\displaystyle \frac{1}{2}}\left({\displaystyle \frac{2s_4^22s_4(s+u_2)+(s+u_2)^2}{(s+u_2)^2}}\right)`$ (E4)
$`\{{\displaystyle \frac{X_tt_2}{(tm_{\stackrel{~}{q}_t}^2)^2}}+{\displaystyle \frac{2X_{tu}sm_1m_2}{(tm_{\stackrel{~}{q}_t}^2)[(s+u_2)(\mathrm{\Delta }_ust_2)+ss_4]}}`$
$`+{\displaystyle \frac{X_uu_2(u_2s_4u_1(s+u_2))}{[(s+u_2)(\mathrm{\Delta }_ust_2)+ss_4]^2}}\}.`$
After partial-fractionation, the collinearly-divergent pieces of the hard matrix elements are
$`\widehat{M}_{11}^q`$ $`=`$ $`(8C_FX_uu_2)\{\mathrm{\Delta }_u\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}^2}}\right)+\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)+{\displaystyle \frac{\mathrm{\Delta }_u2\pi (s_4+m_1^2)(s+u_2)}{s_4[(s+u_2)(\mathrm{\Delta }_ust_2)+ss_4]^2}}`$ (E6)
$`{\displaystyle \frac{2\pi (s_4+m_1^2)}{s_4[(s+u_2)(\mathrm{\Delta }_ust_2)+ss_4]}}\},`$
$`\widehat{M}_{12}^q`$ $`=`$ $`\left({\displaystyle \frac{16C_FX_{tu}sm_1m_2}{(tm_{\stackrel{~}{q}_t}^2)}}\right)\left\{\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right){\displaystyle \frac{2\pi (s_4+m_1^2)}{s_4[(s+u_2)(\mathrm{\Delta }_ust_2)+ss_4]}}\right\},`$ (E7)
$`\widehat{M}_{13}^q`$ $`=`$ $`\left({\displaystyle \frac{4(C_FN_C/2)X_{tu}m_1m_2s_{4\mathrm{\Delta }}}{(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)(tm_{\stackrel{~}{q}_t}^2)}}\right)\{[2\mathrm{\Delta }_um_2^2+2\mathrm{\Delta }_um_1^2+sm_2^2sm_1^2`$ (E10)
$`ss2\mathrm{\Delta }_us_4+ss_42\mathrm{\Delta }_ut_2+2m_2^2t_22m_1^2t_2+st_2+2s_4t_2+2t_2t_2su_2]\widehat{I}({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}})`$
$`[2\mathrm{\Delta }_u]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)+2\widehat{I}(1)\},`$
$`\widehat{M}_{14}^q`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_u}{u_2^2}}\right)\{[4(u_2\mathrm{\Delta }_u)(\mathrm{\Delta }_um_2^2+m_1^2)(\mathrm{\Delta }_um_2^2+m_1^2+u_2)]\widehat{I}\left({\displaystyle \frac{1}{u^{}s_{3\mathrm{\Delta }}}}\right)`$ (E17)
$`+2[6\mathrm{\Delta }_u+4m_1^24m_2^2+4u_2]\widehat{I}\left({\displaystyle \frac{s_{3\mathrm{\Delta }}}{u^{}}}\right)+4\widehat{I}\left({\displaystyle \frac{s_{3\mathrm{\Delta }}^2}{u^{}}}\right)`$
$`+[2(\mathrm{\Delta }_u+m_1^23m_2^2)(\mathrm{\Delta }_u+u_2)`$
$`+(\mathrm{\Delta }_u+m_1^2m_2^2+u_2)(2\mathrm{\Delta }_u+m_1^2m_2^2+ss_4+t_2+u_2)`$
$`+(\mathrm{\Delta }_u+m_1^2m_2^2)(2\mathrm{\Delta }_u+m_1^2m_2^2+ss_4+t_2+3u_2)]\widehat{I}({\displaystyle \frac{1}{s_{3\mathrm{\Delta }}}})`$
$`+2[6\mathrm{\Delta }_u+4m_1^26m_2^2+ss_4+t_2+4u_2]\widehat{I}(1)+6\widehat{I}(s_{3\mathrm{\Delta }})`$
$`+2[\mathrm{\Delta }_u+m_1^23m_2^2]\widehat{I}\left({\displaystyle \frac{u^{}}{s_{3\mathrm{\Delta }}}}\right)+2\widehat{I}(u^{})\}`$
$`+`$ $`\left({\displaystyle \frac{4C_FX_u}{u_2}}\right)\{\mathrm{\Delta }_u[(\mathrm{\Delta }_um_1^2+m_2^2)(2\mathrm{\Delta }_um_1^2+m_2^2s+s_4t_2+u_2)`$ (E24)
$`+(\mathrm{\Delta }_u+m_1^2m_2^2u_2)(2\mathrm{\Delta }_u+m_1^2m_2^2+ss_4+t_2+u_2)]\widehat{I}({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}^2}})`$
$`+[(\mathrm{\Delta }_um_1^2+m_2^2)(2\mathrm{\Delta }_um_1^2+m_2^2s+s_4t_2+u_2)`$
$`+2\mathrm{\Delta }_u(4\mathrm{\Delta }_u3m_1^2+3m_2^2s+s_4t_2+u_2)`$
$`+(\mathrm{\Delta }_u+m_1^2m_2^2u_2)(2\mathrm{\Delta }_u+m_1^2m_2^2+ss_4+t_2+u_2)]\widehat{I}({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}})`$
$`+4\widehat{I}\left({\displaystyle \frac{u_{6\mathrm{\Delta }}}{u^{}}}\right)+2\mathrm{\Delta }_u[3\mathrm{\Delta }_u+2m_1^2+ss_4+t_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}^2}}\right)`$
$`+2[6\mathrm{\Delta }_u2m_1^2s+s_4t_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)6\widehat{I}(1)2\mathrm{\Delta }_u\widehat{I}\left({\displaystyle \frac{u^{}}{u_{6\mathrm{\Delta }}^2}}\right)+2\widehat{I}\left({\displaystyle \frac{u^{}}{u_{6\mathrm{\Delta }}}}\right)\}`$
$`+`$ $`\left({\displaystyle \frac{4C_FX_u}{u_2^2}}\right)\{\mathrm{\Delta }_u[(\mathrm{\Delta }_um_1^2+m_2^2)(2\mathrm{\Delta }_um_1^2+m_2^2s+s_4t_2+u_2)`$ (E29)
$`+(\mathrm{\Delta }_u+m_1^2m_2^2u_2)(2\mathrm{\Delta }_u+m_1^2m_2^2+ss_4+t_2+u_2)]\widehat{I}({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}})`$
$`+2[6\mathrm{\Delta }_u+3m_1^23m_2^2+ss_4+t_2u_2]\widehat{I}\left({\displaystyle \frac{u_{6\mathrm{\Delta }}}{u^{}}}\right)+4\widehat{I}\left({\displaystyle \frac{u_{6\mathrm{\Delta }}^2}{u^{}}}\right)`$
$`+2\mathrm{\Delta }_u[3\mathrm{\Delta }_u+2m_1^2+ss_4+t_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)+2[6\mathrm{\Delta }_u2m_1^2s+s_4t_2]\widehat{I}(1)`$
$`6\widehat{I}(u_{6\mathrm{\Delta }})2\mathrm{\Delta }_u\widehat{I}\left({\displaystyle \frac{u^{}}{u_{6\mathrm{\Delta }}}}\right)+2\widehat{I}(u^{})\},`$
$`\widehat{M}_{15}^q`$ $`=`$ $`\left({\displaystyle \frac{16(C_FN_C/2)X_{tu}m_1m_2s_{4\mathrm{\Delta }}}{s(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\{[(\mathrm{\Delta }_u+t_2)(\mathrm{\Delta }_u+s+t_2)]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)`$ (E31)
$`+[\mathrm{\Delta }_ust_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)\widehat{I}(1)+\widehat{I}\left({\displaystyle \frac{u_{6\mathrm{\Delta }}}{u^{}}}\right)\},`$
$`\widehat{M}_{16}^q`$ $`=`$ $`\left({\displaystyle \frac{8(C_FN_C/2)X_u}{su_2}}\right)\{[(\mathrm{\Delta }_u+m_1^2m_2^2)ss_4+(\mathrm{\Delta }_um_1^2+m_2^2)s(\mathrm{\Delta }_u+u_2)`$ (E38)
$`+2(\mathrm{\Delta }_u+m_1^2m_2^2)(\mathrm{\Delta }_u+u_2)(\mathrm{\Delta }_u+s+t_2+u_2)+u_2(\mathrm{\Delta }_u+u_2)(\mathrm{\Delta }_u+s+t_2+u_2)`$
$`+(\mathrm{\Delta }_um_1^2+m_2^2)(s+u_2)(\mathrm{\Delta }_u+s+t_2+u_2)]\widehat{I}({\displaystyle \frac{1}{u^{}s_{3\mathrm{\Delta }}}})`$
$`+2[3\mathrm{\Delta }_u+m_1^2m_2^2+t_2+2u_2]\widehat{I}\left({\displaystyle \frac{s_{3\mathrm{\Delta }}}{u^{}}}\right)+2\widehat{I}\left({\displaystyle \frac{s_{3\mathrm{\Delta }}^2}{u^{}}}\right)`$
$`+[2\mathrm{\Delta }_u\mathrm{\Delta }_u2\mathrm{\Delta }_um_1^2+2\mathrm{\Delta }_um_2^2+\mathrm{\Delta }_usm_1^2s+m_2^2s2\mathrm{\Delta }_ut_2+m_1^2t_2m_2^2t_23\mathrm{\Delta }_uu_2`$
$`+m_1^2u_2m_2^2u_2+su_2+2t_2u_2+2u_2u_2]\widehat{I}({\displaystyle \frac{1}{s_{3\mathrm{\Delta }}}})`$
$`+[4\mathrm{\Delta }_u+2m_1^22m_2^2s+2t_2+3u_2]\widehat{I}(1)+2\widehat{I}(s_{3\mathrm{\Delta }})+[t_2+u_2]\widehat{I}\left({\displaystyle \frac{u^{}}{s_{3\mathrm{\Delta }}}}\right)\}`$
$`+`$ $`\left({\displaystyle \frac{8(C_FN_C/2)X_u}{su_2}}\right)\{[ss_4(\mathrm{\Delta }_u+m_1^2m_2^2u_2)+\mathrm{\Delta }_uu_2(\mathrm{\Delta }_ust_2)`$ (E45)
$`\mathrm{\Delta }_us(\mathrm{\Delta }_um_1^2+m_2^2+u_2)+2\mathrm{\Delta }_u(\mathrm{\Delta }_u+s+t_2)(\mathrm{\Delta }_um_1^2+m_2^2+u_2)`$
$`+(s+u_2)(\mathrm{\Delta }_u+s+t_2)(\mathrm{\Delta }_u+u_1)]\widehat{I}({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}})+2[3\mathrm{\Delta }_u+m_1^2m_2^2+t_2u_2]\widehat{I}({\displaystyle \frac{u_{6\mathrm{\Delta }}}{u^{}}})`$
$`+2\widehat{I}\left({\displaystyle \frac{u_{6\mathrm{\Delta }}^2}{u^{}}}\right)+[4\mathrm{\Delta }_u\mathrm{\Delta }_u2\mathrm{\Delta }_um_2^2+\mathrm{\Delta }_us+2\mathrm{\Delta }_um_1^2m_1^2s+m_2^2s+s^2`$
$`ss_4+2\mathrm{\Delta }_ut_2m_1^2t_2+m_2^2t_2+st_23\mathrm{\Delta }_uu_2+2su_2+2t_2u_2]\widehat{I}({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}})`$
$`+[8\mathrm{\Delta }_u2m_1^2+2m_2^2s2t_2+3u_2]\widehat{I}(1)4\widehat{I}(u_{6\mathrm{\Delta }})`$
$`+[2\mathrm{\Delta }_u+s+t_2]\widehat{I}\left({\displaystyle \frac{u^{}}{u_{6\mathrm{\Delta }}}}\right)+2\widehat{I}(u^{})\},`$
$`\widehat{M}_{17}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_u}{u_2(s+u_2)}}\right)\{[(\mathrm{\Delta }_um_1^2+m_2^2)ss_4+s_4u_2(\mathrm{\Delta }_uu_2)+2(\mathrm{\Delta }_um_1^2+m_2^2)s_4(\mathrm{\Delta }_u+u_2)`$ (E51)
$`+(\mathrm{\Delta }_u+m_1^2m_2^2)(\mathrm{\Delta }_u+u_2)(s+u_2)`$
$`+(\mathrm{\Delta }_u+m_1^2m_2^2)(s+u_2)(\mathrm{\Delta }_u+s+t_2+u_2)]\widehat{I}({\displaystyle \frac{1}{u^{}s_{3\mathrm{\Delta }}}})`$
$`+2[s+u_2s_4]\widehat{I}\left({\displaystyle \frac{s_{3\mathrm{\Delta }}}{u^{}}}\right)+[2\mathrm{\Delta }_um_2^23\mathrm{\Delta }_us+3m_1^2s3m_2^2s+2\mathrm{\Delta }_us_42m_1^2s_4`$
$`+2m_2^2s_4+m_1^2t_2m_2^2t_2\mathrm{\Delta }_uu_2+2m_1^2u_24m_2^2u_2s_4u_2t_2u_2u_2u_2]\widehat{I}({\displaystyle \frac{1}{s_{3\mathrm{\Delta }}}})`$
$`+[2m_2^2+3s2s_4+u_2]\widehat{I}(1)+[2m_2^2t_2u_2]\widehat{I}\left({\displaystyle \frac{u^{}}{s_{3\mathrm{\Delta }}}}\right)\}`$
$`+`$ $`\left({\displaystyle \frac{4N_CX_u}{u_2(s+u_2)}}\right)\{[2s_4\mathrm{\Delta }_u(\mathrm{\Delta }_u+m_1^2m_2^2u_2)+\mathrm{\Delta }_us_4u_2+ss_4(\mathrm{\Delta }_um_1^2+m_2^2+u_2)`$ (E56)
$`+\mathrm{\Delta }_u(\mathrm{\Delta }_u+u_1)(s+u_2)+(\mathrm{\Delta }_ust_2)(\mathrm{\Delta }_u+u_1)(s+u_2)]\widehat{I}({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}})`$
$`+2[s+u_2s_4]\widehat{I}\left({\displaystyle \frac{u_{6\mathrm{\Delta }}}{u^{}}}\right)+[2\mathrm{\Delta }_um_2^2+\mathrm{\Delta }_us+m_1^2sm_2^2sss2\mathrm{\Delta }_us_4+ss_4`$
$`+m_1^2t_2m_2^2t_2st_2+3\mathrm{\Delta }_uu_22su_22t_2u_2]\widehat{I}({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}})`$
$`+[2m_2^2s+2s_43u_2]\widehat{I}(1)[s+t_2]\widehat{I}\left({\displaystyle \frac{u^{}}{u_{6\mathrm{\Delta }}}}\right)),`$
$`\widehat{M}_{18}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2}{(tm_{\stackrel{~}{q}_t}^2)(s+u_2)}}\right)\times `$ (E59)
$`\{[\mathrm{\Delta }_us+ss+2\mathrm{\Delta }_us_4ss_4+st_22s_4t_2\mathrm{\Delta }_uu_2+su_2+t_2u_2]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)`$
$`+[2\mathrm{\Delta }_u+st_2]\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)2\widehat{I}(1)\},`$
$`\widehat{M}_{22}^q`$ $`=`$ $`\left({\displaystyle \frac{8C_FX_tt_2}{(tm_{\stackrel{~}{q}_t}^2)(tm_{\stackrel{~}{q}_t}^2)}}\right)\left\{\widehat{I}(1){\displaystyle \frac{2\pi (s_4+m_1^2)}{s_4}}\right\},`$ (E60)
$`\widehat{M}_{23}^q`$ $`=`$ $`\left({\displaystyle \frac{4(C_FN_C/2)X_tt_2s_{4\mathrm{\Delta }}}{(tm_{\stackrel{~}{q}_t}^2)(tm_{\stackrel{~}{q}_t}^2)(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\left\{2[m_1^2+m_2^2+s+t_2+u_2]\widehat{I}(1)\right\},`$ (E61)
$`\widehat{M}_{24}^q`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_{tu}m_1m_2}{u_2(tm_{\stackrel{~}{q}_t}^2)}}\right)\{4\widehat{I}\left({\displaystyle \frac{s_{3\mathrm{\Delta }}}{u^{}}}\right)+6\widehat{I}(1)+2\widehat{I}\left({\displaystyle \frac{u^{}}{s_{3\mathrm{\Delta }}}}\right)`$ (E66)
$`+2[3\mathrm{\Delta }_u+m_1^2m_2^2+2ss_4+3t_2+3u_2]\widehat{I}\left({\displaystyle \frac{1}{s_{3\mathrm{\Delta }}}}\right)`$
$`+[4\mathrm{\Delta }_u^22\mathrm{\Delta }_um_1^2+2\mathrm{\Delta }_um_2^24\mathrm{\Delta }_us+m_1^2sm_2^2s+s^2+2\mathrm{\Delta }_us_4ss_46\mathrm{\Delta }_ut_2+2m_1^2t_2`$
$`2m_2^2t_2+3st_22s_4t_2+2t_2^28\mathrm{\Delta }_uu_2+2m_1^2u_2+4u_2^22m_2^2u_2+5su_2`$
$`2s_4u_2+6t_2u_2]\widehat{I}\left({\displaystyle \frac{1}{u^{}s_{3\mathrm{\Delta }}}}\right)\}`$
$`+`$ $`\left({\displaystyle \frac{4C_FX_{tu}m_1m_2}{u_2(tm_{\stackrel{~}{q}_t}^2)}}\right)\{2\mathrm{\Delta }_u\widehat{I}\left({\displaystyle \frac{1}{u_{6\mathrm{\Delta }}}}\right)2\widehat{I}(1)+4\widehat{I}\left({\displaystyle \frac{u_{6\mathrm{\Delta }}}{u^{}}}\right)`$ (E69)
$`+[4\mathrm{\Delta }_u^22\mathrm{\Delta }_um_1^2+2\mathrm{\Delta }_um_2^24\mathrm{\Delta }_us+m_1^2sm_2^2s+s^2+2\mathrm{\Delta }_us_4ss_46\mathrm{\Delta }_ut_2`$
$`+2m_1^2t_22m_2^2t_2+3st_22s_4t_2+2t_2t_2+su_2]\widehat{I}\left({\displaystyle \frac{1}{u^{}u_{6\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{25}^q`$ $`=`$ $`\left({\displaystyle \frac{8(C_FN_C/2)X_ts_{4\mathrm{\Delta }}}{s(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)(tm_{\stackrel{~}{q}_t}^2)}}\right)\{[(s+t_2)(s_4+t_2)]\widehat{I}(1)`$ (E71)
$`+[ss_4+st_22s_4t_2s_4u_2+t_2u_2]\widehat{I}\left({\displaystyle \frac{t^{}}{u^{}}}\right)\},`$
$`\widehat{M}_{26}^q`$ $`=`$ $`\left({\displaystyle \frac{16(C_FN_C/2)X_{tu}m_1m_2}{s(tm_{\stackrel{~}{q}_t}^2)}}\right)\{[(\mathrm{\Delta }_u+t_2+u_2)(\mathrm{\Delta }_u+s+t_2+u_2)]\widehat{I}\left({\displaystyle \frac{1}{u^{}s_{3\mathrm{\Delta }}}}\right)`$ (E73)
$`+\widehat{I}(1)+[\mathrm{\Delta }_u+t_2+u_2]\widehat{I}\left({\displaystyle \frac{1}{s_{3\mathrm{\Delta }}}}\right)+\widehat{I}\left({\displaystyle \frac{s_{3\mathrm{\Delta }}}{u^{}}}\right)\},`$
$`\widehat{M}_{27}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2}{(s+u_2)(tm_{\stackrel{~}{q}_t}^2)}}\right)\{[2s2s_4+t_2+u_2]\widehat{I}\left({\displaystyle \frac{1}{s_{3\mathrm{\Delta }}}}\right)`$ (E75)
$`+[ss_42s_4(\mathrm{\Delta }_u+s+t_2+u_2)+(s+u_2)(\mathrm{\Delta }_u+s+t_2+u_2)]\widehat{I}\left({\displaystyle \frac{1}{u^{}s_{3\mathrm{\Delta }}}}\right)\},`$
$`\widehat{M}_{28}^q`$ $`=`$ $`\left\{{\displaystyle \frac{8N_CX_tt_2}{(tm_{\stackrel{~}{q}_t}^2)(tm_{\stackrel{~}{q}_t}^2)}}\right)(s_4^2\widehat{I}\left({\displaystyle \frac{1}{u^{}u_7}}\right)+(m_1^2s_4)\widehat{I}\left({\displaystyle \frac{1}{u_7}}\right)\widehat{I}(1)\}.`$ (E76)
A squark width $`\mathrm{\Gamma }`$ is included to regularize a possible on-shell squark pole. As discussed in the text, $`\mathrm{\Gamma }`$ serves only to regulate the divergence of these interference terms, and it is taken to be very small compared to all physical masses and momenta.
The finite pieces of the quark emission terms are evaluated directly from Eq. (47) with the angular integrals integrated numerically. They may be expressed as
$`\overline{|^q|}^2`$ $`=`$ $`\left({\displaystyle \frac{\pi ^2\alpha _S\widehat{\alpha }_S}{3}}\right){\displaystyle \underset{i=\mathrm{1..8}}{}}{\displaystyle \underset{j=i\mathrm{..8}}{}}M_{ij}^q,`$ (E77)
with
$`M_{11}^q=M_{12}^q=M_{13}^q=M_{15}^q=M_{22}^q=M_{23}^q=M_{25}^q=M_{26}^q=M_{28}^q`$ $`=`$ $`0,`$ (E78)
$`M_{14}^q`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_tu_7s_{3\mathrm{\Delta }}}{t_2^2(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_{7\mathrm{\Delta }}}}\right)\{2s_{32}^2ss_{32}2s_4s_{32}+2t_2s_{32}2s_{32}u^{}u6s_{32}+t^{}s_5t^{}m_1^2`$ (E80)
$`3t^{}m_2^2+t^{}u_2s_4t_2t_2u^{}\}`$
$`+`$ $`\left({\displaystyle \frac{4C_FX_tu_7s_{3\mathrm{\Delta }}}{t_2(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_{7\mathrm{\Delta }}^2}}\right)\{2s_{32}^2ss_{32}2s_4s_{32}+2t_2s_{32}2s_{32}u^{}u6s_{32}+t^{}s_5t^{}m_1^2`$ (E82)
$`3t^{}m_2^2+t^{}u_2s_4t_2t_2u^{}\},`$
$`M_{16}^q`$ $`=`$ $`\left({\displaystyle \frac{8(C_FN_C/2)X_ts_{3\mathrm{\Delta }}}{s(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_{7\mathrm{\Delta }}t_2}}\right)\{ss_{32}s_4+s_{32}u^{}u_6+ss_{32}u_7+s_{32}t^{}u_7`$ (E84)
$`+2s_{32}u^{}u_7+t_2u^{}u_7t^{}u_2u_7\},`$
$`M_{17}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_ts_{3\mathrm{\Delta }}}{(s+t_2)(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_{7\mathrm{\Delta }}t_2}}\right)\{ss_{32}s_4+s_{32}u^{}u_62s_4s_{32}u_7t^{}u_7m_2^2`$ (E86)
$`m_1^2t^{}u_7s_{32}t^{}u_7+s_5t^{}u_7s_4t_2u_7s_{32}u_6u_7\}`$
$`+`$ $`\left({\displaystyle \frac{4N_CX_ts_{3\mathrm{\Delta }}}{u_6(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_{7\mathrm{\Delta }}(s+t_2)}}\right)\{ss_{32}s_4+s_{32}u^{}u_62s_4s_{32}u_7t^{}u_7m_2^2`$ (E88)
$`m_1^2t^{}u_7s_{32}t^{}u_7+s_5t^{}u_7s_4t_2u_7s_{32}u_6u_7\}`$
$`M_{18}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2}{u_6(um_{\stackrel{~}{q}_u}^2)u_{7\mathrm{\Delta }}(s+t_2)}}\right)\left\{ss_42st^{}2s_4u^{}2t^{}u^{}u^{}u_6+t^{}u_7\right\},`$ (E89)
$`M_{24}^q`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_{tu}m_1m_2s_{3\mathrm{\Delta }}}{u_{7\mathrm{\Delta }}(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)(um_{\stackrel{~}{q}_u}^2)t_2}}\right)\{ss_{32}+ss_4+2su^{}2s_{32}u^{}+2s_4u^{}t_2u^{}`$ (E91)
$`+2u_{}^{}{}_{}{}^{2}+t^{}u_2+u^{}u_6t^{}u_7\},`$
$`M_{27}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}m_1m_2s_{3\mathrm{\Delta }}}{u_6(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)(um_{\stackrel{~}{q}_u}^2)(s+t_2)}}\right)\left\{ss_42s_4u^{}u^{}u_6+t^{}u_7\right\},`$ (E92)
$`M_{33}^q`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_us_4u_2}{(um_{\stackrel{~}{q}_u}^2)^2(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\left\{2m_2^2+2m_1^2s_{32}+s_4s_5u^{}+u_2u_7\right\},`$ (E93)
$`M_{34}^q`$ $`=`$ $`\left({\displaystyle \frac{8(C_FN_C/2)X_{tu}s_{4\mathrm{\Delta }}s_{3\mathrm{\Delta }}m_1m_2u^{}}{(um_{\stackrel{~}{q}_u}^2)^2(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_{7\mathrm{\Delta }}(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\left\{s_5u^{}2m_2^22m_1^2\right\},`$ (E94)
$`M_{35}^q`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_us_4}{(um_{\stackrel{~}{q}_u}^2)(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)s}}\right)\{3m_2^2s+m_1^2sss_{32}ss_5+t_2u^{}+t^{}u_2`$ (E96)
$`2t_2u_2+2u^{}u_22u_2^2+u_2u_6+t_2u_7+2u_2u_7\},`$
$`M_{36}^q`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_{tu}s_{4\mathrm{\Delta }}s_{3\mathrm{\Delta }}m_1m_2}{(um_{\stackrel{~}{q}_u}^2)(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)s}}\right)\{ss_{32}ss_4+2t^{}u^{}t_2u^{}+2u_{}^{}{}_{}{}^{2}t^{}u_2`$ (E98)
$`2u^{}u_2+u^{}u_6+t^{}u_7+2u^{}u_7\},`$
$`M_{37}^q`$ $`=`$ $`\left({\displaystyle \frac{2N_CX_{tu}s_{4\mathrm{\Delta }}s_{3\mathrm{\Delta }}m_1m_2}{(um_{\stackrel{~}{q}_u}^2)(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_6}}\right)\{ss_{32}ss_42m_2^2u^{}6m_1^2u^{}`$ (E100)
$`+2su^{}2s_4u^{}+2s_5u^{}+t_2u^{}t^{}u_2u^{}u_6+t^{}u_7+2u^{}u_7\},`$
$`M_{38}^q`$ $`=`$ $`\left({\displaystyle \frac{2N_CX_us_{4\mathrm{\Delta }}u_2}{(um_{\stackrel{~}{q}_u}^2)^2(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_6}}\right)\{2m_2^2s_46m_1^2s_4+ss_42s_4^2+2s_4s_5`$ (E102)
$`m_2^2t^{}3m_1^2t^{}2s_4t^{}+s_5t^{}+s_4t_2s_{32}u_6u^{}u_6+2s_4u_7+t^{}u_7\},`$
$`M_{44}^q`$ $`=`$ $`\left({\displaystyle \frac{4C_FX_ts_{32}u_7}{u_{7\mathrm{\Delta }}^2(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\left\{2m_2^2+2m_1^2+s_{32}s_4s_5u^{}u_2+u_7\right\},`$ (E103)
$`M_{45}^q`$ $`=`$ $`\left({\displaystyle \frac{4(C_FN_C/2)X_{tu}s_{4\mathrm{\Delta }}s_{3\mathrm{\Delta }}m_1m_2}{u_{7\mathrm{\Delta }}(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)s(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\{ss_{32}+ss_4+2t^{}u^{}+t_2u^{}+2u_{}^{}{}_{}{}^{2}+t^{}u_2`$ (E105)
$`+2u^{}u_2u^{}u_6t^{}u_72u^{}u_7\},`$
$`M_{46}^q`$ $`=`$ $`\left({\displaystyle \frac{4(C_FN_C/2)X_ts_{32}}{u_{7\mathrm{\Delta }}(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)s}}\right)\{sm_2^2+3sm_1^2ss_4ss_5+u^{}u_6+u_2u_6`$ (E107)
$`+t^{}u_7+t_2u_7+2u^{}u_7+2u_2u_72u_6u_72u_7^2\},`$
$`M_{47}^q`$ $`=`$ $`\left({\displaystyle \frac{2N_CX_ts_{32}}{u_{7\mathrm{\Delta }}(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_6}}\right)\{sm_2^2+3sm_1^2ss_4ss_5+2u_7m_2^2+u^{}u_6+u_2u_6`$ (E109)
$`+6u_7m_1^2+2su_72s_4u_72s_5u_7t^{}u_7t_2u_7+2u_7^2\},`$
$`M_{48}^q`$ $`=`$ $`\left({\displaystyle \frac{2N_CX_{tu}s_{3\mathrm{\Delta }}m_1m_2}{u_{7\mathrm{\Delta }}(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_6(um_{\stackrel{~}{q}_u}^2)}}\right)\{ss_{32}ss_4+2u^{}m_2^2+6u^{}m_1^22s_4u^{}`$ (E111)
$`2s_5u^{}2t^{}u^{}t_2u^{}t^{}u_2+u^{}u_6+t^{}u_7+2u^{}u_7\},`$
$`M_{55}^q`$ $`=`$ $`\left({\displaystyle \frac{8C_FX_us_4t_2}{s(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right),`$ (E112)
$`M_{56}^q`$ $`=`$ $`\left({\displaystyle \frac{16C_FX_{tu}s_{4\mathrm{\Delta }}s_{3\mathrm{\Delta }}m_1m_2t^{}}{s(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right),`$ (E113)
$`M_{57}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_{tu}s_{4\mathrm{\Delta }}s_{3\mathrm{\Delta }}m_1m_2}{s(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)u_6(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\left\{ss_4+2st^{}+2su^{}+u^{}u_6+t^{}u_7+2u^{}u_7\right\},`$ (E114)
$`M_{58}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_us_{4\mathrm{\Delta }}}{s(um_{\stackrel{~}{q}_u}^2)u_6(s_{4\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\{ss_4m_2^2ss_4m_1^2+ss_4s_5ss_4u_2s_4u_2u_6`$ (E116)
$`+u^{}u_2u_6s_4t_2u_72s_4u_2u_7t^{}u_2u_7\},`$
$`M_{66}^q`$ $`=`$ $`\left({\displaystyle \frac{8C_FX_ts_{32}su_6}{s^2(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right),`$ (E117)
$`M_{67}^q`$ $`=`$ $`\left({\displaystyle \frac{2N_CX_t}{su_6(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\left\{4sm_1^2+4su_6+4su_7+4u_6u_7+4u_7^2\right\},`$ (E118)
$`M_{68}^q`$ $`=`$ $`\left({\displaystyle \frac{2N_CX_{tu}s_{3\mathrm{\Delta }}m_1m_2}{su_6(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)(um_{\stackrel{~}{q}_u}^2)}}\right)\left\{ss_4+u^{}u_6+t^{}u_7+2u^{}u_7\right\},`$ (E119)
$`M_{77}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_ts_{32}}{u_6^2(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)}}\right)\left\{4m_1^2s+2su_64m_1^2u_7\right\},`$ (E120)
$`M_{78}^q`$ $`=`$ $`\left({\displaystyle \frac{16N_CX_{tu}s_{3\mathrm{\Delta }}m_1m_2u^{}(2m_1^2+u_6)}{u_6^2(s_{3\mathrm{\Delta }}^2+m_{\stackrel{~}{q}_t}^2\mathrm{\Gamma }^2)(um_{\stackrel{~}{q}_u}^2)}}\right),`$ (E121)
$`M_{88}^q`$ $`=`$ $`\left({\displaystyle \frac{4N_CX_uu_2}{u_6^2(um_{\stackrel{~}{q}_u}^2)^2}}\right)\left\{4m_1^2s_4+4m_1^2t^{}2t^{}u_6\right\}.`$ (E122)
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# 1 Introduction
## 1 Introduction
Recent atmospheric and solar neutrino experiments suggest that neutrinos have tiny masses compared to quarks and charged leptons and the MNS mixing matrix is different from the CKM matrix in the quark sector. The tiny masses for neutrinos are expected to be less than $`𝒪(1)`$eV. It is well known that there are mainly two mechanisms to obtain tiny neutrino masses. One is the seesaw mechanism, which requires the existence of right-handed Majorana neutrinos with $`10^{10}10^{15}`$GeV masses. The other mechanism is to generate tiny masses by higher order loop effects by the extension of the standard model(SM). In both senarios light neutrinos are Majorana. In this talk we investigate the second possibility to interpret the tiny neutrino masses and mixing patterns suggested by experiments within the framework of three flavors.
## 2 Zee Model
A model with extension of the Higgs sector by adding one more doublet in $`SU(2)_L`$ and a charged singlet scalar with lepton number 2 was proposed by Zee almost two decades ago. The Zee model is the minimal extension of the SM and the one loop effects by extended Higgs sector can induce small Majorana neutrino masses. In the Zee model the neutrino masses are given by
$$m_{ll^{}}=f_{ll^{}}(m_l^2m_l^2)\frac{\mu v_u}{v_d}F(M_1^2,M_2^2),$$
(1)
where $`M_{1,2}`$ are the masses of the two physical charged scalars. This type of neutrino matrix has the following characteristic features;
1. The diagonal elements are zero due to antisymmetric properties $`f_{ll^{}}=f_{l^{}l}`$.
2. There is no $`CP`$ violation because the three phases in the mass matrix are absorbed by the three Majorana neutrino fields. However, if we take into account of two loop effects there appear tiny non-zero masses at diagonal parts and it is possible to obtain $`CP`$ violation in this model.
3. If we assume naively the relations $`f_{e\mu }f_{e\tau }f_{\mu \tau }`$, it is expected that the relation $`m_{\mu \tau }m_{e\tau }m_{e\mu }`$ due to $`m_{ij}f_{ij}(m_i^2m_j^2)`$. However the mass and mixing matrices are inconsistent with the result of CHOOZ’s experiments.
In this analysis we call the mass matrix given by Eq.(1) as Zee-type neutrino mass matrix and we investigate the consistent solution of Zee type mass matrix with experiments generally by relaxing the relation $`f_{e\mu }f_{e\tau }f_{\mu \tau }`$.
## 3 Bi-maximal MNS Matrix
The mass matrix (1) is symmetric and diagonalization is done by the MNS mixing matrix
$$U_{MNS}=\left(\begin{array}{ccc}c_1c_3& s_1c_3& s_3\\ s_1c_2c_1s_2s_3& c_1c_2s_1s_2s_3& s_2c_3\\ s_1s_2c_1c_2s_3& c_1s_2s_1c_2s_3& c_2c_3\end{array}\right).$$
(2)
We obtain the following conditions due to the property of Zee-type mass matrix;
$`(1,1)`$ $`m_1c_1^2c_3^2+m_2s_1^2c_3^2+m_3s_3^2=0,`$
$`(2,2)`$ $`m_1(s_1c_2+c_1s_2s_3)^2+m_2(c_1c_2s_1s_2s_3)^2+m_3s_2^2c_3^2=0,`$
$`(3,3)`$ $`m_1(s_1s_2c_1c_2s_3)^2+m_2(c_1s_2+s_1c_2s_3)^2+m_3c_2^2c_3^2=0,`$
where $`m_i(i=1,2,3)`$ is the eigenvalues of Eq.(1). Then the relations
$`m_2={\displaystyle \frac{\mathrm{cos}^2\theta _1\mathrm{tan}^2\theta _3}{\mathrm{sin}^2\theta _1\mathrm{tan}^2\theta _3}}m_1,`$ $`m_3=m_1m_2`$
$`\mathrm{cos}2\theta _1\mathrm{cos}2\theta _2\mathrm{cos}2\theta _3`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sin}2\theta _1\mathrm{sin}2\theta _2(3\mathrm{cos}^2\theta _32)\mathrm{sin}\theta _3`$ (3)
are derived. It is noted that the hierarchical mass solution like $`m_3m_2m_1`$ is prohibited in this model by the traceless condition.
Here we input the requisite condition $`\theta _2\pi /4`$ to obtain large mixing suggested by the atmospheric neutrino experiments. From Eq.(3) the possible solutions are $`\theta _30`$, $`\theta _10`$ or $`\theta _3\mathrm{arctan}\sqrt{1/2}`$. However the latter two cases except for $`\theta _1\theta _30`$ require $`|U_{e3}|>0.22`$ and this solution is contradicting with the CHOOZ experiment. Then it is enough to investigate the case of $`\theta _30`$. There are two hierarchical mass squared differences $`\mathrm{\Delta }m_{atm}^2\mathrm{\Delta }m_{solar}`$ and in this model the unique solution is given for $`\theta _30`$ in Eq.(2).
The degenerate solutions with two masses are given for (1)$`\theta _10`$, (2) $`\theta _1\mathrm{arctan}1/\sqrt{2}`$ and (3)$`\theta _1\pi /4`$. The former two cases are inconsistent with the experimental data. The consistent solution is given only for the last case, which implies the solution $`|m_1||m_2||m_3|,\mathrm{\Delta }m_{13}^2\mathrm{\Delta }m_{23}^2=\mathrm{\Delta }m_{atm}^2,\mathrm{\Delta }m_{12}^2=\mathrm{\Delta }m_{}^2`$. The mixing matrix is obtained as
$$U^{Zee}=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}& \frac{1}{\sqrt{2}}& 0\\ \frac{1}{2}& \frac{1}{2}& \frac{1}{\sqrt{2}}\\ \frac{1}{2}& \frac{1}{2}& \frac{1}{\sqrt{2}}\end{array}\right).$$
(4)
This corresponds to the MSW or vacuum large angle solution in solar neutrino experiments. Only the third case($`\theta _1\pi /4,\theta _2\pi /4,\theta _30`$) can give a consistent solution and the MNS mixing matrix is bi-maximal. The neutrino masses are $`|m_1||m_2|\sqrt{\mathrm{\Delta }m_{atm}^2},|m_3|\frac{\mathrm{\Delta }m_{}^2}{2\sqrt{\mathrm{\Delta }m_{atm}^2}}`$ and an anti-hierarchical mass structure is obtained. The relation $`m_1m_2`$ implies pseudo-Dirac solution. For the parameters in the Zee model the consistent solution with experiments suggests a hierarchical relation as $`f_{e\mu }f_{e\tau }f_{\mu \tau }`$.
## 4 Summary
The unique solution of the model is the bi-maximal one, which corresponds to MSW or vacuum large mixing angle solution in solar neutrino oscillation. The neutrino masses are $`|m_1||m_2|0.020.08\mathrm{eV}(|m_3||m_1||m_2|)`$. The neutrinoless double $`\beta `$ decay is induced by the (1,1) element in the mass matrix and $`m_\nu =|_iU_{ei}^2m_{\nu i}|=|(1,1)|0`$ prohibits the neutrinoless double $`\beta `$ decay.
The solution implies the relation $`f_{e\mu }f_{e\tau }f_{\mu \tau }`$ in the Zee model. The conservation of the new lepton number $`L_{new}L_eL_\mu L_\tau `$ requires $`f_{\mu \tau }=0`$ means $`m_{\mu \tau }=0`$. The realistic model is suggested in the extended Zee model with $`L_{new}=2`$ doubly charged Higgs scalar field and 2-loop effects derive non-zero $`m_{\mu \tau }`$. Also the two loop effects in the original Zee model are studied.
Also the models which give the Zee type mass matrix are discussed in the viewpoint of $`R`$-parity violating SUSY model or gauge mediated SUSY model. Also three active and one sterile neutrino model is discussed in Ref.. The phenomenology of $`h^+`$ is also analyzed.
The bi-maximal solution in Zee type mass matrix requires two degenerate heavy masses and one light neutrino. Solar neutrino oscillation experiment can resolve which mixing matrix is viable.
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# 𝛼_𝑆 EVOLUTION FROM 35 GEV TO 202 GEV AND FLAVOUR INDEPENDENCE
## 1 Motivation
Three fundamental properties follow from QCD: (i) the scale dependence of the renormalised coupling strength, (ii) the flavour independence of the coupling apart from effects due to finite quark masses, and (iii) the scale dependence of the renormalised quark masses. It constitutes a significant experimental test of QCD if the strong interaction obeys these properties.
The large range of energy covered by e<sup>+</sup>e<sup>-</sup> colliders makes an investigation of the energy evolution of $`\alpha _S`$ possible. Although the expected running of the coupling is more pronounced towards lower centre-of-mass energies, the uniform analyses at LEP provide significant QCD tests up to the highest energies accessible.
## 2 $`\alpha _S(200\mathrm{GeV})`$
The excellent performance of the LEP collider in 1999 provided data at $`\sqrt{s}=192`$, 196, 200, and 202 GeV comprising together about 220 pb<sup>-1</sup> per experiment. The results at each of the four $`\sqrt{s}`$ will be combined to derive the coupling strength at $`\sqrt{s}=198`$ GeV which is the luminosity weighted average energy.
Even though the selection of hadronic final states is trivial, background contributions from essentially W pair production and large initial state radiation have to be rejected. W pair events are excluded using the wide separation of the four jets and requiring di-jet masses different from the W mass. The demand that the effective centre-of-mass energy of the hadronic final state is within 10-20% of the nominal value is applied against initial state radiation events. These criteria select more than 75% of the real hadronic final states and reduce the contribution from other processes to less than 20%.
From the selected events the observables thrust ($`T`$), heavy jet mass ($`M_H`$), total and wide jet broadening ($`B_T`$, $`B_W`$), $`C`$ parameter ($`C`$), and the value of the jet resolution parameter at the transition from 3 to 2 jets using the Durham jet algorithm ($`y_3`$) are measured. The strong coupling constant is found from fitting the QCD prediction, convolved with the hadronisation correction, to the distribution of an observable. The QCD prediction is the matched resummed next-to-leading logarithmic approximation (NLLA) with the full second order matrix element ($`𝒪(\alpha _S^2)`$). The hadronisation correction is taken from various Monte Carlo event generators which although they were tuned to describe the data measured at 91 GeV yield a good representation of the data at energies of around 200 GeV.
The left part of Fig. 1 shows the fit results of the LEP experiments $`^\mathrm{?}`$ for $`\alpha _S`$ at $`\sqrt{s}=192`$-196 and 200-202 GeV. Combining the values yields $`\alpha _S(198\mathrm{GeV})=0.109\pm 0.001_{(\mathrm{expt})}\pm 0.005_{(\mathrm{theo})}`$.
## 3 Energy evolution
The LEP experiments contributed a large number of $`\alpha _S`$ determinations for $`\sqrt{s}m_\mathrm{Z}`$ (see $`^\mathrm{?}`$ and references therein). Exploiting initial and hard final state photon radiation, $`\sqrt{s}`$ as low as 30 GeV are accessible and have been investigated by L3 and DELPHI $`^\mathrm{?}`$. These determinations of $`\alpha _S`$ below the Z mass are complemented by results of experiments at lower centre-of-mass energies $`^{\mathrm{?},\mathrm{?}}`$ which, even though already completed since long, re-analysed their data to employ the matched resummed NLLA and $`𝒪(\alpha _S^2)`$ predictions for an $`\alpha _S`$ determination.
The right part of Fig. 1 shows the $`\alpha _S`$ values versus the centre-of-mass energy. The values with their experimental errors are fitted with the 4-loop formula for the scale dependence $`^\mathrm{?}`$ of $`\alpha _S`$. To estimate the theory uncertainty the scale uncertainties of the single measurements are considered to be fully correlated. The fit yields $`\alpha _S(m_\mathrm{Z})=0.1208\pm 0.0006_{(\mathrm{expt})}\pm 0.0048_{(\mathrm{theo})}`$ which is not dominated by the single measurement at 91.2 GeV and which agrees with the world average of $`0.1184\pm 0.0031`$ $`^\mathrm{?}`$.
Recently jet observables for the Durham and Cambridge jet algorithms have been investigated using data of the OPAL and the JADE experiments $`^\mathrm{?}`$. The analysis treated the data of both experiments and estimated the errors of $`\alpha _S`$ in a similar way. Fig. 2 shows the 2-, 3-, 4-, and 5-jet fractions for the Durham jet algorithm at three different $`\sqrt{s}`$. Predictions of several models are overlaid.
The coupling strength has been obtained from fits of the matched resummed NLLA plus second order matrix element to the differential 2-jet rate ($`D_2`$) and the jet multiplicity ($`N`$) of both jet finders. Fig. 3 shows the results with the world average overlaid.
Combining all eight determinations and taking correlations into account yields $`\alpha _S(m_\mathrm{Z})=0.1187\pm 0.0010_{(\mathrm{expt})}^{+0.0032}_{0.0016(\mathrm{theo})}`$ which is in excellent agreement with the world average $`^\mathrm{?}`$ and has a very small total error.
## 4 Flavour independence
Finite quark masses affect the result of $`\alpha _S`$ determinations. In particular bottom quark events at the Z mass yield a 7% lower value of $`\alpha _S`$ if the quark mass effect is neglected $`^\mathrm{?}`$. To account for the mass effect for an inclusive determination which neglected this effect, the value of $`\alpha _S`$ has to be increased by about 1%, which is covered by the typical total error.
Being precisely confirmed for the heavy charm and bottom quarks, the flavour independence has been scarcely tested for the light quarks at high energies. At $`\sqrt{s}0`$ evidence for the flavour independence comes from e.g. isospin invariance and approximate $`SU(3)_{\mathrm{flav}}`$ symmetry. The challenge for an investigation of the flavour independence at $`\sqrt{s}=m_\mathrm{Z}`$ is to separate u, d, and s quark events. Using the leading particle effect $`^\mathrm{?}`$ for K<sup>±</sup>, K$`{}_{S}{}^{}{}_{}{}^{0}`$, and all kinds of charged particles OPAL $`^\mathrm{?}`$ selected events which are enriched differently in u, d, and s quarks and thus allows for a statistical decomposition of the contribution of each of the three light quark flavours. $`\alpha _S`$ is determined from the charged multiplicities, $`N`$, obtained for each quark flavour from the decomposition, using $`N\alpha _S^B\mathrm{exp}(C/\sqrt{\alpha _S})`$ where $`B`$ and $`C`$ are known from QCD calculations $`^\mathrm{?}`$. This yielded the preliminary ratios: $`\alpha _S^\mathrm{u}/\alpha _S^\mathrm{d}=0.88\pm 0.08`$, $`\alpha _S^\mathrm{s}/\alpha _S^\mathrm{d}=0.96\pm 0.06`$, and $`\alpha _S^\mathrm{s}/\alpha _S^\mathrm{u}=1.09\pm 0.06`$ which are consistent with flavour independence at a level of better than 10% and constitute an improvement over the previous OPAL result $`^\mathrm{?}`$.
## 5 Summary
QCD is in a very good shape! The fundamental properties of the theory, i.e. the running of $`\alpha _S`$, its flavour independence and also the running of the renormalised quark masses are observed and confirmed in experimental investigations. At the highest energies of LEP the value of the coupling is determined to be $`\alpha _S(198\mathrm{GeV})=0.109\pm 0.005`$ (prelim.) which agrees with the expected running. The value of the strong coupling constant is now very precisely determined from a combined analysis of OPAL and JADE data covering centre-of-mass energies from $`\sqrt{s}=22`$ through $`189`$ GeV to be $`\alpha _S(m_\mathrm{Z})=0.1187_{0.0019}^{+0.0034}`$. The flavour independence of the coupling at high energies is now confirmed for the light quarks at the level of better than 10%.
## References
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# Black hole evolution by spectral methods
## I Introduction and summary
A major thrust of research in classical general relativity in the past decade has been to devise algorithms to solve Einstein’s equations numerically. Despite advances in our analytic understanding of general relativity, we still do not know what all the features of the theory really are. Numerical solutions will continue to provide fresh insights into the theory as they have in the past, for example, critical behavior in black hole formation and the formation of toroidal black holes .
New urgency has been injected into numerical relativity by the imminent deployment of LIGO. The prime target for LIGO is coalescence of binary neutron star and black hole systems. The waveform is reasonably well predicted by the post-Newtonian approximation when the binary components are at large separation. However, extracting the most important physics requires us to be able to deal with fully non-linear general relativity as the system spirals together and coalesces. Moreover, a number of people believe that there is a significant event rate for the coalescence of massive black hole systems ($`20M_{}`$) . In this case, LIGO is most sensitive to waves emitted from the strong field regime. Indeed, without some theoretical guidance as to what to expect from this regime, it is possible we may miss these events entirely .
However, the goal of developing a general algorithm that can solve Einstein’s equations for two black holes has remained elusive. All attempts to date have been plagued by instabilities. These instabilities are caused by an interplay of three factors: (1) Einstein’s equations are an overdetermined system, with the evolution equations subject to constraints. So if, for example, you choose to solve only the evolution equations, then there can be unstable solutions that are in fact solutions of the evolution equations, but do not satisfy the constraints. Small numerical errors may cause these solutions to appear and swamp the true solution (“constraint-violating modes”). (2) The coordinate freedom inherent in the theory means that it is very easy to impose coordinate conditions that lead to numerical instabilities (“gauge modes”). (3) Experience has shown that the kind of boundary conditions we choose and how we implement them can affect the stability of an algorithm enormously.
Similar instabilities have hampered efforts to solve the related problem of binary neutron stars; only very recently has there been some success in finding stable algorithms. However, black hole evolutions face an additional obstacle that is absent in the case of neutron stars: for neutron stars the gravitational field is everywhere regular, but for black holes one must somehow deal with the physical singularity that lurks inside each hole.
There are two main approaches for handling these singularities. The first is to use gauge conditions (e.g., maximal slicing) that avoid the singularities altogether. Such conditions, however, lead to large gradients in the gravitational field variables near the horizon. These grow exponentially in time and ultimately cannot be resolved by the numerical evolution, causing the code to crash. The alternative approach is to excise the region containing the singularity from the computational domain and evolve only the exterior region. If the excision boundary is placed inside the horizon of the black hole, causality assures us that we do not need to impose a physical boundary condition there.
However, black hole excision is only known to be mathematically well-posed if the evolution equations are hyperbolic with characteristic speeds less than or equal to $`c`$. In this case, the structure of the equations guarantees that even unphysical modes present in the solution (gauge modes, constraint-violating modes) behave causally and cannot propagate out of the horizon. For many representations of general relativity such as the usual ADM formulation, the evolution equations are of no mathematical type for which well-posedness has been proven, so the suitability of these formulations for black hole excision must be determined empirically on a case-by-case basis. It is in part for this reason that much attention has been recently focused on hyperbolic representations of Einstein’s equations .
It is still unclear whether hyperbolic formulations are computationally advantageous. However, it has been shown that the formulation of the evolution equations can affect stability. For example, some instabilities can be eliminated by changing from one formulation of Einstein’s equations to another or by modifying the evolution equations to change the spectrum of unphysical modes .
Nevertheless, it is likely that many instabilities encountered in practice are due to the numerical implementation of the evolution equations and of the boundary conditions. Even for well-understood systems, it is far from guaranteed that any given numerical approximation to the continuum equations will be stable. The well-known Courant instability is a trivial example. Hence it is prudent to explore alternative numerical methods that may offer a shorter path to the goal of long-term stability.
Traditionally, black hole spacetimes have been evolved using finite-difference (FD) methods. Current FD codes for evolving black hole spacetimes with excised horizons are mostly based on a numerical technique known as causal differencing , which allows one to update the fundamental variables in time while avoiding numerical problems associated with superluminal grid speeds. This technique has been used successfully to propagate an excised hole across a grid, even when grid points fall into or emerge from the horizon .
However, causal differencing is complicated because it requires interpolation, and it has to deal with points “missing” in some irregular fashion near the excision boundary. The FD operator that performs the interpolation depends on the shape of the excision boundary. Furthermore, even for a given excision boundary and a given target point, the operator is not unique. Hence one must construct a large number of interpolation operators, each of which involves some arbitrary choice, in order to perform interpolations on the entire grid. It is only by trial and error that one finds operators that result in a stable evolution scheme. For simulations of a single spherical black hole on a Cartesian three-dimensional grid, we have found a case in which changing a single interpolation operator that is used only for a single target point on the grid makes the difference between a stable and an unstable code.
Another limitation of FD methods is the difficulty of imposing boundary conditions at the outer boundary of the calculation. There are two aspects to this problem: First, one must formulate a procedure for handling the boundary, such as imposing an analytic condition (e.g., a Sommerfeld condition) on the fundamental variables, matching to a wave perturbation described by the Zerilli equation , or matching to a characteristic evolution code that propagates the solution out to null infinity . Second, one must construct a FD approximation of either the analytic boundary condition or the matching condition. It can be difficult to find such an approximation that yields a stable evolution.
In this paper we explore an alternative computational strategy: a pseudospectral collocation (PSC) scheme. In our PSC evolution scheme, the solutions to a set of hyperbolic differential equations are approximated as series expansions in a set of orthogonal basis functions (e.g., Chebyshev polynomials) in space, and these coefficients are integrated forward in time using the method of lines. PSC has three important advantages over FD: (1) No ad hoc interpolation operators are required to determine field values at an arbitrary point because the solution provided by PSC is an analytic function given everywhere on the computational domain. (2) Boundary conditions are imposed directly on the basis functions, with no approximations, in a straightforward manner. (3) For smooth solutions, PSC will converge to the actual solution exponentially as the number of basis functions is increased. A FD solution, on the other hand, never converges faster than algebraically with the number of grid points. Thus, for a given accuracy, PSC requires far less CPU time and memory than FD methods.
PSC has been applied successfully to solve problems in many fields, including fluid dynamics, meteorology, seismology, and relativistic astrophysics (cf. ). For example, PSC has been applied successfully to model stellar core collapse and construct equilibrium sequences of irrotational binary neutron stars . For black hole spacetimes PSC has been applied successfully to solve initial data for the standard field equations and the conformal field equations , to find apparent horizons , to solve the shift vector equation for a Kerr black hole , and to evolve Einstein’s equations in null quasi-spherical coordinates .
Here we evolve a spherically symmetric black hole spacetime in one spatial dimension by applying PSC methods to a hyperbolic formulation of Einstein’s equations, the “Einstein-Christoffel” (EC) system . A hyperbolic formulation provides a well-defined prescription for imposing boundary conditions. At a boundary, the fundamental fields can be decomposed into characteristic fields, which in the case of the EC system propagate either along the light cone or normal to the spatial foliation. Boundary conditions are imposed on the incoming (with respect to the computational domain) characteristic fields, but not on the outgoing fields. After imposing the boundary condition, the fundamental fields are reconstituted from the characteristic decomposition. If the excision boundary is placed inside the event horizon of a black hole, then (for appropriate coordinate systems) all characteristic fields are outgoing at this boundary, so no boundary conditions are needed there. Thus, black hole excision is a trivial operation. At the outer boundary of the domain, boundary conditions corresponding to no incoming radiation can be imposed on the incoming characteristic fields.
We find that our PSC method is able to evolve a spherically symmetric black hole spacetime forever. Furthermore, we find that the solution converges exponentially to the exact solution as the number of basis functions is increased. We discuss the time-stepping algorithms, outer boundary conditions, and gauge conditions required for stable evolution and how these depend on the particular slicing of the Schwarzschild geometry we wish to reproduce. We also show that our PSC method can handle dynamics by evolving a black hole spacetime containing a scalar field. Finally, we outline a strategy for applying this method to two black holes in three spatial dimensions using multiple domains, and we present tests of this domain decomposition idea in spherical symmetry.
In Sec. II we list the basic equations we use to evolve a spherically symmetric spacetime, boundary conditions, gauge conditions, initial data, and diagnostics. In Sec. III we introduce PSC and describe our numerical methods. In Sec. IV we present numerical evolutions of single black hole spacetimes in spherical symmetry. Finally, in Sec. V we discuss our results and the generalization of our methods to multiple dimensions and to binary black hole spacetimes.
## II Basic equations
### A Einstein-Christoffel system
We adopt the “Einstein-Christoffel” (EC) hyperbolic representation of Einstein’s equations . In this formulation the evolution equations are written in first-order
symmetric hyperbolic form, and all characteristic curves are directed either along the light cone or normal to the spatial foliation. The EC system takes as fundamental quantities the familiar three-metric and extrinsic curvature plus only 18 additional “connection” variables (in three spatial dimensions). Unlike some hyperbolic representations of general relativity, the EC system requires no derivatives of the stress-energy tensor.
We write the metric in the usual 3+1 form
$$ds^2=N^2dt^2+g_{ij}(dx^i+\beta ^idt)(dx^j+\beta ^jdt),$$
(1)
where $`g_{ij}`$ is the three-metric, $`\beta ^i`$ is the shift vector, and $`N`$ is the lapse function. In the EC formulation it is not the lapse function $`N`$ that is freely specifiable; instead, one arbitrarily prescribes the densitized lapse function $`\alpha `$, defined by
$$\alpha \frac{N}{\sqrt{g}},$$
(2)
where $`g`$ is the determinant of the three-metric. The use of a densitized lapse does not fix the temporal gauge freedom in any way: one can in principle obtain any lapse function $`N`$ by an appropriate choice of $`\alpha `$.
To write down the EC evolution equations, first define the new variables
$$f_{kij}\mathrm{\Gamma }_{(ij)k}+g_{ki}g^{lm}\mathrm{\Gamma }_{[lj]m}+g_{kj}g^{lm}\mathrm{\Gamma }_{[li]m},$$
(3)
where $`\mathrm{\Gamma }_{ij}^k`$ is the affine connection associated with $`g_{ij}`$, and parentheses and brackets denote symmetrization and antisymmetrization, respectively. The quantities $`f_{kij}`$ will be taken as fundamental variables along with $`g_{ij}`$ and the extrinsic curvature $`K_{ij}`$.
The EC evolution equations can be written in the form
$`\widehat{_0}g_{ij}=`$ $`2NK_{ij},`$ (5)
$`\widehat{_0}K_{ij}+Ng^{kl}_lf_{kij}=`$ $`N\{g^{kl}(K_{kl}K_{ij}2K_{ki}K_{lj})+g^{kl}g^{mn}(4f_{kmi}f_{[ln]j}+4f_{km[n}f_{l]ij}f_{ikm}f_{jln}`$ (10)
$`+8f_{(ij)k}f_{[ln]m}+4f_{km(i}f_{j)ln}8f_{kli}f_{mnj}+20f_{kl(i}f_{j)mn}13f_{ikl}f_{jmn})`$
$`_i_j\mathrm{ln}\alpha (_i\mathrm{ln}\alpha )(_j\mathrm{ln}\alpha )+2g_{ij}g^{kl}g^{mn}(f_{kmn}_l\mathrm{ln}\alpha f_{kml}_n\mathrm{ln}\alpha )`$
$`+g^{kl}[(2f_{(ij)k}f_{kij})_l\mathrm{ln}\alpha +4f_{kl(i}_{j)}\mathrm{ln}\alpha 3(f_{ikl}_j\mathrm{ln}\alpha +f_{jkl}_i\mathrm{ln}\alpha )]`$
$`8\pi S_{ij}+4\pi g_{ij}T\},`$
$`\widehat{_0}f_{kij}+N_kK_{ij}=`$ $`N\{g^{mn}[4K_{k(i}f_{j)mn}4f_{mn(i}K_{j)k}+K_{ij}(2f_{mnk}3f_{kmn})]`$ (14)
$`+2g^{mn}g^{pq}[K_{mp}(g_{k(i}f_{j)qn}2f_{qn(i}g_{j)k})+g_{k(i}K_{j)m}(8f_{npq}6f_{pqn})`$
$`+K_{mn}(4f_{pq(i}g_{j)k}5g_{k(i}f_{j)pq})]K_{ij}_k\mathrm{ln}\alpha `$
$`+2g^{mn}(K_{m(i}g_{j)k}_n\mathrm{ln}\alpha K_{mn}g_{k(i}_{j)}\mathrm{ln}\alpha )+16\pi g_{k(i}J_{j)}\}.`$
Here the symbol $`\widehat{_0}`$ is the time derivative operator normal to the spatial foliation, defined by
$$\widehat{_0}_t\mathrm{\pounds }_\beta ,$$
(15)
where $`\mathrm{\pounds }`$ denotes a Lie derivative. The matter terms are
$`\rho `$ $`n^\mu n^\nu T_{\mu \nu },`$ (17)
$`T`$ $`{}_{}{}^{(4)}g_{}^{\mu \nu }T_{\mu \nu },`$ (18)
$`J_i`$ $`n^\mu \gamma _i^\nu T_{\mu \nu },`$ (19)
$`S_{ij}`$ $`\gamma _i^\nu \gamma _j^\mu T_{\mu \nu },`$ (20)
where $`T_{\mu \nu }`$ is the stress-energy tensor, $`{}_{}{}^{(4)}g_{\mu \nu }^{}`$ is the four-metric, $`n^\mu `$ is the unit normal to the spatial foliation, and $`\gamma _i^\nu `$ is the spatial projection operator $`n^\nu n_i+{}_{}{}^{(4)}g_{i}^{\nu }`$.
Note that for each $`\{i,j\}`$ pair, the evolution equations (II A) can be written in the form (dropping the $`i`$ and $`j`$ indices on $`g_{ij}`$, $`K_{ij}`$, and $`f_{kij}`$)
$`\widehat{_0}g`$ $`2NK,`$ (22)
$`\widehat{_0}K+Ng^{kl}_lf_k`$ $`R,`$ (24)
$`\widehat{_0}f_k+N_kK`$ $`S,`$ (25)
where $`R`$ and $`S`$ are nonlinear terms that contain no derivatives of the fundamental variables (but may contain spatial derivatives of the arbitrary gauge function $`\alpha `$). Except for the right-hand sides, equations (22) are just $`\mathrm{}g=0`$ written in first order form. Thus one can think of the EC system (II A) as a set of six (one for each $`\{i,j\}`$ pair) coupled quasilinear scalar wave equations with nonlinear source terms.
A solution of the evolution equations (II A) is not a solution to Einstein’s equations unless twenty-two constraints are also satisfied. These are the Hamiltonian constraint
$`𝒞`$ $`g^{ij}g^{kl}\{2(_kf_{ijl}_if_{jkl})+K_{ik}K_{jl}K_{ij}K_{kl}+g^{mn}[f_{ikm}(5f_{jln}6f_{ljn})+13f_{ikl}f_{jmn}`$ (28)
$`+f_{ijk}(8f_{mln}20f_{lmn})]\}+16\pi \rho =0,`$
the three momentum constraints
$`𝒞_i`$ $`g^{kl}\{g^{mn}[K_{ik}(3f_{lmn}2f_{mnl})K_{km}f_{iln}]+_iK_{kl}_kK_{il}\}+8\pi J_i=0,`$ (29)
and the eighteen constraints
$$𝒞_{kij}_kg_{ij}2f_{kij}+4g^{lm}(f_{lm(i}g_{j)k}g_{k(i}f_{j)lm})=0$$
(30)
that relate $`f_{kij}`$ to spatial derivatives of the three-metric. If the constraints are satisfied for the initial data, they are preserved by the evolution equations for all time. As for any formulation of Einstein’s equations, however, numerical approximations may spoil this constraint-preserving property.
### B Reduction to spherical symmetry
The most general spherically symmetric metric can be written in the form
$`ds^2=`$ $`N^2dt^2+g_{rr}(dr+\beta ^rdt)^2`$ (32)
$`+g_Tr^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2),`$
where the transverse metric component is defined by
$$g_T\frac{g_{\theta \theta }}{r^2}=\frac{g_{\varphi \varphi }}{r^2\mathrm{sin}^2\theta }.$$
(33)
The two nonvanishing independent components of the extrinsic curvature are $`K_{rr}`$ and the transverse extrinsic curvature
$$K_T\frac{K_{\theta \theta }}{r^2}=\frac{K_{\varphi \varphi }}{r^2\mathrm{sin}^2\theta }.$$
(34)
The nonvanishing components of $`f_{kij}`$ are $`f_{rrr}`$, the transverse component
$$f_{rT}\frac{f_{r\theta \theta }}{r^2}=\frac{f_{\theta \theta r}}{2r^2}=\frac{f_{r\varphi \varphi }}{r^2\mathrm{sin}^2\theta }=\frac{f_{\varphi \varphi r}}{2r^2\mathrm{sin}^2\theta },$$
(35)
and the additional components
$`f_{rr\theta }=`$ $`g_{rr}\mathrm{cot}\theta ,`$ (37)
$`f_{\theta \theta \theta }=`$ $`2r^2g_T\mathrm{cot}\theta ,`$ (38)
$`f_{\theta \varphi \varphi }=`$ $`r^2g_T\mathrm{sin}\theta \mathrm{cos}\theta ,`$ (39)
$`f_{\varphi \theta \varphi }=`$ $`r^2g_T\mathrm{sin}\theta \mathrm{cos}\theta .`$ (40)
The evolution equations for these additional components (35) are automatically obeyed if the evolution equations for the metric are satisfied. We therefore do not treat these quantities as independent variables, and wherever they appear in the equations we replace them with the appropriate metric components using equations (35). Of course, all angular dependence due to these terms drops out.
We therefore take as fundamental variables the six quantities $`g_{rr}`$, $`g_T`$, $`K_{rr}`$, $`K_T`$, $`f_{rrr}`$, and $`f_{rT}`$. Using (II A), we obtain the following evolution equations for these variables:
$`_tg_{rr}\beta ^r_rg_{rr}=`$ $`2NK_{rr}+2g_{rr}_r\beta ^r,`$ (42)
$`_tg_T\beta ^r_rg_T=`$ $`2NK_T+2{\displaystyle \frac{\beta ^r}{r}}g_T,`$ (43)
$`_tK_{rr}\beta ^r_rK_{rr}+{\displaystyle \frac{N}{g_{rr}}}_rf_{rrr}=`$ $`N[2f_{rr}^r(f_{rr}^r+{\displaystyle \frac{1}{r}}{\displaystyle \frac{4f_{rT}}{g_T}}){\displaystyle \frac{6}{r^2}}+K_{rr}(2{\displaystyle \frac{K_T}{g_T}}K_r^r)6\left({\displaystyle \frac{f_{rT}}{g_T}}\right)^2_r^2\mathrm{ln}\stackrel{~}{\alpha }`$ (45)
$`(_r\mathrm{ln}\stackrel{~}{\alpha })^2+({\displaystyle \frac{4}{r}}f_{rr}^r)_r\mathrm{ln}\stackrel{~}{\alpha }]+2K_{rr}_r\beta ^r+4\pi N(Tg_{rr}2S_{rr}),`$
$`_tK_T\beta ^r_rK_T+{\displaystyle \frac{N}{g_{rr}}}_rf_{rT}=`$ $`N\left(K_TK_r^r+{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{2f_{rT}^2}{g_{rr}g_T}}{\displaystyle \frac{f_{rT}}{g_{rr}}}_r\mathrm{ln}\stackrel{~}{\alpha }\right)+{\displaystyle \frac{2\beta ^r}{r}}K_T,`$ (46)
$`_tf_{rrr}\beta ^r_rf_{rrr}+N_rK_{rr}=`$ $`N\left[4g_{rr}{\displaystyle \frac{K_T}{g_T}}\left(3{\displaystyle \frac{f_{rT}}{g_T}}f_{rr}^r+{\displaystyle \frac{2}{r}}_r\mathrm{ln}\stackrel{~}{\alpha }\right)K_{rr}\left(10{\displaystyle \frac{f_{rT}}{g_T}}+f_{rr}^r{\displaystyle \frac{2}{r}}+_r\mathrm{ln}\stackrel{~}{\alpha }\right)\right]`$ (49)
$`+3f_{rrr}_r\beta ^r+g_{rr}_r^2\beta ^r+16\pi NJ_rg_{rr},`$
$`_tf_{rT}\beta ^r_rf_{rT}+N_rK_T=`$ $`N\left[K_T\left(2{\displaystyle \frac{f_{rT}}{g_T}}f_{rr}^r_r\mathrm{ln}\stackrel{~}{\alpha }\right)\right]+\left(_r\beta ^r+{\displaystyle \frac{2\beta ^r}{r}}\right)f_{rT}.`$ (51)
Here
$$\stackrel{~}{\alpha }\alpha r^2\mathrm{sin}\theta =\frac{N}{g_T\sqrt{g_{rr}}},$$
(52)
and we have explicitly included the terms involving the Lie derivative of the shift vector.
The six fundamental variables obey four constraints that can be obtained from (II A):
$`𝒞`$ $`{\displaystyle \frac{_rf_{rT}}{g_{rr}g_T}}{\displaystyle \frac{1}{2r^2g_T}}+{\displaystyle \frac{f_{rT}}{g_{rr}g_T}}\left({\displaystyle \frac{2}{r}}+{\displaystyle \frac{7f_{rT}}{2g_T}}f_{rr}^r\right){\displaystyle \frac{K_T}{g_T}}\left(K_r^r+{\displaystyle \frac{K_T}{2g_T}}\right)+4\pi \rho =0,`$ (54)
$`𝒞_r`$ $`{\displaystyle \frac{_rK_T}{g_T}}+{\displaystyle \frac{2K_T}{rg_T}}{\displaystyle \frac{f_{rT}}{g_T}}\left(K_r^r+{\displaystyle \frac{K_T}{g_T}}\right)+4\pi J_r=0,`$ (55)
$`𝒞_{rrr}`$ $`_rg_{rr}+{\displaystyle \frac{8g_{rr}f_{rT}}{g_T}}2f_{rrr}=0,`$ (56)
$`𝒞_{rT}`$ $`_rg_T+{\displaystyle \frac{2g_T}{r}}2f_{rT}=0.`$ (57)
We do not explicitly solve the constraints during our evolution, but instead we use them as error estimators.
### C Boundary conditions
Boundary conditions are imposed on the above evolution equations (II B) via characteristic decomposition. Consider a first-order symmetrizable hyperbolic system
$$_tU+A^i_iU=R,$$
(58)
where $`U`$ is the vector of variables, $`R`$ is a vector, and the three $`A^i`$ are matrices. Then for a particular unit vector $`\xi _i`$, the solutions $`U_c`$ to the eigenvalue problem
$$A^i\xi _iU_c=v_cU_c$$
(59)
define the characteristic fields normal to the direction $`\xi _i`$, and the eigenvalues $`v_c`$ define the characteristic speeds of these fields. Each of the characteristic fields $`U_c`$ can be thought of as a plane wave solution moving in the direction $`\xi _i`$ with speed $`v_c`$. One is allowed to impose boundary conditions only on characteristic fields that propagate into the computational domain, but not on fields that propagate out of the domain.
For the evolution equations (II B), the characteristic fields in the radial direction ($`\xi _r=\sqrt{g_{rr}}`$) are
$`U_r^0g_{rr}`$ $`(v_c=\beta ^r),`$ (61)
$`U_t^0g_T`$ $`(v_c=\beta ^r),`$ (62)
$`U_r^\pm K_{rr}\pm {\displaystyle \frac{f_{rrr}}{\sqrt{g_{rr}}}}`$ $`(v_c=\beta ^r\pm \stackrel{~}{\alpha }g_T),`$ (63)
$`U_T^\pm K_T\pm {\displaystyle \frac{f_{rT}}{\sqrt{g_{rr}}}}`$ $`(v_c=\beta ^r\pm \stackrel{~}{\alpha }g_T).`$ (64)
The characteristic speeds of the metric variables correspond to propagation along the timelike normal to the foliation and the characteristic speeds of the other quantities correspond to propagation along the light cone. Thus, if the inner boundary of our domain moves along a spacelike trajectory, all characteristic speeds are negative (with respect to $`r`$) there, so no boundary conditions need to be imposed. At the outer boundary, boundary conditions are imposed only on those quantities with negative characteristic speeds (usually $`U_r^0`$, $`U_t^0`$, $`U_r^{}`$, and $`U_T^{}`$).
We have experimented with three types of outer boundary conditions. The first, which we call the freezing boundary condition, is
$$_tU_c=0$$
(65)
applied to all incoming characteristic fields $`U_c`$. This corresponds to no incoming radiation at the boundary; however, for nonlinear or inhomogeneous problems this is not strictly correct unless the boundary is at infinity.
The second is the Robin condition
$$_r[r^n(UU_{\mathrm{}})]=0,$$
(66)
which assumes that $`U`$ behaves like
$$U_{\mathrm{}}+\frac{\text{const}}{r^n}$$
(67)
at large $`r`$. For a given incoming variable $`U_c`$, appropriate values of the parameters $`U_{\mathrm{}}`$ and $`n`$ can be found from the analytic representations of the Schwarzschild geometry in section II D below.
Finally, the constraints (II B) can be used to derive mixed Neumann-Dirichlet boundary conditions for four of the characteristic fields: Equations (56) and (57) can be used directly as boundary conditions on $`U_r^0`$ and $`U_t^0`$, and Equations (54) and (55) can be combined to yield boundary conditions on $`U_T^{}`$ and $`U_T^+`$. We use only three of these boundary conditions—the ones for $`U_r^0`$, $`U_t^0`$, and $`U_T^{}`$—because $`U_T^+`$ is outgoing at the outer boundary and therefore needs no boundary condition there. Similarly, in three spatial dimensions the constraints can be used to derive twenty-two relations among the thirty characteristic fields, some of which can be used as boundary conditions on the incoming fields.
We describe our numerical implementation of boundary conditions in section III D.
### D Coordinate systems
It is convenient to choose a coordinate system in which the Schwarzschild geometry is time-independent. Furthermore, since we wish to include the apparent horizon in our computational domain, we must choose coordinates such that the spacelike slices labeled by constant values of coordinate $`t`$ penetrate the horizon and are nonsingular there. Here we list several coordinate systems that satisfy these properties.
#### 1 Kerr-Schild coordinates
In this coordinate system, also referred to as ingoing Eddington-Finkelstein coordinates , ingoing null rays have unit coordinate speed. In addition, the radial coordinate $`r`$ is chosen such that $`4\pi r^2`$ is the surface area of a sphere at that radius. In this coordinate system the Schwarzschild solution takes the form
$`g_{rr}=`$ $`1+{\displaystyle \frac{2M}{r}},`$ (69)
$`g_T=`$ $`1,`$ (70)
$`\stackrel{~}{\alpha }=`$ $`\left(1+{\displaystyle \frac{2M}{r}}\right)^1,`$ (71)
$`\beta ^r=`$ $`{\displaystyle \frac{2M}{r}}\left(1+{\displaystyle \frac{2M}{r}}\right)^1,`$ (72)
$`K_{rr}=`$ $`{\displaystyle \frac{2M}{r^2}}\left(1+{\displaystyle \frac{M}{r}}\right)\left(1+{\displaystyle \frac{2M}{r}}\right)^{1/2},`$ (73)
$`K_T=`$ $`{\displaystyle \frac{2M}{r^2}}\left(1+{\displaystyle \frac{2M}{r}}\right)^{1/2},`$ (74)
$`f_{rrr}=`$ $`{\displaystyle \frac{1}{r}}\left(4+{\displaystyle \frac{7M}{r}}\right),`$ (75)
$`f_{rT}=`$ $`{\displaystyle \frac{1}{r}},`$ (76)
where $`M`$ is the mass of the hole. The event horizon is coincident with the apparent horizon and is located at $`r=2M`$.
#### 2 Painlevé-Gullstrand coordinates
In this coordinate system the spatial three-metric is flat and the Schwarzschild solution is particularly simple:
$`g_{rr}=`$ $`1,`$ (78)
$`g_T=`$ $`1,`$ (79)
$`\stackrel{~}{\alpha }=`$ $`1,`$ (80)
$`\beta ^r=`$ $`\sqrt{{\displaystyle \frac{2M}{r}}},`$ (81)
$`K_{rr}=`$ $`\sqrt{{\displaystyle \frac{M}{2r^3}}},`$ (82)
$`K_T=`$ $`\sqrt{{\displaystyle \frac{2M}{r^3}}},`$ (83)
$`f_{rrr}=`$ $`{\displaystyle \frac{4}{r}},`$ (84)
$`f_{rT}=`$ $`{\displaystyle \frac{1}{r}}.`$ (85)
The horizon is again located at $`r=2M`$.
#### 3 Harmonic time slicing, areal radial coordinates
If one requires the time coordinate to satisfy $`\mathrm{}t=0`$, the radial coordinate $`r`$ to correspond to the areal radius, and the coordinate system to be regular at the horizon, then the Schwarzschild solution takes the form
$`g_{rr}=`$ $`\left(1+{\displaystyle \frac{2M}{r}}\right)\left(1+{\displaystyle \frac{4M^2}{r^2}}\right),`$ (87)
$`g_T=`$ $`1,`$ (88)
$`\stackrel{~}{\alpha }=`$ $`\left(1+{\displaystyle \frac{2M}{r}}\right)^1\left(1+{\displaystyle \frac{4M^2}{r^2}}\right)^1,`$ (89)
$`\beta ^r=`$ $`{\displaystyle \frac{4\stackrel{~}{\alpha }M^2}{r^2}},`$ (90)
$`K_{rr}=`$ $`{\displaystyle \frac{4M^2}{r^3}}\sqrt{\stackrel{~}{\alpha }}\left(2+{\displaystyle \frac{3M}{r}}+{\displaystyle \frac{4M^2}{r^2}}+{\displaystyle \frac{4M^3}{r^3}}\right),`$ (91)
$`K_T=`$ $`{\displaystyle \frac{4M^2}{r^3}}\sqrt{\stackrel{~}{\alpha }},`$ (92)
$`f_{rrr}=`$ $`{\displaystyle \frac{4}{r}}+{\displaystyle \frac{7M}{r^2}}+{\displaystyle \frac{12M^2}{r^3}}+{\displaystyle \frac{20M^3}{r^4}},`$ (93)
$`f_{rT}=`$ $`{\displaystyle \frac{1}{r}}.`$ (94)
The horizon is at $`r=2M`$.
#### 4 Fully harmonic coordinates
The Schwarzschild solution can also be written in a coordinate system where all coordinates satisfy $`\mathrm{}x^\mu =0`$ and are regular at the event horizon :
$`g_{rr}=`$ $`1+ϵ+ϵ^2+ϵ^3,`$ (96)
$`g_T=`$ $`\left(1+{\displaystyle \frac{M}{r}}\right)^2,`$ (97)
$`\stackrel{~}{\alpha }=`$ $`\left(1+{\displaystyle \frac{M}{r}}\right)^1\left(1+ϵ^2\right)^1\left(1+{\displaystyle \frac{3M}{r}}\right)^1,`$ (98)
$`\beta ^r=`$ $`ϵ^2\left(1+{\displaystyle \frac{M}{r}}\right)\left(1+ϵ^2\right)^1\left(1+{\displaystyle \frac{3M}{r}}\right)^1,`$ (99)
$`K_{rr}=`$ $`{\displaystyle \frac{K_T}{g_T}}\left(2+{\displaystyle \frac{3ϵ}{2}}+ϵ^2+{\displaystyle \frac{ϵ^3}{2}}\right),`$ (101)
$`K_T=`$ $`{\displaystyle \frac{4M^2}{r^3}}\sqrt{\stackrel{~}{\alpha }},`$ (102)
$`f_{rrr}=`$ $`{\displaystyle \frac{1}{g_T}}\left[{\displaystyle \frac{4}{r}}(1+ϵ^2)+{\displaystyle \frac{M}{r^2}}\left(112ϵ+9ϵ^2\right)\right],`$ (103)
$`f_{rT}=`$ $`{\displaystyle \frac{1}{r}}+{\displaystyle \frac{M}{r^2}},`$ (104)
where
$$ϵ\frac{2M}{r}\left(1+\frac{M}{r}\right)^1.$$
(105)
Here the horizon is located at $`r=M`$.
### E Einstein-Klein-Gordon system
To add dynamics to the spherically symmetric problem, we introduce a Klein-Gordon scalar field $`\varphi `$ with stress-energy
$$4\pi T_{ab}(_a\varphi )(_b\varphi )\frac{1}{2}g_{ab}(_c\varphi )(^c\varphi ).$$
(106)
Defining the quantities
$`\mathrm{\Pi }`$ $`N^1\widehat{_0}\varphi ,`$ (108)
$`\mathrm{\Phi }_i`$ $`_i\varphi ,`$ (109)
the matter terms (II A) are given by
$`4\pi \rho =`$ $`{\displaystyle \frac{1}{2}}\left(\mathrm{\Pi }^2+\mathrm{\Phi }^i\mathrm{\Phi }_i\right),`$ (111)
$`4\pi J_i=`$ $`\mathrm{\Pi }\mathrm{\Phi }_i,`$ (112)
$`4\pi T=`$ $`\mathrm{\Pi }^2\mathrm{\Phi }^i\mathrm{\Phi }_i,`$ (113)
$`4\pi S_{ij}=`$ $`\mathrm{\Phi }_i\mathrm{\Phi }_j+{\displaystyle \frac{1}{2}}g_{ij}(\mathrm{\Pi }^2\mathrm{\Phi }^i\mathrm{\Phi }_i).`$ (114)
The scalar field obeys $`\mathrm{}\varphi =0`$, which in spherical symmetry takes the form
$`_t\mathrm{\Pi }\beta ^r_r\mathrm{\Pi }+{\displaystyle \frac{N}{g_{rr}}}_r\mathrm{\Phi }_r=`$ $`N\left[\mathrm{\Pi }\left({\displaystyle \frac{K_{rr}}{g_{rr}}}+{\displaystyle \frac{2K_T}{g_T}}\right){\displaystyle \frac{\mathrm{\Phi }_r}{g_{rr}}}\left({\displaystyle \frac{4f_{rT}}{g_T}}{\displaystyle \frac{2}{r}}+_r\mathrm{ln}\stackrel{~}{\alpha }\right)\right],`$ (116)
$`_t\mathrm{\Phi }_r\beta ^r_r\mathrm{\Phi }_r+N_r\mathrm{\Pi }=`$ $`N\mathrm{\Pi }\left({\displaystyle \frac{f_{rrr}}{g_{rr}}}{\displaystyle \frac{2f_{rT}}{g_T}}{\displaystyle \frac{2}{r}}+_r\mathrm{ln}\stackrel{~}{\alpha }\right)+\mathrm{\Phi }_r_r\beta ^r.`$ (117)
Only the derivatives of the scalar field $`\varphi `$ appear in equations (106) and (II E), so $`\varphi `$ itself need not be evolved.
For the evolution equations (II E), the characteristic fields in the radial direction ($`\xi _r=\sqrt{g_{rr}}`$) are
$$U_\varphi ^\pm \mathrm{\Pi }\pm \frac{\mathrm{\Phi }_r}{\sqrt{g_{rr}}},$$
(118)
with characteristic speeds $`\beta ^r\pm \stackrel{~}{\alpha }g_T`$.
### F Apparent horizon
A marginal outer trapped 2-surface is defined by the equation
$$D_is^i+\left(s^is^jg^{ij}\right)K_{ij}=0,$$
(119)
where $`D_i`$ is the covariant derivative compatible with the three-metric, and $`s^i`$ is the outward-pointing spatial unit normal of the surface. In spherical symmetry, equation (119) reduces to
$$\frac{f_{rT}}{\sqrt{g_{rr}}}K_T=0.$$
(120)
The apparent horizon is the outermost surface at which (120) is satisfied. On each time slice, the coordinate radius of the horizon $`r_{\text{ah}}`$ is located by solving (120) using a standard root-finding algorithm.
If the horizon is to remain at a fixed coordinate radius as the spacetime evolves, the following relation must be obeyed at the horizon:
$$\frac{\beta ^r}{\stackrel{~}{\alpha }g_T}=\frac{1+8\pi r^2g_T(\rho J_r/\sqrt{g_{rr}})}{18\pi r^2g_T(\rho J_r/\sqrt{g_{rr}})}.$$
(121)
Equation (121) can be derived by setting the time derivative of (120) equal to zero and substituting the evolution and constraint equations to eliminate time and spatial derivatives.
### G Gauge conditions
The question of which gauge conditions one should impose on a numerically generated spacetime is one of the key unsolved problems in numerical relativity. In principle, the coordinate invariance of general relativity allows one to make this choice arbitrarily. However, a poor choice may not only obscure the physics one is searching for in the simulation, but may also allow rapidly-growing gauge modes that halt the code altogether.
#### 1 Algebraic conditions
The simplest gauge choices we consider here are algebraic ones: we set $`\stackrel{~}{\alpha }`$ and $`\beta ^i`$ equal to their analytic values for some parameterization of the Schwarzschild solution in section II D that we wish to reproduce numerically. While these gauge conditions are obviously applicable only to test problems, they provide a simplified setting in which to study the properties of our evolution scheme. Such conditions have been used extensively in 3D test problems.
Next we consider other algebraic gauge conditions that are independent of a particular analytic solution, but are still of limited generality. For instance, one might require that the radial coordinate remains areal, or in other words, that $`g_T`$ is time-independent. Using our variables, this condition can be written
$$\beta ^rf_{rT}\stackrel{~}{\alpha }g_T\sqrt{g_{rr}}K_T=0.$$
(122)
One might also require that the ingoing coordinate speed of light takes on a prescribed value $`c_{}`$:
$$\stackrel{~}{\alpha }g_T+\beta ^r+c_{}=0.$$
(123)
These conditions are not generalizable to two black holes because the first relies on the notion of an areal radial coordinate and the second assumes that a unique “ingoing” direction exists at every point in spacetime. Nevertheless, these and similar gauge conditions have proven useful for studies of single-black-hole spacetimes . Furthermore, if imposed only at one point, they can be used as boundary conditions on more general elliptic gauge choices, described below.
#### 2 Elliptic conditions
We also explore gauge conditions that should be applicable to general spacetimes. For the shift vector, we consider two elliptic equations due to : minimal strain
$`_r^2\beta ^r`$ $`+\left({\displaystyle \frac{f_{rrr}}{g_{rr}}}{\displaystyle \frac{2f_{rT}}{g_T}}\right)_r\beta ^r+\left[{\displaystyle \frac{_rf_{rrr}}{g_{rr}}}{\displaystyle \frac{4_rf_{rT}}{g_T}}2\left({\displaystyle \frac{f_{rrr}}{g_{rr}}}\right)^2+{\displaystyle \frac{2f_{rT}}{g_T}}\left({\displaystyle \frac{5f_{rrr}}{g_{rr}}}{\displaystyle \frac{f_{rT}}{g_T}}{\displaystyle \frac{4}{r}}\right)\right]\beta ^r`$ (125)
$`{\displaystyle \frac{K_{rr}g_T}{\sqrt{g_{rr}}}}_r\stackrel{~}{\alpha }+\stackrel{~}{\alpha }\sqrt{g_{rr}}g_T\left[{\displaystyle \frac{_rK_{rr}}{g_{rr}}}+{\displaystyle \frac{2K_Tf_{rT}}{(g_T)^2}}+{\displaystyle \frac{K_{rr}}{g_{rr}}}\left({\displaystyle \frac{f_{rrr}}{g_{rr}}}+{\displaystyle \frac{2}{r}}{\displaystyle \frac{8f_{rT}}{g_T}}\right)\right]=0,`$
which minimizes changes in the three-metric in a global sense, and minimal distortion,
$`_r^2\beta ^r`$ $`+\left({\displaystyle \frac{f_{rrr}}{g_{rr}}}{\displaystyle \frac{2f_{rT}}{g_T}}\right)_r\beta ^r+\left[{\displaystyle \frac{_rf_{rrr}}{g_{rr}}}{\displaystyle \frac{5_rf_{rT}}{g_T}}2\left({\displaystyle \frac{f_{rrr}}{g_{rr}}}\right)^2+{\displaystyle \frac{f_{rT}}{g_T}}\left({\displaystyle \frac{11f_{rrr}}{g_{rr}}}{\displaystyle \frac{5f_{rT}}{g_T}}{\displaystyle \frac{10}{r}}\right)\right]\beta ^r`$ (127)
$`+\left({\displaystyle \frac{K_T}{g_T}}{\displaystyle \frac{K_{rr}}{g_{rr}}}\right)\sqrt{g_{rr}}g_T_r\stackrel{~}{\alpha }+\stackrel{~}{\alpha }\sqrt{g_{rr}}g_T\left[{\displaystyle \frac{_rK_{rr}}{g_{rr}}}+{\displaystyle \frac{K_T}{g_T}}\left({\displaystyle \frac{f_{rrr}}{g_{rr}}}{\displaystyle \frac{2}{r}}\right)+{\displaystyle \frac{K_{rr}}{g_{rr}}}\left({\displaystyle \frac{f_{rrr}}{g_{rr}}}+{\displaystyle \frac{2}{r}}{\displaystyle \frac{8f_{rT}}{g_T}}\right)\right]=0,`$
which similarly minimizes changes of the conformal three-metric.
For the densitized lapse function, we consider the stationary mean curvature condition $`_tK=0`$, which in terms of our variables can be written
$`_r^2\stackrel{~}{\alpha }`$ $`+({\displaystyle \frac{f_{rrr}}{g_{rr}}}+{\displaystyle \frac{2f_{rT}}{g_T}}{\displaystyle \frac{4}{r}})_r\stackrel{~}{\alpha }+\stackrel{~}{\alpha }[{\displaystyle \frac{_rf_{rrr}}{g_{rr}}}+{\displaystyle \frac{2_rf_{rT}}{g_T}}2({\displaystyle \frac{f_{rrr}}{g_{rr}}}+{\displaystyle \frac{f_{rT}}{g_T}})({\displaystyle \frac{f_{rrr}}{g_{rr}}}{\displaystyle \frac{5f_{rT}}{g_T}})+{\displaystyle \frac{6}{r^2}}`$ (130)
$`{\displaystyle \frac{2g_{rr}}{r^2g_T}}g_{rr}({\displaystyle \frac{K_{rr}}{g_{rr}}}+{\displaystyle \frac{2K_T}{g_T}})^2{\displaystyle \frac{2f_{rrr}}{rg_{rr}}}+8\pi S_{rr}4\pi g_{rr}T]`$
$`{\displaystyle \frac{\beta ^r\sqrt{g_{rr}}}{g_T}}\left[{\displaystyle \frac{_rK_{rr}}{g_{rr}}}+{\displaystyle \frac{2_rK_T}{g_T}}+{\displaystyle \frac{4K_T}{g_T}}\left({\displaystyle \frac{1}{r}}{\displaystyle \frac{f_{rT}}{g_T}}\right)+{\displaystyle \frac{2K_{rr}}{g_{rr}}}\left({\displaystyle \frac{4f_{rT}}{g_T}}{\displaystyle \frac{f_{rrr}}{g_{rr}}}\right)\right]=0.`$
This condition was also discussed by , and is best known in the special case $`K=0`$, when it reduces to the familiar maximal slicing condition. Use of the stationary mean curvature lapse combined with either the minimal strain or minimal distortion shift vectors has been recently encouraged by .
Each of the elliptic equations (125), (127), and (130) requires two boundary conditions. At the horizon, we either set $`\stackrel{~}{\alpha }`$ and $`\beta ^r`$ to prescribed values, or we impose (121) and (123). At the outer boundary, we again can set $`\stackrel{~}{\alpha }`$ and $`\beta ^r`$ to prescribed values, we can impose (122) and (123), or we can impose Robin conditions of the form (66) on $`\stackrel{~}{\alpha }`$ and $`\beta ^r`$.
### H Mass
It is useful for diagnostic purposes to compute the mass of the spacetime. In spherical symmetry, the total mass inside an invariant spherical surface labeled by coordinate $`r`$ is well-defined and given by the Misner-Sharp formula , which for our variables reads
$$M_{\text{MS}}(r)\frac{r\sqrt{g_T}}{2}\left[1+\frac{r^2}{g_T}\left(K_T^2\frac{f_{rT}^2}{g_{rr}}\right)\right].$$
(131)
## III Pseudospectral collocation methods
### A Introduction
Consider a system of $`L`$ evolution equations of the form
$$_tf^{(\mathrm{})}=^{(\mathrm{})}[\{f^{(\mathrm{})}\}]$$
(132)
for $`1\mathrm{}L`$, where $`\{f^{(\mathrm{})}(\stackrel{}{x},t)\}`$ is the solution, and $`^{(\mathrm{})}[\{f^{(\mathrm{})}\}]`$ are (possibly nonlinear) functions of $`\{f^{(\mathrm{})}\}`$ and their spatial derivatives. Approximate each function $`f^{(\mathrm{})}`$ of the solution as a finite sum of basis functions $`\varphi _k^{(\mathrm{})}(\stackrel{}{x})`$,
$$f_N^{(\mathrm{})}(\stackrel{}{x},t)=\underset{k=0}{\overset{N1}{}}\stackrel{~}{f}_k^{(\mathrm{})}(t)\varphi _k^{(\mathrm{})}(\stackrel{}{x}).$$
(133)
For smooth functions as $`N\mathrm{}`$ the approximation is exact. Corresponding to the approximate solution $`\{f_N^{(\mathrm{})}\}`$ is a residual
$$R_N^{(\mathrm{})}=_tf_N^{(\mathrm{})}^{(\mathrm{})}[\{f_N^{(\mathrm{})}\}],$$
(134)
for each evolution equation.
In PSC the spectral coefficients $`\stackrel{~}{f}_k^{(\mathrm{})}(t)`$ are determined by demanding that the residuals $`R_N^{(\mathrm{})}`$ vanish at a fixed set of $`N`$ collocation points $`\stackrel{}{x}_n`$. In other words, it is demanded that the system of differential equations (132) be satisfied exactly at the collocation points $`\{\stackrel{}{x}_n\}`$. The choice of the collocation points is intimately related to the choice of basis functions used in the approximate solution. In the following subsection, we discuss how they are chosen.
### B Expansion basis and collocation points
For the remainder of this section, we will restrict ourselves to problems with one spatial dimension (1D). The choice of an expansion basis depends upon the particular problem being solved. For example, the natural expansion basis for a 1D problem with periodic boundary conditions is a Fourier series. For more general boundary conditions, such as the ones we will impose in our black hole evolutions, Chebyshev polynomials are a robust choice for the basis functions. Chebyshev polynomials are defined on the interval
$$𝕀=[1,1]$$
(135)
by
$$T_k(x)=\mathrm{cos}(k\mathrm{cos}^1x).$$
(136)
A function $`f`$ on $`𝕀`$ is approximated as <sup>*</sup><sup>*</sup>*For Chebyshev bases the conventional notation is that $`k`$ runs from $`0`$ to $`N`$, not $`N1`$; thus, there are $`N+1`$ coefficients and collocation points.
$$f_N(x,t)=\underset{k=0}{\overset{N}{}}\stackrel{~}{f}_k(t)T_k(x).$$
(137)
Note that in order to use this expansion, we must specify a mapping from our physical domain $`[r_{\text{min}},r_{\text{max}}]`$ to $`𝕀`$. The simplest choice is a linear mapping, but other choices may work better.
For a Chebyshev expansion, a convenient choice of the collocation points is
$$x_n=\mathrm{cos}\frac{\pi n}{N}.$$
(138)
At these collocation points, the Chebyshev polynomials satisfy the discrete orthogonality relation
$$\delta _{jk}=\frac{2}{N\overline{c}_k}\underset{n=0}{\overset{N}{}}\frac{1}{\overline{c}_n}T_j(x_n)T_k(x_n),$$
(139)
where
$$\overline{c}_k=\{\begin{array}{cc}2,\hfill & k=0\text{ or }N\hfill \\ 1,\hfill & 1kN1.\hfill \end{array}$$
(140)
Using the orthogonality relation, the spectral coefficients are given by
$$\stackrel{~}{f}_k=\frac{2}{N\overline{c}_k}\underset{n=0}{\overset{N}{}}\frac{1}{\overline{c}_n}f_N(x_n)T_k(x_n).$$
(141)
Since
$$T_k(x_n)=\mathrm{cos}\frac{\pi kn}{N},$$
(142)
fast cosine transforms can be used to compute (137) at the collocation points and to evaluate (141).
In PSC, the focus is not on the set of spectral coefficients $`\{\stackrel{~}{f}_k(t)\}`$, but on the equivalent set $`\{f(x_n,t)\}`$, the approximate solution evaluated at the collocation points. In particular, the approximate solution to (132) would be given by evolving
$$_tf_N^{(\mathrm{})}(x_n,t)=^{(\mathrm{})}(x_n,t),$$
(143)
for $`1\mathrm{}L`$ and $`0nN`$. Given initial conditions $`f^{(\mathrm{})}(x,0)`$, and appropriate boundary conditions, Equation (143) can be evolved forward in time using the method of lines, described in section III D. Since the focus is on grid-point values, and not the spectral coefficients, it is possible to reuse large amounts of code developed for FD methods.
### C Computation of derivatives
The main differences between PSC and FD in evolving (143) are the choice of collocation (grid) points $`x_n`$, how spatial derivatives are computed, and how boundary conditions are imposed. In PSC, spatial derivatives are computed analytically from the series expansion
$$\frac{f_N(x,t)}{x}=\underset{k=0}{\overset{N}{}}\stackrel{~}{f}_k(t)\frac{dT_k(x)}{dx}.$$
(144)
This derivative can be written as another sum over Chebyshev polynomials
$$\frac{f_N(x,t)}{x}=\underset{k=0}{\overset{N}{}}\stackrel{~}{f}_k^{}(t)T_k(x),$$
(145)
by using the simple recursion relation
$$c_k\stackrel{~}{f}_{k}^{}{}_{}{}^{}(t)=\stackrel{~}{f}_{k+2}^{}{}_{}{}^{}(t)+2(k+1)\stackrel{~}{f}_{k+1}(t),$$
(146)
where
$$c_k=\{\begin{array}{cc}2,\hfill & k=0\hfill \\ 1,\hfill & k1.\hfill \end{array}$$
(147)
Evaluating a derivative requires two fast transforms; the first to compute the spectral coefficients needed in the recursion relation (146), the second to evaluate (145).
### D Time evolution and application of boundary conditions
We evolve our hyperbolic system using the method of lines. In this method, we cast our system into the form (143) and use a standard ODE solver to integrate the equation in time. For most of the results presented in this paper, we have used a fourth-order explicit Runge-Kutta method. One of the drawbacks of using PSC is that the Chebyshev collocation points (138) are clustered near the domain boundaries. This places a more severe Courant stability limit $`\mathrm{\Delta }tO(N^2)`$ on a wave equation than for FD, where $`\mathrm{\Delta }t\mathrm{\Delta }xO(N^1)`$. Because of the superior spatial convergence of PSC, however, this restriction is not as severe as it may seem at first glance. In fact, to retain the accuracy gained by the spatial resolution, it may be necessary to use a time step smaller than that demanded by stability. Moreover, if the stability restriction on the time limit becomes too severe for practical evolutions, one can implement implicit or semi-implicit time-stepping schemes.
One of the advantages of PSC over FD is in how the boundary conditions are applied. In FD, derivatives are approximated by differences of field variables at grid points. The pattern of grid points used must typically be modified at the boundaries of the numerical grid. Consequently, boundary conditions can be difficult to formulate in FD. In PSC, on the other hand, the approximate solution is given over the entire domain. As seen in the previous section, derivatives are computed analytically; therefore nothing special needs to be done to compute the derivative at a boundary. Furthermore, since there are collocation points on the boundary of the domain, the application of boundary conditions is straightforward in PSC. One simply demands that the approximate solution satisfy the exact boundary condition at the boundary collocation point.
The boundary conditions are applied during the time step by modifying $`^{(\mathrm{})}[\{f^{(\mathrm{})}\}]`$ (cf. equation 132) at the boundary points so that the boundary conditions are satisfied. In this paper we are interested in applying boundary conditions on a hyperbolic system of evolution equations. As described in section II C, the solution to a hyperbolic system can be written in terms of characteristic fields that propagate with corresponding characteristic speeds. Physically we know that boundary conditions need only be applied to the incoming characteristic fields.
Therefore, to impose a boundary condition at a domain boundary $`x=x_b`$, we first compute the time derivatives of the characteristic fields $`U_c(x_b)`$ at the boundary. We then apply a boundary condition to those fields that are propagating into the domain; the remaining fields are untouched. Finally, we reconstruct the time derivatives of the fundamental variables at $`x_b`$ and use these values in the time update. Failure to impose a boundary condition on an incoming field or imposition of a boundary condition on an outgoing field almost always leads to an unstable evolution.
A Dirichlet boundary condition $`u_c(x_b,t)=g(t)`$ is applied by ensuring that the time derivative of the incoming characteristic field at the boundary collocation point is set to $`dg/dt`$. A boundary condition such as Neumann or Robin that involves the spatial derivative of the characteristic field is enforced by replacing the spatial derivative at the boundary with the appropriate value in order to satisfy the boundary condition.
### E Multiple domains
In order to use a PSC method for problems of dimension $`d`$ greater than unity the computational domain must be sufficiently simple that it can be mapped to $`𝕀^d`$ or $`𝕀^{d2}\times S^2`$ (where $`S^2`$ are two-spheres). For three dimensions, this typically means a cube, a sphere, or a spherical shell. If the computational domain is more complicated, then it must be decomposed into sub-domains that can each be mapped to one of these domains. For example, in two dimensions an L-shaped region can be decomposed into two adjacent rectangles.
The binary black hole problem will need to be solved using multiple domains. Therefore we test our ability to handle multiple domains on our one-dimensional problems. The use of multiple domains also provides a natural way of making our code run in parallel. We use KeLP to handle communication between multiple domains and for parallelization of our code.
The extension of our method from one domain to multiple domains is straightforward. We evolve each domain independently with communication done only at the boundaries. At the domain boundaries we compute the time derivatives of the characteristic fields in each domain. We then replace the time derivatives of the incoming characteristic fields at the boundary with the time derivatives of the outgoing characteristic fields of the neighboring domain. If there is no neighboring domain at a particular boundary, the external boundary condition is applied as described in section III D.
### F Solving elliptic equations
In addition to evolving our hyperbolic system of evolution equations (II A), we may need to solve elliptic equations in order to construct initial data for the Einstein-Klein-Gordon system or to enforce elliptic gauge conditions. Consider a linear elliptic equation of the form
$$(u(x))=f(x),$$
(148)
where $`u(x)`$ is the solution we are seeking. This can be cast as a matrix problem where, unlike for FD, the matrix corresponding to the linear operator $``$ is full. In 1D we solve this matrix equation directly, but for higher-dimensional problems, it will be more efficient to use an iterative method.
A nonlinear elliptic equation such as the Hamiltonian constraint can be solved either by the methods described in , or by linearizing the nonlinear system and iterating the linearized equations until a solution is found. The latter method is employed in the work described here.
### G Filtering
The errors in a spectral method are dominated by two types of terms of roughly equal magnitude. Truncation error arises from the neglect of the high-frequency terms that are not retained in the truncated series. Aliasing error occurs because each neglected high-frequency mode is indistinguishable from some retained lower-frequency mode when sampled only at the collocation points; for example, the functions $`\mathrm{sin}(\pi x/5)`$ and $`\mathrm{sin}(9\pi x/5)`$ take the same values on a grid of $`N`$ points $`x\{0,1,2,\mathrm{},N1\}`$. Because of aliasing error, power in the high-frequency mode, instead of being completely neglected, ends up contributing to the lower-frequency mode.
When solving a nonlinear system of equations it becomes important to control the aliasing error. This can be done by filtering the high-frequency modes of the retained series. For quadratic nonlinearities it is sufficient to zero the top third of the spectral coefficients to eliminate aliasing . In our 1D evolutions, we have found it necessary to filter only gauge variables that are computed from an elliptic equation. In effect we are smoothing the solutions to the elliptic gauge equations to eliminate high-frequency noise. Our preliminary investigations suggest that more extensive filtering may be required to produce stable evolutions in 3D.
## IV Numerical results
### A Schwarzschild black hole
In this section we evolve a time-independent slicing of a Schwarzschild black hole. We begin our numerical evolutions with initial data corresponding to one of the slicings given in Sec. II D. If the evolution equations are integrated exactly, the solution will remain time-independent. We can test the convergence of our method by measuring the deviation of the solution from the initial data at a given coordinate time, or by measuring the constraint quantities (II B), which are zero for the exact solution. For all evolutions, the interior of the hole is excised, and no boundary condition is applied at the inner boundary because all characteristic fields are outgoing (off the domain) there. In Table I we list the input parameters and the results for selected evolutions.
#### 1 Analytic gauge conditions
The simplest gauge treatment is to fix the gauge variables $`\stackrel{~}{\alpha }`$ and $`\beta ^r`$ to their initial values during the entire evolution. For example, one can begin the evolution with Kerr-Schild initial data (II D 1) and set $`\stackrel{~}{\alpha }`$ and $`\beta ^r`$ according to the analytic expressions (7172) for all time. In Figs. 1 and 2 we plot the norm of the Hamiltonian constraint and the deviation of $`g_{rr}`$ from the analytic solution versus time for such an evolution (run I from Table I). Each plot shows results for several spatial resolutions $`N_r`$ run at a fixed time resolution $`\mathrm{\Delta }t=0.007M`$. The features near $`t=10M`$ correspond to a small error pulse that begins at the outer boundary at $`t=0`$, grows like $`r^2`$ as it propagates inwards, and eventually falls into the hole. After several crossing times, the evolution settles into a steady state that converges to the analytic solution as one increases the spatial resolution. We end the evolutions at $`t=11000M`$ even though they clearly would have proceeded further. The convergence rate is exponential until machine roundoff errors dominate, as illustrated in Fig. 3. Repeating the evolutions shown in Figs. 13 for Painlevé-Gullstrand initial data (run I from Table I) yields similar results.
In Fig. 4 we show the norm of the Hamiltonian constraint for run I of Table I. This evolution is identical to run I except the initial data, as well as the values of $`\stackrel{~}{\alpha }`$ and $`\beta ^r`$ for all time, correspond to a time-independent harmonic slice of the Schwarzschild geometry (II D 3). Rather than settling to a steady state, the numerical solution grows exponentially at late times, eventually crashing the code. This is caused by a combination of high-frequency numerical instabilities, rapidly-growing gauge modes, and rapidly-growing constraint violating modes, all of which can be suppressed by appropriate changes in the evolution algorithm, as described below. Evolutions of fully harmonic initial data (II D 4) behave similarly. It is not known why these instabilities are absent in evolutions of Kerr-Schild and Painlevé-Gullstrand initial data. However, the dependence of stability on the choice of initial data should not be too surprising if one thinks of the initial data as a background solution and the numerical evolution as a perturbation on this background: in general, modifying the background solution can change the stability of perturbations.
For the evolutions shown in Fig. 4, freezing boundary conditions (65) are imposed on the incoming characteristic fields $`U_r^0`$, $`U_t^0`$, $`U_r^{}`$, and $`U_T^{}`$. One can suppress the constraint-violating modes seen in Fig. 4 by replacing the freezing boundary conditions on $`U_r^0`$, $`U_t^0`$, and $`U_T^{}`$ with constraint boundary conditions as discussed in Section II C. The resulting evolutions are shown in Figs. 5 and 6. Except for the evolution with $`N_r=32`$ discussed below, the Hamiltonian constraint $`𝒞`$ settles to a steady state that converges exponentially to zero. The same is true for the other three constraints $`𝒞_{rT}`$, $`𝒞_{rrr}`$, and $`𝒞_r`$. However, the metric quantities and other fundamental variables grow approximately quadratically with time, eventually causing the simulations to terminate. Because the constraints remain satisfied, we attribute this quadratic growth to a gauge mode.
The $`N_r=32`$ case shown in Figs. 5 and 6 suffers from high-frequency noise that grows exponentially in time. We have experimented with various methods of damping this noise, including filtering the fundamental variables after each time step and adding numerical dissipation terms to the equations. However, we have obtained best results by changing our fourth order Runge-Kutta time-stepping algorithm to an implicit backwards Euler scheme, which is much more dissipative. Figures 7 and 8 show the results of this modification. The evolution now satisfies the constraints at late times for sufficiently fine resolution, but still suffers from a quadratically growing gauge mode that causes the coarser resolution runs to crash. This gauge mode can be suppressed by applying active gauge conditions, as shown in Section IV A 2 below. Evolutions of fully harmonic initial data (II D 4) produce results similar to those shown in Figs. 48.
We note that even with analytic gauge conditions and freezing outer boundary conditions, evolutions of harmonic and fully harmonic initial data such as those shown in Fig. 4 become stable when the outer boundary is moved sufficiently close to the black hole (see runs I and I). A similar dependence on the outer boundary location has also been reported by others. A possible explanation for this is discussed briefly in: For a nonzero shift vector, any unstable zero-speed modes present in the solution will propagate inward from the outer boundary with speed $`\beta ^r`$, growing as they propagate. If the domain is sufficiently small, these modes do not have time to grow appreciably before they are swallowed by the horizon. As discussed previously, we find that constraint boundary conditions suppress exponentially growing modes, and thus allow evolutions with a larger outer boundary radius (runs I, I, I and I).
#### 2 Elliptic gauge conditions
Although choosing a time-independent $`\beta ^r`$ and $`\stackrel{~}{\alpha }`$ is the simplest gauge condition to implement, for an evolving numerical solution such a choice does not actively enforce any particular coordinate condition. In fact, it is remarkable that many of the cases discussed in Section IV A 1 remain stable when the coordinates experience small deviations from the exact solution. For more than one black hole in three spatial dimensions, one will almost certainly need general gauge conditions designed to prevent large changes in the numerical solution of a stationary or quasi-stationary spacetime.
Figure 9 shows the norm of the Hamiltonian constraint for an evolution of Painlevé-Gullstrand initial data. The gauge variables $`\beta ^r`$ and $`\stackrel{~}{\alpha }`$ are computed by solving the minimal strain and stationary mean curvature equations (125) and (130) after each time step. These elliptic equations require boundary conditions. We impose (121) and (123) at the current location of the horizon, which we recompute after every time step. For (123) we choose $`c_{}=2`$, which is the value of $`c_{}`$ at the horizon for the analytic solution (II D 2). At the outer boundary, we set $`\stackrel{~}{\alpha }=1`$ and we impose a Robin condition (66) on $`\beta ^r`$ with $`\beta _{\mathrm{}}^r=0`$, $`n=1/2`$. As seen in the figure, the evolution remains stable and convergent. To achieve stability, we find it necessary to apply a simple $`2/3`$ cutoff filter to $`\stackrel{~}{\alpha }`$ and $`\beta ^r`$ each time they are computed, and to impose constraint boundary conditions on $`U_r^0`$ and $`U_t^0`$ (but not on $`U_T^{}`$). Figure 10 shows the error in $`g_{rr}`$ for the same evolution. For the highest resolution, one can see a linearly-growing gauge mode. Although modes that grow linearly will eventually terminate a simulation, they pose no difficulty for long-term evolutions because a much longer run time can be achieved by a modest increase in resolution.
Similar results for the case of harmonic initial data are shown in Figs. 11 and 12. The evolution is stable and convergent, and the rapidly-growing gauge mode that terminated the simulation in the case of time-independent $`\beta ^r`$ and $`\stackrel{~}{\alpha }`$ (Section IV A 1, Fig. 8) now grows only linearly with time. As in the case shown in Fig. 8, we use a backwards Euler scheme for time evolution. For a fourth-order Runge-Kutta time discretization, results are similar except the evolutions with $`N_r=25`$, $`30`$, and $`32`$ are unstable.
### B Black hole plus scalar wave
In this subsection, we add dynamics to our spherically symmetric spacetime by including a Klein-Gordon scalar field as a matter source, as described in Section II E. To set up initial data, we first choose arbitrary initial values of $`\mathrm{\Pi }`$ and $`\mathrm{\Phi }_r`$. We then solve the Hamiltonian and momentum constraints via the standard York-Lichnerowicz conformal decomposition, using one of the time-independent representations of the Schwarzschild geometry discussed in Section II D as a background solution. Once we obtain the conformal factor and the trace-free longitudinal part of the extrinsic curvature from the constraints, we reconstruct the fundamental variables.
For the evolutions described here, the scalar field $`\mathrm{\Pi }`$ is initially a Gaussian centered at $`r=20M`$ with a width of $`5M`$ and an amplitude of $`0.02/M`$, and $`\mathrm{\Phi }_r`$ is zero everywhere. The outer boundary is located at $`r=120M`$ and the inner boundary is at $`1.75M`$. Here $`M`$ is the mass of the background solution, which is different than the actual mass of the spacetime. We choose a Kerr-Schild background, analytic gauge conditions, and a fourth-order Runge-Kutta time stepping algorithm. At the outer boundary we apply freezing boundary conditions (65) to $`U_\varphi ^{}`$ and $`U_r^{}`$, and constraint boundary conditions to $`U_r^0`$, $`U_t^0`$ and $`U_T^{}`$. There is no boundary condition imposed at the inner boundary.
To demonstrate our ability to handle multiple domains, for this evolution we cover the entire domain with $`8`$ equal-sized abutting subdomains, each using $`45`$ spectral coefficients. At each domain boundary, the incoming characteristic quantities in each domain are set equal to the corresponding outgoing quantities of the neighboring domain.
In Fig. 13 we plot the mass contained within radius $`r`$ as a function of $`r`$ for selected times. Initially the mass of the black hole is $`0.97M`$ and the mass of the entire spacetime is $`1.52M`$. The scalar field energy concentrated near $`r=20M`$ accounts for $`0.55M`$. As the evolution proceeds, the initial Gaussian scalar field pulse divides into incoming and outgoing pieces. The outgoing piece propagates to infinity, while the incoming piece is partially reflected off the Schwarzschild potential and partially swallowed by the black hole. At $`t=90M`$ the mass of the black hole has reached its final value of $`1.15M`$, and the initial outgoing pulse and the reflected pulse have not yet reached the outer boundary of the domain. By $`t=180M`$ the remaining scalar radiation has left the domain.
Figure 14 shows the coordinate radius and the areal radius of the apparent horizon versus time. The area of the horizon increases between $`t=15M`$ and $`t=30M`$ as the scalar field pulse falls into the black hole, and the areal radius asymptotically approaches the value $`2.31M`$, which is twice the mass of the final black hole as expected. Unlike the areal radius, the coordinate radius decreases with time until $`t=10M`$, after which it increases and eventually asymptotes to $`r=2.66M`$. One can see from the figure that on the initial slice the coordinate $`r`$ is nearly areal, as it would be for the Kerr-Schild background solution without a scalar field. At late times the deviation from an areal radial coordinate is large.
Figure 15 shows the norm of the Hamiltonian constraint as a function of time for the same evolution, as well as for other evolutions with different spatial resolutions but the same $`\mathrm{\Delta }t`$. The plots exhibit exponential convergence, even with nontrivial dynamics and multiple domains. There is, however, a small failure of convergence in the highest resolutions around $`t=120M`$. This is because the boundary condition (65), which is applied to $`U_r^{}`$ and $`U_\varphi ^{}`$ at the outer boundary $`r_b`$, is not strictly correct while a wave is passing through the boundary, and the error this introduces scales like $`r_{b}^{}{}_{}{}^{2}`$. We have verified this by repeating the evolutions from Fig. 15 with the inner boundary a factor of two closer, at $`60M`$. This is shown in Fig. 16. The resolution per subdomain is the same as in Fig. 15, but we use $`4`$ equal-sized subdomains instead of $`8`$. The small nonconvergent feature in Fig. 16 is approximately a factor of four larger than in Fig. 15, and occurs at an earlier time, $`t=60M`$, because the wave pulse reaches the outer boundary a factor of two earlier.
## V Discussion
We have found that, at least for spherical symmetry, applying PSC methods to hyperbolic formulations of general relativity can achieve stable evolutions of black holes with horizon excision. Excision itself is trivial as long as one uses a formulation in which all characteristic speeds are causal. Using realistic elliptic gauge conditions, our evolutions are limited only by linearly growing gauge modes that converge exponentially to zero with increasing resolution. These modes create no difficulty for long-term simulations because a small increase in resolution enables one to run much farther in time. We note that even when errors grow exponentially in time (e.g., Fig. 4), the high accuracy provided by PSC allows us to evolve to times of hundreds or sometimes thousands of $`M`$.
A hyperbolic formulation of Einstein’s equations provides a straightforward way in which to formulate and implement boundary conditions using the complete set of characteristic eigenfields provided by hyperbolicity. In principle, our method can be applied to any hyperbolic formulation, but so far we have only used the EC system. We have not investigated whether PSC can be used with non-hyperbolic formulations of Einstein’s equations such as ADM, as this would require a different treatment of the boundary conditions. It is entirely possible, however, that one might find boundary conditions that result in stable evolutions for such a formulation.
There has also been some concern about using hyperbolic representations of general relativity with complicated gauge conditions. This is because hyperbolic formulations of Einstein’s equations formally require the gauge quantities (shift and densitized lapse in the case of EC) to be prescribed functions of space and time, and not evolved quantities that couple to the fundamental variables. However, as long as the gauge variables are held fixed during each entire time step, as discussed by and , we find no fundamental difficulty in applying elliptic gauge conditions during our simulations.
By using the constraints as boundary conditions on the hyperbolic evolution equations, we have found that one can improve evolutions of the EC system. We have found similar improvement for finite-difference evolutions as well. Even in the general 3D case, applying constraint boundary conditions on the metric variables is straightforward. However, casting the Hamiltonian and momentum constraints as boundary conditions may be more difficult, because unlike in spherical symmetry, these constraints involve contractions of derivatives of fundamental variables.
The numerical techniques discussed in this paper should be generalizable to three spatial dimensions. For PSC evolutions of two black holes with excised horizons it will be necessary to use multiple computational domains (see Fig. 17). In spherical symmetry, we have shown how multiple domains can be easily implemented in a natural way by using characteristic fields to provide inter-domain boundary conditions. This method is directly applicable to abutting domains in 3D, and the extension from abutting domains to overlapping domains is straightforward . Work on 3D black hole evolutions using PSC is in progress.
###### Acknowledgements.
This work was supported in part by NSF grants PHY-9800737 and PHY-9900672 and NASA Grant NAG5-7264 to Cornell University, and NSF grants PHY-9802571 and PHY-9988581 to Wake Forest University. Computations were performed on the National Computational Science Alliance SGI Origin2000, and on the Wake Forest University Department of Physics IBM SP2 with support from an IBM SUR grant.
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# Sylvester Waves in the Coxeter Groups
## 1 Introduction
More than hundred years ago J.J.Sylvester stated and proved a theorem about restricted partition number $`W(s,𝐝^m)`$ of positive integer $`s`$ with respect to the $`m`$-tuple of positive integers $`𝐝^m=\{d_1,d_2,\mathrm{},d_m\}`$:
Theorem. The number $`W(s,𝐝^m)`$ of ways in which $`s`$ can be composed of (not necessarily distinct) $`m`$ integers $`d_1,d_2,\mathrm{},d_m`$ is made up of a finite number of waves
$$W(s,𝐝^m)=\underset{q}{\overset{maxq}{}}W_q(s,𝐝^m),W_q(s,𝐝^m)=\underset{k}{\overset{maxk}{}}W_{p_k|q}(s,𝐝^m),$$
(1)
where $`q`$ run over all distinct factors in $`d_1,d_2,\mathrm{},d_m`$ and $`W_{p_k|q}(s,𝐝^m)`$ denotes the coefficient of $`t^1`$ in the series expansion in ascending powers of $`t`$ of
$$F(s,𝐝^m,k;t)=e^{sw_k}\underset{r=1}{\overset{m}{}}\frac{1}{1e^{d_ru_k}},w_k=2\pi i\frac{p_k}{q}+t,u_k=2\pi i\frac{p_k}{q}t,$$
(2)
and $`p_1,p_2,\mathrm{},p_{maxk}`$ are all numbers (unity included) less than $`q`$ and prime to it.
$`W(s,𝐝^m)`$ is also a number of sets of positive integer solutions $`(x_1,x_2,\mathrm{},x_m)`$ of equation $`_r^md_rx_r=s`$. It is known that $`W(s,𝐝^m)`$ is equal to the coefficient of $`t^s`$ in the expansion of generating function
$$M(𝐝^m,t)=\underset{r=1}{\overset{m}{}}\frac{1}{1t^{d_r}}=\underset{s=0}{\overset{\mathrm{}}{}}W(s,𝐝^m)t^s.$$
(3)
If the exponents $`d_1,d_2,\mathrm{},d_m`$ become the series of integers $`1,2,3,\mathrm{},m`$, the number of waves is $`m`$ and $`W(s,𝐝^m)`$ of $`s`$ is usually referred to as a restricted partition number $`𝒫_m(s)`$ of $`s`$ into parts none of which exceeds $`m`$.
Another definition of $`W(s,𝐝^m)`$ comes from the polynomial invariant of finite reflection groups. Let $`M(𝐝^m,t)`$ is a Molien function of such a group $`G`$, $`d_r`$ are the degrees of basic invariants, and $`m`$ is the number of basic invariants . Then $`W(s,𝐝^m)`$ gives a number of algebraic independent polynomial invariants of the $`s`$-degree for group $`G`$.
Throughout his papers J.J.Sylvester gave different names for $`W(s,𝐝^m)`$ : quotity, denumerant, quot-undulant and quot-additant. Sometime after he discarded some of them. Because of a wide usage of $`W(s,𝐝^m)`$ not only as a partition number we shall call $`W(s,𝐝^m)`$ a Sylvester wave.
The Sylvester theorem is a very powerful tool not only in the trivial situation when $`m`$ is finite but also it was used for the purposes of asymptotic evaluations $`𝒫_m(s)`$, as well as for the main term of the Hardy-Ramanujan formulas for unrestricted partition number $`𝒫(s)`$ .
Recent progress in the self-dual problem of effective isotropic conductivity in two-dimensional three-component regular checkerboards and its further extension on the $`m`$-component anisotropic cases have shown an existence of algebraic equations with permutation invariance with respect to the action of the finite group $`G`$ permuting $`m`$ components. $`G`$ is a subgroup of symmetric group $`𝒮_m`$ and the coefficients in the equations are build out of algebraic independent polynomial invariants for group $`G`$. Here $`W(s,𝐝^m)`$ measures a degree of non-universality of the algebraic solution with respect to the different kinds of $`m`$-color plane groups.
Several proofs of Sylvester theorem are known ,. All of them make use of the Cauchy<sup>,</sup>s theory of residues. The recursion relations imposed on $`W(s,𝐝^m)`$ provide a combinatorial version of Sylvester formula. The classical example for the elementary (complex-variable-free) derivation was shown by Erd$`\ddot{o}`$s for the main term of the Hardy-Ramanujan formula. Recently an elementary derivation of Szekeres<sup>,</sup> formula for $`W(s,𝐝^m)`$ based on the recursion satisfied by $`W(s,𝐝^m)`$ was elaborated in . In this paper we give a new derivation of the Sylvester waves based on the recursion relation for $`W(s,𝐝^m)`$. We find also its zeroes and prove a lemma on parity properties of the Sylvester waves. Finally we present a list of the first ten Sylvester waves $`W(s,𝒮_m),m=1,\mathrm{},10`$ for symmetric groups $`𝒮_m`$ and for all Coxeter groups. In the Appendix we prove a conjecture on asymptotic behaviour of the least common multiple lcm($`1,2,\mathrm{},N`$) of the series of natural numbers.
## 2 Recursion relation for $`W(s,𝐝^m)`$.
We start with a recursion that follows from (3)
$$M(𝐝^m,t)M(𝐝^{m1},t)=t^{d_m}M(𝐝^m,t),$$
(4)
and after inserting the series expansions into the last equation we arrive at
$$W(s,𝐝^m)=W(s,𝐝^{m1})+W(sd_m,𝐝^m),d_ms,$$
(5)
where $`s`$ is assumed to be real. We apply now the recursive procedure (5) several times
$$W(s,𝐝^m)=\underset{p=0}{\overset{r_m}{}}W(spd_m,𝐝^{m1})+W(s(r_m+1)d_m,𝐝^m).$$
(6)
Let us consider the generic form of $`W(k\tau \{𝐝^m\}+s,𝐝^m),s<\tau \{𝐝^m\}`$ where $`k,s`$ and $`\tau \{𝐝^m\}`$ are the independent positive integers. We will choose them in such a way that
$$k\tau \{𝐝^m\}+s(r_m+1)d_m=(k1)\tau \{𝐝^m\}+s,\tau \{𝐝^m\}=(r_m+1)d_m.$$
(7)
Thus the relation (6) reads
$`W(k\tau \{𝐝^m\}+s,𝐝^m)`$ $`=`$ $`W((k1)\tau \{𝐝^m\}+s,𝐝^m)+`$ (8)
$`{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}W(k\tau \{𝐝^m\}pd_m+s,𝐝^{m1}),\delta _m={\displaystyle \frac{\tau \{𝐝^m\}}{d_m}}.`$
As follows from (7), in order to return via the recursive procedure from $`W(k\tau \{𝐝^m\}+s,𝐝^m)`$ to $`W((k1)\tau \{𝐝^m\}+s,𝐝^m)`$ we must use $`\tau \{𝐝^m\}`$ which have $`d_m`$ as a divisor. Due to the arbitrariness of $`d_m`$ it is easy to conclude that all exponents $`d_1,d_2,\mathrm{},d_m`$ serve as the divisors of $`\tau \{𝐝^m\}`$. In other words $`\tau \{𝐝^m\}`$ is the least common multiple lcm of the exponents $`d_1,d_2,\mathrm{},d_m`$
$$\tau \{𝐝^m\}=\mathrm{𝗅𝖼𝗆}(d_1,d_2,\mathrm{},d_m).$$
(9)
Actually $`\tau \{𝐝^m\}`$ does play a role of the ”period ” of $`W(s,𝐝^m)`$. But strictly speaking it is not a periodic function with respect to the integer variable $`s`$ as could be seen from (8). The rest of the paper clarifies this hidden periodicity.
As we have mentioned above, $`W(s,𝐝^m)`$ gives a number of algebraic independent polynomial invariants of the $`s`$-degree for the group $`G`$. The situation becomes more transparent if we deal with the irreducible Coxeter group where the degrees $`d_r`$ and the number of basic invariants $`m`$ are well known.
Table 1. The ”periods $`\tau (G)`$ of $`W(s,𝐝^m)`$ for the irreducible Coxeter groups. $`G`$ $`A_m`$ $`B_m`$ $`D_m`$ $`G_2`$ $`F_4`$ $`E_6`$ $`\tau (G)`$ $`(m+1)`$ 2$`(m)`$ 2$`(m)`$ $`\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}6}`$ 24 360 $`G`$ $`E_7`$ $`E_8`$ $`H_3`$ $`H_4`$ $`I_2`$(2m) $`I_2`$(2m+1) $`\tau (G)`$ 2520 2520 30 60 2 m 2 (2m+1)
where $`(m)`$=lcm(1,2,3,…,$`m`$) is the least common multiple of the series of the natural numbers.
$`(m)`$ can be viewed as $`\tau (𝒮_m)`$ for symmetric group $`𝒮_m`$ or, in other words, as a ”period ” of the restricted partition number $`𝒫_m(s)`$. This makes it possible to pose a question about asymptotic behaviour of $`\tau (𝒮_m)`$ with $`m\mathrm{}`$. $`(m)`$ is a very fast growing function: $`(1)`$=1, $`(10)`$=2520, $`(20)`$=232792560, $`(30)`$=2329089562800 etc. Actually $`\frac{\mathrm{ln}(m)}{m}`$ oscillates infinitely many times around 1 and the function $`(m)`$ has an exponential increase with the asymptotic law <sup>1</sup><sup>1</sup>1It seems to be strange but we have not found throughout the textbooks on number theory any discussion about the asymptotics of lcm(1,2,3,…,$`m`$). The formula (10) was conjectured by one of the authors (LGF) based on the numerical calculations and proved by Z.Rudnick (Tel-Aviv Univ., Israel) which had communicated this proof to us. The proof is given in Appendix A.
$$\underset{m\mathrm{}}{lim}\frac{\mathrm{ln}(m)}{m}=1.$$
(10)
## 3 Polynomial representation for $`W(s,𝐝^m)`$.
Making use of the relations (8,9) we obtain the exact formula for $`W(k\tau \{𝐝^m\}+s,𝐝^m)`$ for different $`𝐝^m`$. We will treat it in an ascending order in the number $`m`$ of exponents. The first steps are simple and they yield $`\underset{¯}{𝐝^1=(d_1)}`$ $`,\tau \{𝐝^1\}>s0`$
$$W(kd_1+s,𝐝^1)=W(s,𝐝^1)=\mathrm{\Psi }_{d_1}(s)=\{\begin{array}{cc}1,\hfill & s=0(modd_1)\hfill \\ 0,\hfill & s0(modd_1)\hfill \end{array}$$
(11)
$`\mathrm{\Psi }_{d_1}(s)`$ may be represented as a sum of prime roots of unit of degree $`d_1`$:
$$\mathrm{\Psi }_{d_1}(s)=\frac{1}{d_1}\underset{k=0}{\overset{d_11}{}}\mathrm{exp}(\frac{2\pi iks}{d_1})=\frac{1}{d_1}\{\begin{array}{cc}1+\mathrm{cos}\pi s+2_{k=1}^{d_1/21}\mathrm{cos}\frac{2\pi ks}{d_1},\hfill & \text{even }d_1\hfill \\ 1+2_{k=1}^{(d_11)/2}\mathrm{cos}\frac{2\pi ks}{d_1},\hfill & \text{odd }d_1\hfill \end{array}.$$
$`\underset{¯}{𝐝^2=(d_1,d_2)}`$ $`,\tau \{𝐝^2\}>s0`$
$$W(k\tau \{𝐝^2\}+s,𝐝^2)=W(s,𝐝^2)+k\underset{p=0}{\overset{\delta _21}{}}W(|spd_2|,𝐝^1).$$
(12)
$`\underset{¯}{𝐝^3=(d_1,d_2,d_3)}`$ $`,\tau \{𝐝^3\}>s0`$
$`W(k\tau \{𝐝^3\}+s,𝐝^3)`$ $`=`$ $`W(s,𝐝^3)+k{\displaystyle \underset{p=0}{\overset{\delta _31}{}}}W(|spd_3|,𝐝^2)+`$ (13)
$`{\displaystyle \frac{k(k+1)}{2}}{\displaystyle \frac{\tau \{𝐝^3\}}{\tau \{𝐝^2\}}}{\displaystyle \underset{p=0}{\overset{\delta _31}{}}}{\displaystyle \underset{q=0}{\overset{\delta _21}{}}}W(|spd_3qd_2|,𝐝^1).`$
Now it is simple to deduce by induction that in the general case $`W(k\tau \{𝐝^m\}+s,𝐝^m)`$ has a polynomial representation with respect to $`k`$
$$W(k\tau \{𝐝^m\}+s,𝐝^m)=A_{m1}^m(s)k^{m1}+A_{m2}^m(s)k^{m2}+\mathrm{}+A_1^m(s)k+A_0^m(s,𝐝^m),$$
(14)
where $`A_{mr}^m(s)`$ is based on the $`\tau \{𝐝^r\}`$-periodic functions as well as the entire $`W(s,𝐝^m)`$ is based on the $`\tau \{𝐝^m\}`$-periodic functions. The coefficient in the leading term can be written in a closed form
$`A_{m1}^m(s)={\displaystyle \frac{1}{(m1)!}}{\displaystyle \frac{\tau ^{m2}\{𝐝^m\}}{\tau \{𝐝^2\}\tau \{𝐝^3\}\mathrm{}\tau \{𝐝^{m1}\}}}\times `$
$`\times {\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}{\displaystyle \underset{q=0}{\overset{\delta _{m1}1}{}}}\mathrm{}{\displaystyle \underset{v=0}{\overset{\delta _21}{}}}W(|spd_mqd_{m1}\mathrm{}vd_2|,𝐝^1).`$ (15)
With $`d_1=1`$ we have $`W(|spd_mqd_{m1}\mathrm{}vd_2|,1)=1`$, which makes $`A_{m1}^m(s)`$ independent of $`s`$ and gives an asymptotics of $`W(s,𝐝^m)`$ for $`sm`$
$$A_{m1}^m(s)=\frac{\tau ^{m1}\{𝐝^m\}}{(m1)!m!},W(s,𝐝^m)\stackrel{s\mathrm{}}{}\frac{s^{m1}}{(m1)!m!}.$$
(16)
Now we are ready to prove the statement about splitting of $`W(s,𝐝^m)`$ into periodic and non-periodic parts.
Lemma 3.1. The Sylvester wave $`W(s,𝐝^m)`$ can be represented in the following way
$$W(s,𝐝^m)=Q_m^m(s)+\underset{j=1}{\overset{m1}{}}Q_j^m(s)s^{mj},$$
(17)
where $`Q_j^m(s)`$ is a periodic function with the period $`\tau \{𝐝^j\}=\mathrm{𝗅𝖼𝗆}(d_1,d_2,\mathrm{},d_j)`$.
Proof. We start with the identity for the polynomial representation for $`W(k\tau \{𝐝^m\}+s,𝐝^m)`$
$$W((k+1)\tau \{𝐝^m\}+s,𝐝^m)=W(k\tau \{𝐝^m\}+s+\tau \{𝐝^m\},𝐝^m),$$
that can be transformed, using (14), into
$`A_{m1}^m(s)(k+1)^{m1}+A_{m2}^m(s)(k+1)^{m2}+\mathrm{}+A_1^m(s)(k+1)+W(s,𝐝^m)=`$
$`A_{m1}^m(s+\tau \{𝐝^m\})k^{m1}+A_{m2}^m(s+\tau \{𝐝^m\})k^{m2}+\mathrm{}+A_1^m(s+\tau \{𝐝^m\})k+`$
$`W(s+\tau \{𝐝^m\},𝐝^m).`$ (18)
The last identity generates a finite number of coupled difference equations for the coefficients $`A_r^m(s)`$
$$A_{mr}^m(s+\tau \{𝐝^m\})=\underset{j=1}{\overset{r}{}}C_{mj}^{mr}A_{mj}^m(s),\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}1}rm,$$
(19)
where $`C_n^k`$ denotes a binomial coefficient. The first equation $`(r=1)`$
$$A_{m1}^m(s+\tau \{𝐝^m\})=A_{m1}^m(s)$$
declares that $`A_{m1}^m(s)`$ is an arbitrary $`\tau \{𝐝^m\}`$-periodic function. We can specify the last statement taking into account (14) that actually $`A_{m1}^m(s)`$ is $`\tau \{𝐝^1\}`$-periodic function which will be denoted as $`Q_1^m(s)`$. The second equation $`(r=2)`$
$$A_{m2}^m(s+\tau \{𝐝^m\})=A_{m2}^m(s)+(m1)A_{m1}^m(s)$$
can be solved completely
$$A_{m2}^m(s)=Q_2^m(s)+(m1)sQ_1^m(s),$$
(20)
where $`Q_2^m(s+\tau \{𝐝^2\})=Q_2^m(s)`$. Continuing this procedure, it is not difficult to prove by induction that for any $`r`$ we have
$$A_{mr}^m(s)=\underset{j=1}{\overset{r}{}}C_{mj}^{mr}Q_j^m(s)s^{rj},$$
(21)
where $`Q_j^m(s+\tau \{𝐝^j\})=Q_j^m(s)`$. Since $`W(s,𝐝^m)=A_0^m(s)`$ we arrive finally at (17) by inserting $`r=m`$ into equation (21), that splits $`W(s,𝐝^m)`$, in accordance with the Sylvester theorem, into periodic and non-periodic parts.$`\mathrm{}`$
## 4 Partition identities and zeroes of $`W(s,𝐝^m)`$.
In this section we assume that the variable $`s`$ has only integer values.
Consider a new quantity
$$V(s,𝐝^m)=W(s\xi \{𝐝^m\},𝐝^m),\xi \{𝐝^m\}=\frac{1}{2}\underset{i=1}{\overset{m}{}}d_i.$$
(22)
Lemma 4.1. $`V(s,𝐝^m)`$ has the following parity properties:
$$V(s,𝐝^{2m})=V(s,𝐝^{2m}),V(s,𝐝^{2m+1})=V(s,𝐝^{2m+1}).$$
(23)
Proof. A basic recursion relation (5) can be rewritten for $`V(s,𝐝^m)`$
$$V(s,𝐝^m)V(sd_m,𝐝^m)=V(s\frac{d_m}{2},𝐝^{m1}).$$
(24)
The last relation produces two equations in a new variable $`q=s\frac{d_m}{2}`$
$`V(q,𝐝^{m1})`$ $`=`$ $`V(q+{\displaystyle \frac{d_m}{2}},𝐝^m)V(q{\displaystyle \frac{d_m}{2}},𝐝^m),`$ (25)
$`V(q,𝐝^{m1})`$ $`=`$ $`V(q+{\displaystyle \frac{d_m}{2}},𝐝^m)V(q{\displaystyle \frac{d_m}{2}},𝐝^m).`$ (26)
Hence if $`V(q,𝐝^m)`$ is an even function of $`q`$, then $`V(q,𝐝^{m1})`$ is an odd one, and vice versa. Because $`V(q,𝐝^1)`$ is an even function, we arrive at (23). $`\mathrm{}`$
Corollary. If $`s_1+s_2+2\xi \{𝐝^m\}=0`$, then
$$W(s_1,𝐝^m)=(1)^{m+1}W(s_2,𝐝^m)$$
Proof. This follows from the parity properties and after substitution two new variables $`s_1=s\xi \{𝐝^m\},s_2=s\xi \{𝐝^m\}`$ into (23). $`\mathrm{}`$
Lemma 4.2. Let $`m`$-tuple $`\{𝐝^m\}`$ generates the Sylvester wave $`W(s,𝐝^m)`$ . Then for every integer $`p`$ a $`m`$-tuple $`\{p𝐝^m\}=\{pd_1,pd_2,\mathrm{},pd_m\}`$ generates the following Sylvester wave
$$W(s,p𝐝^m)=\mathrm{\Psi }_p(s)W(\frac{s}{p},𝐝^m),\text{or}V(s,p𝐝^m)=\mathrm{\Psi }_p(sp\xi \{𝐝^m\})V(\frac{s}{p},𝐝^m),$$
(27)
where the periodic function $`\mathrm{\Psi }_p(s)=\mathrm{\Psi }_p(s+p)`$ is defined in (11).
Proof. According to the definition (3)
$$\underset{s}{}W(s,p𝐝^m)t^s=\underset{s}{}W(s,𝐝^m)t^{ps}=\underset{s^{}}{}W(\frac{s^{}}{p},𝐝^m)t^s^{}$$
Equating powers of $`t`$ in the latter equation and taking into account that $`s^{}/p`$ must be integer we obtain (27). $`\mathrm{}`$
Lemma 4.3. Let $`m`$-tuple $`\{𝐝^m\}`$ generates the Sylvester wave $`W(s,𝐝^m)`$ . Then $`W(s,𝐝^m)`$ has the following zeroes:
* If all exponents $`d_r`$ are mutually prime numbers, then the zeroes $`𝔰_0(𝐝^m)`$ read
$`𝔰_0(𝐝^m)`$ $`=`$ $`1,2,\mathrm{},{\displaystyle \underset{r=1}{\overset{m}{}}}d_r+1,\text{if}m=2k+1,`$
$`𝔰_0(𝐝^m)`$ $`=`$ $`1,2,\mathrm{},{\displaystyle \underset{r=1}{\overset{m}{}}}d_r+1,\xi \{𝐝^m\},\text{if}m=2k;`$ (28)
* If all exponents $`d_r`$ have a maximal common factor $`p`$, then $`W(s,𝐝^m)`$ has infinite number of zeroes $`𝔖_1(𝐝^m)`$ which are distributed in the following way
$$𝔖_1(𝐝^m)=𝔰_1(𝐝^m)\{/p\},$$
where $`\{/p\}`$ denotes a set of integers $``$ with deleted integers of modulo $`p`$
$$\{/p\}=\{\mathrm{},p1,p+1,\mathrm{},1,1,\mathrm{},p1,p+1,\mathrm{}\}$$
and
$`𝔰_1(𝐝^m)`$ $`=`$ $`p,2p,\mathrm{},{\displaystyle \underset{r=1}{\overset{m}{}}}d_r+p,\text{if}m=2k+1,`$
$`𝔰_1(𝐝^m)`$ $`=`$ $`p,2p,\mathrm{},{\displaystyle \underset{r=1}{\overset{m}{}}}d_r+p,\xi \{𝐝^m\},\text{if}m=2k.`$ (29)
Proof. Consider again the relation (6) which we rewrite as follows
$$\underset{s=0}{\overset{\mathrm{}}{}}W(s,𝐝^m)t^s=\frac{1}{1t^{d_m}}\underset{s^{}=0}{\overset{\mathrm{}}{}}W(s^{},𝐝^{m1})t^s^{}$$
(30)
assuming that the exponents in $`𝐝^m`$ are sorted in the ascending order. Note that the influence of the new $`d_m`$ exponent appears only in terms $`t^s`$ with $`sd_m`$. This enables us to deduce that the values of $`W(s,𝐝^{m1})`$ and $`W(s,𝐝^m)`$ coincide at integer positive values $`s=0,1,\mathrm{},d_m1`$. This means that for $`0sd_m1`$ we have $`W(s,𝐝^m)=W(s,𝐝^{m1})`$. Recalling the main recursion relation (5) we conclude that
$$W(s,𝐝^m)=0(d_ms1).$$
Using the last relation for $`m`$ and $`m1`$ in (5) we can find also
$$W(sd_m,𝐝^m)=0(d_{m1}s1)W(s,𝐝^m)=0(d_{m1}d_ms1).$$
Repeating this procedure and taking into account that at the last step it leads to the zeroes of $`\mathrm{\Psi }_{d_1}`$ which are located at $`(1d_1s1)`$, we get the set of the zeroes for $`W(s,𝐝^m)`$ with odd number of exponents $`m=2k+1`$
$$W(s,𝐝^m)=0(1\underset{i=1}{\overset{m}{}}d_is1).$$
(31)
The eveness of $`m`$ gives one more zero of $`W(s,𝐝^m)`$ which arises from the parity properties of $`V(s,𝐝^m)`$, namely, $`V(0,𝐝^{2k})=0`$. The last equality immediately generates a zero $`\xi \{𝐝^{2k}\}`$ of $`W(s,𝐝^{2k})`$ that together with (31) proves the first part (28) of Lemma 3.
The second part of Lemma 3 follows from (27) and from the first part of (28) because a set of integers $`\{/p\}`$ represents the zeroes of the periodic function $`\mathrm{\Psi }_p(s)`$. $`\mathrm{}`$
The complexity of the exponents sequence $`\{𝐝^m\}`$ and its large length make the calculative procedure of restoration of $`Q_j^m(s)`$ very cumbersome. Therefore it is important to find the inner properties of $`\{𝐝^m\}`$ when this procedure could be essentially reduced.
Lemma 4.4. Let $`m`$-tuple $`\{𝐝^m\}=\{d_1,d_2,\mathrm{},d_r,d_r,\mathrm{},d_m\}`$ contains an exponent $`d_r`$ twice. Then the Sylvester wave $`V(s,𝐝^m)`$ is related to the Sylvester wave $`V(s,𝐝^{m_1})`$ produced by the the non-degenerated tuple $`\{𝐝^{m_1}\}=\{d_1,d_2,\mathrm{},d_r,..,.d_m,2d_r\}`$ as follows
$$V(s,𝐝^m)=V(s\frac{d_r}{2},𝐝^{m_1})+V(s+\frac{d_r}{2},𝐝^{m_1}).$$
(32)
Proof. According to the definition (3)
$$(1+t^{d_r})\underset{s}{}W(s,𝐝^{m_1})t^s=\underset{s}{}W(s,𝐝^m)t^s.$$
Taking into account that $`\xi \{𝐝^{m_1}\}\xi \{𝐝^m\}=d_r/2`$ and equating powers of $`t`$ in the latter equation we obtain the stated relation (32) according to the definition (22). $`\mathrm{}`$
We will make worth of relation (32) during the evaluation of the expression $`V(s,𝐝^m)`$ for the Coxeter group D<sub>m</sub>.
## 5 Recursion formulas for $`V(s,𝐝^m)`$.
The shift (22) transforms the relation (8) into
$$V(s+\tau \{𝐝^m\},𝐝^m)=V(s,𝐝^m)+\underset{p=0}{\overset{\delta _m1}{}}V(s+\tau \{𝐝^m\}\lambda _pd_m,𝐝^{m1}),\lambda _p=p+\frac{1}{2}$$
(33)
and the relation (17) into
$$V(s,𝐝^m)=R_m^m(s)+\underset{j=1}{\overset{m1}{}}R_j^m(s)s^{mj},$$
(34)
where
$$R_j^m(s)=\underset{i=1}{\overset{j}{}}C_{mi}^{ji}(\xi \{𝐝^m\})^{ji}Q_i^m(s\xi \{𝐝^m\}),$$
i.e., $`R_1^m(s)=Q_1^m(s\xi \{𝐝^m\});R_2^m(s)=Q_2^m(s\xi \{𝐝^m\})(m1)\xi \{𝐝^m\}Q_1^m(s\xi \{𝐝^m\})`$ etc. This means that the functions $`R_j^m(s)`$ and $`Q_j^m(s)`$ have the same period $`\tau \{𝐝^j\}`$.
Inserting the expansion (34) into the relation (33) and equating powers of $`s`$ we can obtain for $`k=1,2,\mathrm{},m1`$
$$\underset{j=1}{\overset{k}{}}C_{mj}^{m1k}R_j^m(s)\tau \{𝐝^m\}^{k+1j}=\underset{p=0}{\overset{\delta _m1}{}}\underset{j=1}{\overset{k}{}}R_j^{m1}(s\lambda _pd_m)C_{m1j}^{m1k}(\tau \{𝐝^m\}\lambda _pd_m)^{kj}.$$
(35)
For the first successive values of $`k`$ the latter equation (35) gives
$`R_1^m(s)`$ $`=`$ $`{\displaystyle \frac{1}{(m1)\tau \{𝐝^m\}}}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}R_1^{m1}(s\lambda _pd_m),`$
$`R_2^m(s)`$ $`=`$ $`{\displaystyle \frac{1}{(m2)\tau \{𝐝^m\}}}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}R_2^{m1}(s\lambda _pd_m)+{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}({\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _p}{\delta _m}})R_1^{m1}(s\lambda _pd_m),`$
$`R_3^m(s)`$ $`=`$ $`{\displaystyle \frac{1}{(m3)\tau \{𝐝^m\}}}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}R_3^{m1}(s\lambda _pd_m)+{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}({\displaystyle \frac{1}{2}}{\displaystyle \frac{\lambda _p}{\delta _m}})R_2^{m1}(s\lambda _pd_m)+`$ (36)
$`{\displaystyle \frac{m2}{2}}\tau \{𝐝^m\}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}({\displaystyle \frac{1}{6}}{\displaystyle \frac{\lambda _p}{\delta _m}}+{\displaystyle \frac{\lambda _p^2}{\delta _m^2}})R_1^{m1}(s\lambda _pd_m).`$
It is easy to see that in the summands of the latter formulas (36) there appear the Bernoulli polynomials $`_i(1\frac{\lambda _p}{\delta _m})`$ : $`_0(x)=1,_1(x)=x1/2,_2(x)=x^2x+1/6,_3(x)=x^33/2x^2+1/2x`$, etc . Continuing the evaluation of the general expression for $`R_j^m(s),\mathrm{\hspace{0.33em}1}<j<m`$, we arrive at
Lemma 5.1. $`R_j^m(s)`$ for $`1j<m`$ is given by the formula
$$R_j^m(s)=\frac{1}{mj}\underset{l=0}{\overset{j1}{}}(\tau \{𝐝^m\})^{l1}C_{m1j+l}^l\underset{p=0}{\overset{\delta _m1}{}}_l(1\frac{\lambda _p}{\delta _m})R_{jl}^{m1}(s\lambda _pd_m).$$
(37)
Proof. Before going to the proof we recall two identities for the Bernoulli polynomials ,
$$_l(x+y)_l(x)=\underset{j=1}{\overset{l}{}}C_l^jy^j_{lj}(x),_l(1+x)_l(x)=lx^{l1}.$$
(38)
Using the definition (34) we check that formula (37) satisfies (33).
$`V(s,𝐝^m)=R_m^m(s)+{\displaystyle \underset{j=1}{\overset{m1}{}}}s^j{\displaystyle \underset{l=j}{\overset{m1}{}}}C_l^j{\displaystyle \frac{(\tau \{𝐝^m\})^{lj1}}{l}}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}_{lj}(1{\displaystyle \frac{\lambda _p}{\delta _m}})R_{ml}^{m1}(s\lambda _pd_m)=`$
$`R_m^m(s)+{\displaystyle \underset{l=1}{\overset{m1}{}}}{\displaystyle \frac{(\tau \{𝐝^m\})^{l1}}{l}}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}R_{ml}^{m1}(s\lambda _pd_m){\displaystyle \underset{j=1}{\overset{l}{}}}C_l^j\left({\displaystyle \frac{s}{\tau \{𝐝^m\}}}\right)^j_{lj}(1{\displaystyle \frac{\lambda _p}{\delta _m}})=`$
$`R_m^m(s)+{\displaystyle \underset{l=1}{\overset{m1}{}}}{\displaystyle \frac{(\tau \{𝐝^m\})^{l1}}{l}}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}R_{ml}^{m1}(s\lambda _pd_m)\left[_l(1+{\displaystyle \frac{s\lambda _pd_m}{\tau \{𝐝^m\}}})_l(1{\displaystyle \frac{\lambda _p}{\delta _m}})\right],`$ (39)
where we use the first of the identities (38). Having in mind the $`\tau \{𝐝^m\}`$-periodicity of functions $`R_j^m(s)`$ and $`R_j^{m1}(s)`$ and the second identity (38) we may rewrite the difference in the l.h.s of relation (33) in the following form:
$`V(s,𝐝^m)V(s\tau \{𝐝^m\},𝐝^m)=`$ (40)
$`{\displaystyle \underset{l=1}{\overset{m1}{}}}{\displaystyle \frac{(\tau \{𝐝^m\})^{l1}}{l}}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}R_{ml}^{m1}(s\lambda _pd_m)\left[_l(1{\displaystyle \frac{\lambda _p}{\delta _m}}+{\displaystyle \frac{s}{\tau \{𝐝^m\}}})_l({\displaystyle \frac{\lambda _p}{\delta _m}}+{\displaystyle \frac{s}{\tau \{𝐝^m\}}})\right]`$
$`{\displaystyle \underset{l=1}{\overset{m1}{}}}{\displaystyle \frac{(\tau \{𝐝^m\})^{l1}}{l}}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}R_{ml}^{m1}(s\lambda _pd_m)l\left({\displaystyle \frac{s\lambda _pd_m}{\tau \{𝐝^m\}}}\right)^{l1}=`$
$`{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}{\displaystyle \underset{l=0}{\overset{m2}{}}}(s\lambda _pd_m)^lR_{m1l}^{m1}(s\lambda _pd_m)={\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}V(s\lambda _pd_m,𝐝^{m1}).\mathrm{}`$
The formula (37) enables to restore all terms $`R_k^m(s)`$ except the last $`R_m^m(s)`$. Actually we can learn about it from the following consideration. Let us separate $`R_{mk}^m(s)`$ in the following way
$$R_{mk}^m(s)=_{mk}^m(s)+r_{mk}^m(s),\mathrm{\hspace{0.33em}\hspace{0.33em}0}km1,$$
(41)
where
$`_{mk}^m(s)`$ $`=`$ $`{\displaystyle \underset{l=1}{\overset{mk1}{}}}{\displaystyle \frac{(\tau \{𝐝^m\})^{l1}}{l+k}}C_{l+k}^k{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}_l(1{\displaystyle \frac{\lambda _p}{\delta _m}})R_{mkl}^{m1}(s\lambda _pd_m)`$ (42)
$`r_{mk}^m(s)`$ $`=`$ $`{\displaystyle \frac{1}{k\tau \{𝐝^m\}}}{\displaystyle \underset{p=0}{\overset{\delta _m1}{}}}R_{mk}^{m1}(s\lambda _pd_m),r_{mk}^m(s)=r_{mk}^m(sd_m),(k0)`$ (43)
The representation (41) and $`d_m`$-periodicity of the function $`r_{mk}^m(s)`$ make possible to prove the following
Lemma 5.2. $`R_{mk}^m(s)`$ for $`0km1`$ and $`_{mk}^m(s)`$ for $`0<km1`$ satisfy the recursion relation
$`R_{mk}^m(s)R_{mk}^m(sd_m)=_{mk}^m(s)_{mk}^m(sd_m)=`$
$`{\displaystyle \underset{j=k+1}{\overset{m1}{}}}\left\{(d_m)^{jk}C_j^kR_{mj}^m(sd_m)+({\displaystyle \frac{d_m}{2}})^{j1k}C_{j1}^kR_{mj}^{m1}(s{\displaystyle \frac{d_m}{2}})\right\}.`$ (44)
Proof. Inserting (34) into (24), expanding the powers of binomials into sums and equating the powers of $`s`$ in the latter equation we obtain the relation (5) for the function $`R_{mk}^m(s)`$, $`0km1`$. Using the definition (41) we immediately arrive at the relation for the function $`_{mk}^m(s)`$ , $`0<km1`$. $`\mathrm{}`$
In the special case $`k=0`$ the general relation (5) produces the recursion for $`R_m^m(s)`$
$$R_m^m(s)R_m^m(sd_m)=\underset{j=1}{\overset{m1}{}}\left\{(d_m)^jR_{mj}^m(sd_m)+(\frac{d_m}{2})^{j1}R_{mj}^{m1}(s\frac{d_m}{2})\right\}.$$
(45)
We can not use directly (42) for $`k=0`$ since $`r_m^m(s)`$ can not be derived from (43). But it is a good mathematical intuition to exploit the formula (42) for $`k=0`$ in order to prove
Lemma 5.3. $`_m^m(s)`$ is given by the formula
$$_m^m(s)=\underset{l=1}{\overset{m1}{}}\frac{(\tau \{𝐝^m\})^{l1}}{l}\underset{p=0}{\overset{\delta _m1}{}}_l(1\frac{\lambda _p}{\delta _m})R_{ml}^{m1}(s\lambda _pd_m).$$
(46)
Proof. In order to prove that $`_m^m(s)`$ given by (46) satisfies the difference equation (45) we consider a difference $`_m^m(s)_m^m(sd_m)=\mathrm{\Delta }_m(s)=\mathrm{\Delta }_m^1(s)+\mathrm{\Delta }_m^2(s)`$:
$$\mathrm{\Delta }_m(s)=\underset{l=1}{\overset{m1}{}}\frac{(\tau \{𝐝^m\})^{l1}}{l}\underset{p=0}{\overset{\delta _m1}{}}_l(1\frac{\lambda _p}{\delta _m})\left[R_{ml}^{m1}(s\lambda _pd_m)R_{ml}^{m1}(s\lambda _{p+1}d_m)\right]$$
with
$$\mathrm{\Delta }_m^1(s)=\underset{l=1}{\overset{m1}{}}\frac{(\tau \{𝐝^m\})^{l1}}{l}\left\{_l(1\frac{1}{2\delta _m})_l(\frac{1}{2\delta _m})\right\}R_{ml}^{m1}(s\frac{d_m}{2}),$$
$$\mathrm{\Delta }_m^2(s)=\underset{l=1}{\overset{m1}{}}\frac{(\tau \{𝐝^m\})^{l1}}{l}\underset{p=1}{\overset{\delta _m}{}}\left\{_l(1\frac{\lambda _p}{\delta _m})_l(1\frac{\lambda _p}{\delta _m}+\frac{1}{\delta _m})\right\}R_{ml}^{m1}(s\lambda _pd_m).$$
The first term $`\mathrm{\Delta }_m^1(s)`$ is calculated with the help of one of the identities (38):
$$\mathrm{\Delta }_m^1(s)=\underset{l=1}{\overset{m1}{}}(\frac{d_m}{2})^{l1}R_{ml}^{m1}(s\frac{d_m}{2}).$$
(47)
Using another identity from (38) we may write for $`\mathrm{\Delta }_m^2(s)`$:
$$\mathrm{\Delta }_m^2(s)=\underset{l=1}{\overset{m1}{}}\underset{j=1}{\overset{l}{}}\frac{(\tau \{𝐝^m\})^{l1}}{l}C_l^j(\frac{1}{\delta _m})^j\underset{p=1}{\overset{\delta _m}{}}_{lj}(1\frac{\lambda _{p1}}{\delta _m})R_{ml}^{m1}(s\lambda _pd_m).$$
Changing here summation order $`_{k=l+1}^{m1}_{j=l+1}^k=_{j=l+1}^{m1}_{k=j}^{m1}`$ and comparing the inner sum with (37) we arrive at
$$\mathrm{\Delta }_m^2(s)=\underset{j=1}{\overset{m1}{}}(d_m)^jR_{mj}^m(sd_m)$$
(48)
Then (47) and (48) prove the Lemma. $`\mathrm{}`$
From this Lemma follows an existence of $`d_m`$-periodic function $`r_m^m(s)=r_m^m(sd_m)`$ which could not be derived from (43). Unknown function $`r_m^m(s)`$ corresponds to vanishing harmonics in the r.h.s. of equation (5). We are free to choose any basic system of continuous $`\tau \{𝐝^m\}`$-periodic functions. This arbitrariness can affect behaviour of $`W(s,𝐝^m)`$ only for non-integer $`s`$ that does not violate the recursion relation (5). In the rest of the paper we will choose a basic system of the simplest periodic functions sin and cos.
The function $`r_m^m(s)`$ corresponds to the harmonics of the type
$$\left\{\begin{array}{c}\mathrm{sin}\\ \mathrm{cos}\end{array}\right\}\frac{2\pi n}{d_m}s$$
Because the parity properties of $`R_m^m(s)`$ coincide with that of $`V(s,𝐝^m)`$ itself we can rewrite (34) in the following form
$`V(s,𝐝^{2m})`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2m1}{}}}R_j^{2m}(s)s^{2mj}+_{2m}^{2m}(s)+{\displaystyle \underset{n}{}}\rho _n^{2m}\mathrm{sin}{\displaystyle \frac{2\pi n}{d_{2m}}}s,`$ (49)
$`V(s,𝐝^{2m+1})`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{2m}{}}}R_j^{2m+1}(s)s^{2m+1j}+_{2m+1}^{2m+1}(s)+{\displaystyle \underset{n}{}}\rho _n^{2m+1}\mathrm{cos}{\displaystyle \frac{2\pi n}{d_{2m+1}}}s.`$ (50)
In order to produce $`r_m^m(s)`$ we use some of zeroes $`𝔰`$, described in the preceding Section, constructing a system of linear equations for $`[(m+1)/2]`$ coefficients $`\rho _n`$; $`n`$ runs from $`1`$ to $`m/2`$ in (49 and from $`0`$ to $`(m1)/2`$ in (50). We use a trivial identity $`V(\xi (𝐝^m),𝐝^m)=1`$, and choose the values of $`s`$ out of the set $`𝔰`$, adding homogeneous equations to arrive at a non-degenerate inhomogeneous system of linear equations. This system is solved further to produce the final expression for corresponding Sylvester wave. These explicit expressions are given in the next Section. Appendix B presents two instructive examples of the above procedure.
## 6 Sylvester waves $`V(s,G)`$.
We start with the symmetric group $`𝒮_m`$ because of two reasons: first, of their relation with restricted partition numbers and , second, they arranged natural basis to utilize the Sylvester waves $`V(s,G)`$ in all Coxeter groups.
### 6.1 Symmetric groups $`𝒮_m`$.
Making use of the procedure developed in the previous section we present here first ten Sylvester waves $`V(s,𝒮_m),m=1,\mathrm{},10`$. <sup>2</sup><sup>2</sup>2Having in mind the results of Sylvester , and Glaisher for restricted partition numbers for $`m9`$ we repeat them adding a formula for $`m=10`$. The list of $`V(s,𝒮_m)`$ can be simply continued up to any finite $`m`$ with the help of the symbolic code written in Mathematica language .
$`\underset{¯}{G=𝒮_m}`$ $`,d_r=1,2,3,\mathrm{},m,\xi (𝒮_m)=\frac{m(m+1)}{4},`$
$`V(s,𝒮_1)`$ $`=`$ $`1,`$
$`V(s,𝒮_2)`$ $`=`$ $`{\displaystyle \frac{s}{2}}{\displaystyle \frac{1}{4}}\mathrm{sin}\pi s,`$
$`V(s,𝒮_3)`$ $`=`$ $`{\displaystyle \frac{s^2}{12}}{\displaystyle \frac{7}{72}}{\displaystyle \frac{1}{8}}\mathrm{cos}\pi s+{\displaystyle \frac{2}{9}}\mathrm{cos}{\displaystyle \frac{2\pi s}{3}},`$
$`V(s,𝒮_4)`$ $`=`$ $`{\displaystyle \frac{s^3}{144}}{\displaystyle \frac{s}{96}}(5+3\mathrm{cos}\pi s)+{\displaystyle \frac{1}{8}}\mathrm{sin}{\displaystyle \frac{\pi s}{2}}{\displaystyle \frac{2}{9\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{2\pi s}{3}},`$ (51)
$`V(s,𝒮_5)`$ $`=`$ $`{\displaystyle \frac{s^4}{2880}}{\displaystyle \frac{11s^2}{1152}}{\displaystyle \frac{s}{64}}\mathrm{sin}\pi s+{\displaystyle \frac{17083}{691200}}{\displaystyle \frac{2}{27}}\mathrm{cos}{\displaystyle \frac{2\pi s}{3}}+`$
$`{\displaystyle \frac{1}{8\sqrt{2}}}\mathrm{cos}{\displaystyle \frac{\pi s}{2}}+{\displaystyle \frac{2}{25}}(\mathrm{cos}{\displaystyle \frac{2\pi s}{5}}+\mathrm{cos}{\displaystyle \frac{4\pi s}{5}}),`$
$`V(s,𝒮_6)`$ $`=`$ $`{\displaystyle \frac{s^5}{86400}}{\displaystyle \frac{91s^3}{103680}}+{\displaystyle \frac{s^2}{768}}\mathrm{sin}\pi s+{\displaystyle \frac{s}{829440}}(919110240\mathrm{cos}{\displaystyle \frac{2\pi s}{3}})`$
$`{\displaystyle \frac{161}{9216}}\mathrm{sin}\pi s{\displaystyle \frac{1}{16\sqrt{2}}}\mathrm{sin}{\displaystyle \frac{\pi s}{2}}{\displaystyle \frac{1}{81\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{2\pi s}{3}}{\displaystyle \frac{1}{18}}\mathrm{sin}{\displaystyle \frac{\pi s}{3}}`$
$`{\displaystyle \frac{2}{25\sqrt{5}}}(\mathrm{sin}{\displaystyle \frac{\pi }{5}}\mathrm{sin}{\displaystyle \frac{4\pi s}{5}}+\mathrm{sin}{\displaystyle \frac{2\pi }{5}}\mathrm{sin}{\displaystyle \frac{2\pi s}{5}}),`$
$`V(s,𝒮_7)`$ $`=`$ $`{\displaystyle \frac{s^6}{3628800}}{\displaystyle \frac{s^4}{20736}}+{\displaystyle \frac{s^2}{38400}}(71+25\mathrm{cos}\pi s){\displaystyle \frac{s}{81\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{2\pi s}{3}}`$
$`{\displaystyle \frac{52705}{6096384}}{\displaystyle \frac{77}{4608}}\mathrm{cos}\pi s{\displaystyle \frac{1}{32}}\mathrm{cos}{\displaystyle \frac{\pi s}{2}}{\displaystyle \frac{5}{486}}\mathrm{cos}{\displaystyle \frac{2\pi s}{3}}{\displaystyle \frac{1}{18}}\mathrm{cos}{\displaystyle \frac{\pi s}{3}}+`$
$`{\displaystyle \frac{2}{25\sqrt{5}}}(\mathrm{cos}{\displaystyle \frac{2\pi s}{5}}\mathrm{cos}{\displaystyle \frac{4\pi s}{5}})+{\displaystyle \frac{2}{49}}(\mathrm{cos}{\displaystyle \frac{2\pi s}{7}}+\mathrm{cos}{\displaystyle \frac{4\pi s}{7}}+\mathrm{cos}{\displaystyle \frac{6\pi s}{7}}),`$
$`V(s,𝒮_8)`$ $`=`$ $`{\displaystyle \frac{s^7}{203212800}}{\displaystyle \frac{17s^5}{9676800}}+{\displaystyle \frac{s^3}{8294400}}(1343+225\mathrm{cos}\pi s)+`$
$`s({\displaystyle \frac{16133}{4976640}}{\displaystyle \frac{1}{256}}\mathrm{cos}{\displaystyle \frac{\pi s}{2}}+{\displaystyle \frac{1}{243}}\mathrm{cos}{\displaystyle \frac{2\pi s}{3}}{\displaystyle \frac{31}{12288}}\mathrm{cos}\pi s)+`$
$`{\displaystyle \frac{1}{32}}(\mathrm{sin}{\displaystyle \frac{\pi s}{4}}\mathrm{sin}{\displaystyle \frac{3\pi s}{4}}){\displaystyle \frac{1}{128}}\mathrm{sin}{\displaystyle \frac{\pi s}{2}}+{\displaystyle \frac{1}{162\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{2\pi s}{3}}+{\displaystyle \frac{1}{18\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{\pi s}{3}}+`$
$`{\displaystyle \frac{4}{125}}(\mathrm{sin}{\displaystyle \frac{2\pi }{5}}\mathrm{sin}{\displaystyle \frac{4\pi s}{5}}\mathrm{sin}{\displaystyle \frac{\pi }{5}}\mathrm{sin}{\displaystyle \frac{2\pi s}{5}})`$
$`{\displaystyle \frac{1}{49}}(\mathrm{sin}{\displaystyle \frac{2\pi s}{7}}\mathrm{csc}{\displaystyle \frac{\pi }{7}}\mathrm{sin}{\displaystyle \frac{4\pi s}{7}}\mathrm{csc}{\displaystyle \frac{2\pi }{7}}+\mathrm{sin}{\displaystyle \frac{6\pi s}{7}}\mathrm{csc}{\displaystyle \frac{3\pi }{7}}),`$
$`V(s,𝒮_9)`$ $`=`$ $`{\displaystyle \frac{s^8}{14631321600}}{\displaystyle \frac{19s^6}{418037760}}+{\displaystyle \frac{145597s^4}{16721510400}}+{\displaystyle \frac{s^3}{73728}}\mathrm{sin}\pi s`$
$`s^2({\displaystyle \frac{67293991}{140460687360}}+{\displaystyle \frac{1}{4374}}\mathrm{cos}{\displaystyle \frac{2\pi s}{3}})`$
$`s({\displaystyle \frac{1}{256\sqrt{2}}}\mathrm{sin}{\displaystyle \frac{\pi s}{2}}+{\displaystyle \frac{1}{1458\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{2\pi s}{3}}+{\displaystyle \frac{205}{98304}}\mathrm{sin}\pi s)+{\displaystyle \frac{199596951167}{56184274944000}}+`$
$`{\displaystyle \frac{1}{64}}(\mathrm{cos}{\displaystyle \frac{\pi s}{4}}\mathrm{csc}{\displaystyle \frac{\pi }{8}}\mathrm{cos}{\displaystyle \frac{3\pi s}{4}}\mathrm{csc}{\displaystyle \frac{3\pi }{8}})+{\displaystyle \frac{2}{125}}(\mathrm{cos}{\displaystyle \frac{4\pi s}{5}}\mathrm{cos}{\displaystyle \frac{2\pi s}{5}})`$
$`{\displaystyle \frac{5}{512\sqrt{2}}}\mathrm{cos}{\displaystyle \frac{\pi s}{2}}+{\displaystyle \frac{257}{17496}}\mathrm{cos}{\displaystyle \frac{2\pi s}{3}}+{\displaystyle \frac{1}{36\sqrt{3}}}\mathrm{cos}{\displaystyle \frac{\pi s}{3}}+`$
$`{\displaystyle \frac{2}{81}}(\mathrm{cos}{\displaystyle \frac{2\pi s}{9}}+\mathrm{cos}{\displaystyle \frac{4\pi s}{9}}+\mathrm{cos}{\displaystyle \frac{8\pi s}{9}})`$
$`{\displaystyle \frac{1}{98}}(\mathrm{cos}{\displaystyle \frac{2\pi s}{7}}\mathrm{csc}{\displaystyle \frac{\pi }{7}}\mathrm{csc}{\displaystyle \frac{2\pi }{7}}+\mathrm{cos}{\displaystyle \frac{4\pi s}{7}}\mathrm{csc}{\displaystyle \frac{2\pi }{7}}\mathrm{csc}{\displaystyle \frac{3\pi }{7}}+\mathrm{cos}{\displaystyle \frac{6\pi s}{7}}\mathrm{csc}{\displaystyle \frac{3\pi }{7}}\mathrm{csc}{\displaystyle \frac{\pi }{7}}),`$
$`V(s,𝒮_{10})`$ $`=`$ $`{\displaystyle \frac{s^9}{1316818944000}}{\displaystyle \frac{11s^7}{12541132800}}+{\displaystyle \frac{113113s^5}{358318080000}}{\displaystyle \frac{\mathrm{sin}\pi s}{2949120}}s^4`$
$`{\displaystyle \frac{18063859s^3}{468202291200}}+s^2({\displaystyle \frac{1}{4374\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{2\pi s}{3}}+{\displaystyle \frac{143}{1179648}}\mathrm{sin}\pi s)+`$
$`s[{\displaystyle \frac{273512277643}{240789749760000}}+{\displaystyle \frac{1}{512\sqrt{2}}}\mathrm{cos}{\displaystyle \frac{\pi s}{2}}+{\displaystyle \frac{7}{13122}}\mathrm{cos}{\displaystyle \frac{2\pi s}{3}}+`$
$`{\displaystyle \frac{1}{625}}(\mathrm{cos}{\displaystyle \frac{4\pi s}{5}}\mathrm{cos}{\displaystyle \frac{2\pi s}{5}})]{\displaystyle \frac{2877523}{707788800}}\mathrm{sin}\pi s{\displaystyle \frac{1211}{52488\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{2\pi s}{3}}`$
$`{\displaystyle \frac{5}{1024\sqrt{2}}}\mathrm{sin}{\displaystyle \frac{\pi s}{2}}{\displaystyle \frac{1}{108}}\mathrm{sin}{\displaystyle \frac{\pi s}{3}}+{\displaystyle \frac{1}{64\sqrt{2}}}(\mathrm{csc}{\displaystyle \frac{3\pi }{8}}\mathrm{sin}{\displaystyle \frac{3\pi s}{4}}\mathrm{csc}{\displaystyle \frac{\pi }{8}}\mathrm{sin}{\displaystyle \frac{\pi s}{4}})+`$
$`{\displaystyle \frac{1}{50}}(\mathrm{sin}{\displaystyle \frac{3\pi s}{5}}\mathrm{sin}{\displaystyle \frac{\pi s}{5}}){\displaystyle \frac{2\sqrt{2}}{625}}({\displaystyle \frac{\sqrt{5}+2}{\sqrt{5+\sqrt{5}}}}\mathrm{sin}{\displaystyle \frac{2\pi s}{5}}+{\displaystyle \frac{\sqrt{5}2}{\sqrt{5\sqrt{5}}}}\mathrm{sin}{\displaystyle \frac{4\pi s}{5}})`$
$`{\displaystyle \frac{1}{196}}\mathrm{csc}{\displaystyle \frac{\pi }{7}}\mathrm{csc}{\displaystyle \frac{2\pi }{7}}\mathrm{csc}{\displaystyle \frac{3\pi }{7}}(\mathrm{sin}{\displaystyle \frac{6\pi s}{7}}+\mathrm{sin}{\displaystyle \frac{4\pi s}{7}}\mathrm{sin}{\displaystyle \frac{2\pi s}{7}})+`$
$`{\displaystyle \frac{1}{81}}(\mathrm{csc}{\displaystyle \frac{4\pi }{9}}\mathrm{sin}{\displaystyle \frac{8\pi s}{9}}+\mathrm{csc}{\displaystyle \frac{2\pi }{9}}\mathrm{sin}{\displaystyle \frac{4\pi s}{9}}+\mathrm{csc}{\displaystyle \frac{\pi }{9}}\mathrm{sin}{\displaystyle \frac{2\pi s}{9}}).`$
### 6.2 Coxeter groups.
Let us define two auxiliary functions
$`U_+(s,p,G)`$ $`=`$ $`V(s+p,G)+V(sp,G),`$
$`U_{}(s,p,G)`$ $`=`$ $`V(s+p,G)V(sp,G)`$ (52)
with obvious properties
$`U_+(s,p,𝐝^m/d_r)`$ $`=`$ $`U_{}(s,p+{\displaystyle \frac{d_r}{2}},𝐝^m)U_{}(s,p{\displaystyle \frac{d_r}{2}},𝐝^m),U_+(s,0,G)=2V(s,G),`$
$`U_{}(s,p,𝐝^m/d_r)`$ $`=`$ $`U_+(s,p+{\displaystyle \frac{d_r}{2}},𝐝^m)U_+(s,p{\displaystyle \frac{d_r}{2}},𝐝^m),U_{}(s,{\displaystyle \frac{d_r}{2}},𝐝^m)=V(s,𝐝^m/d_r),`$
where $`(m1)`$-$`tuple`$ $`\{𝐝^m/d_r\}=\{d_1,d_2,\mathrm{},d_{r1},d_{r+1},\mathrm{},d_m\}`$ doesn’t contain $`d_r`$-exponent.
Sylvester waves for the Coxeter groups are given below expressed through the relations elaborated in the previous Sections.
$`\underset{¯}{G=A_m}`$ $`,d_r=2,3,\mathrm{},m+1;\xi (A_m)=\frac{1}{4}m(m+3)`$
$`V(s,A_m)`$ $`=`$ $`U_{}(s,{\displaystyle \frac{1}{2}},𝒮_m).`$ (53)
$`\underset{¯}{G=B_m}`$ $`,d_r=2,4,6,\mathrm{},2m;\xi (B_m)=\frac{1}{2}m(m+1)`$
$`V(s,B_m)`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{\Psi }_2(s\xi (B_m))U_+({\displaystyle \frac{s}{2}},0,𝒮_m).`$ (54)
In the list for $`D_m`$ groups the degree $`m`$ occurs twice when $`m`$ is even. This is the only case involving such a repetition.
$`\underset{¯}{G=D_m}`$ $`,d_r=2,4,6,\mathrm{},2(m1),m,m3;\xi (D_m)=\frac{1}{2}m^2,`$
$`V(s,D_{2m})`$ $`=`$ $`\mathrm{\Psi }_2(s)U_+({\displaystyle \frac{s}{2}},{\displaystyle \frac{m}{2}},𝒮_{2m}),`$ (55)
$`V(s,D_{2m+1})`$ $`=`$ $`{\displaystyle \underset{s_1=0}{\overset{s\xi (D_{\mathit{2}m+\mathit{1}})}{}}}V(s+{\displaystyle \frac{2m+1}{2}}s_1,B_{2m})\mathrm{\Psi }_{2m+1}(s_1),`$
$`V(s,D_3)`$ $`=`$ $`V(s,A_3),`$
$`V(s,D_5)`$ $`=`$ $`U_{}(s,{\displaystyle \frac{11}{2}},𝒮_8)U_{}(s,{\displaystyle \frac{9}{2}},𝒮_8)U_{}(s,{\displaystyle \frac{5}{2}},𝒮_8)+U_{}(s,{\displaystyle \frac{3}{2}},𝒮_8).`$
$`\underset{¯}{G=G_\mathit{2}}`$ $`,d_r=2,6;\xi (G_\mathit{2})=4,`$
$`V(s,G_2)`$ $`=`$ $`\mathrm{\Psi }_2(s)U_{}({\displaystyle \frac{s}{2}},1,𝒮_3).`$ (56)
$`\underset{¯}{G=F_\mathit{4}}`$ $`,d_r=2,6,8,12;\xi (F_\mathit{4})=14,`$
$`V(s,F_4)`$ $`=`$ $`\mathrm{\Psi }_2(s)[U_+({\displaystyle \frac{s}{2}},{\displaystyle \frac{7}{2}},𝒮_6)U_+({\displaystyle \frac{s}{2}},{\displaystyle \frac{3}{2}},𝒮_6)].`$ (57)
$`\underset{¯}{G=E_\mathit{6}}`$ $`,d_r=2,5,6,8,9,12;\xi (E_\mathit{6})=21,`$
$`V(s,E_6)`$ $`=`$ $`U_+(s,18,𝒮_{12})U_+(s,17,𝒮_{12})U_+(s,15,𝒮_{12})+`$ (58)
$`U_+(s,13,𝒮_{12})+U_+(s,5,𝒮_{12})U_+(s,2,𝒮_{12}).`$
$`\underset{¯}{G=E_\mathit{7}}`$ $`,d_r=2,6,8,10,12,14,18;\xi (E_\mathit{7})=35,`$
$`V(s,E_7)`$ $`=`$ $`\mathrm{\Psi }_2(s1)[U_+({\displaystyle \frac{s}{2}},5,𝒮_9)U_+({\displaystyle \frac{s}{2}},3,𝒮_9)].`$ (59)
$`\underset{¯}{G=E_\mathit{8}}`$ $`,d_r=2,8,12,14,18,20,24,30;\xi (E_\mathit{8})=64,`$
$`V(s,E_8)`$ $`=`$ $`\mathrm{\Psi }_2(s)[U_{}({\displaystyle \frac{s}{2}},28,𝒮_{15})+U_{}({\displaystyle \frac{s}{2}},21,𝒮_{15})+U_{}({\displaystyle \frac{s}{2}},12,𝒮_{15})+`$ (60)
$`U_{}({\displaystyle \frac{s}{2}},11,𝒮_{15})U_{}({\displaystyle \frac{s}{2}},8,𝒮_{15})U_{}({\displaystyle \frac{s}{2}},7,𝒮_{15})`$
$`U_{}({\displaystyle \frac{s}{2}},6,𝒮_{15})U_{}({\displaystyle \frac{s}{2}},26,𝒮_{15})U_{}({\displaystyle \frac{s}{2}},25,𝒮_{15})].`$
$`\underset{¯}{G=H_\mathit{3}}`$ $`,d_r=2,6,10;\xi (H_\mathit{3})=9,`$
$`V(s,H_3)`$ $`=`$ $`\mathrm{\Psi }_2(s1)[U_+({\displaystyle \frac{s}{2}},3,𝒮_5)U_+({\displaystyle \frac{s}{2}},1,𝒮_5)].`$ (61)
$`\underset{¯}{G=H_\mathit{4}}`$ $`,d_r=2,12,20,30;\xi (H_\mathit{3})=32,`$
$`V(s,H_4)`$ $`=`$ $`U_+(s,32,E_8)U_+(s,24,E_8)U_+(s,18,E_8)U_+(s,14,E_8)+`$ (62)
$`U_+(s,10,E_8)U_+(s,8,E_8)+U_+(s,6,E_8)+U_+(s,0,E_8).`$
$`\underset{¯}{G=I_m}`$ $`,d_r=2,m;\xi (I_m)=1+\frac{1}{2}m`$
$`V(s,I_m)`$ $`=`$ $`{\displaystyle \underset{s_1=0}{\overset{s\xi (I_m)}{}}}\mathrm{\Psi }_2(s\xi (I_m)s_1)\mathrm{\Psi }_m(s_1),`$ (63)
$`V(s,I_\mathit{2})`$ $`=`$ $`V(s,B_\mathit{1}),V(s,I_\mathit{3})=V(s,A_\mathit{2}),V(s,I_\mathit{4})=V(s,B_\mathit{2}),`$
$`V(s,I_\mathit{5})`$ $`=`$ $`U_+(s,{\displaystyle \frac{7}{2}},A_\mathit{4})U_+(s,{\displaystyle \frac{1}{2}},A_\mathit{4}),`$
$`V(s,I_\mathit{6})`$ $`=`$ $`V(s,G_\mathit{2}),V(s,I_\mathit{8})=U_+(s,5,B_\mathit{4})U_+(s,1,B_\mathit{4})`$
$`V(s,I_{\mathit{10}})`$ $`=`$ $`U_{}(s,3,H_\mathit{3}),V(s,I_{\mathit{12}})=U_+(s,7,F_\mathit{4})U_+(s,1,F_\mathit{4}).`$
## 7 Acknowledgement
We would like to thank Prof. Z.Rudnik for the communication of the proof related to the asymptotic behaviour of the least common multiple lcm($`1,2,\mathrm{},m,N`$).
This research was supported in part by grants from the Tel Aviv University Research Authority and the Gileadi Fellowship program of the Ministry of Absorption of the State of Israel (LGF).
## Appendix A Asymptotic behaviour of $`\mathrm{𝗅𝖼𝗆}`$(1,2,…,N).
The arithmetical function least common multiple lcm($`1,2,\mathrm{},N`$) =$`(N)`$ of the series of the natural numbers takes a specific place among the other arithmetical functions. It can neither be represented as the Cauchy integral of the generating function with subsequent evaluation with Hardy-Ramanujan circle method like different partition functions $`p(N),q(N)`$, nor has it its genesis in Riemann’s Zeta-function like many arithmetical functions $`\mu (N),\nu (N),\varphi (N),d(N)`$. $`(N)`$ appears naturally in the theory of restricted partition numbers as periods of Sylvester waves in symmetric groups $`𝒮_N`$ .
Numerical calculations of $`\frac{1}{N}\mathrm{ln}[(N)]`$ in the range $`0<N<550\times 10^3`$ give an oscillating behaviour around 1 with asymptotic approach to this value (Fig. 1).
This enabled us to conjecture an asymptotic law
$$\underset{N\mathrm{}}{lim}\frac{\mathrm{ln}(N)}{N}=1.$$
(A1)
In the rest of this Appendix we give a proof of this statement. Before going to the proof we recall some facts of the prime number theory:
$`\mathrm{𝖥𝟣}.`$ The Prime Number Theorem (PNT)
$$\text{if}\pi (N)=\underset{p_iN}{}1,\text{then}\pi (N)\stackrel{N\mathrm{}}{}\frac{N}{\mathrm{ln}N}.$$
(A2)
where a sum is running over all primes $`p_i`$ up to $`N`$.
$`\mathrm{𝖥𝟤}.`$ Let us set after Chebyshev
$$\theta (N)=\underset{p_iN}{}\mathrm{ln}p_i,$$
(A3)
then PNT is equivalent to $`\theta (N)\stackrel{N\mathrm{}}{}N`$.
$`\mathrm{𝖥𝟥}.`$ The Rieman hypothesis is equivalent to
$$\theta (N)=N+O(\sqrt{N}\mathrm{ln}N)$$
(A4)
Now it follows
$`\mathrm{𝖫𝖾𝗆𝗆𝖺}𝖠.`$
$$\underset{N\mathrm{}}{lim}\frac{\mathrm{ln}(N)}{N}=1.$$
$`andassumingtheRiemanhypothesis`$
$$\mathrm{ln}(N)=N+O(\sqrt{N}\mathrm{ln}N).$$
$`\mathrm{𝖯𝗋𝗈𝗈𝖿}\mathrm{𝗈𝖿}\mathrm{𝖫𝖾𝗆𝗆𝖺}𝖠.`$
We write the prime decomposition of $`(N)`$ as $`(N)=p^{k_p}`$. Clearly , for a prime to divide $`(N)`$ it has to be at most less than $`N`$. Moreover the highest $`k`$ power of $`p`$ dividing one of the integers 1, 2, …, $`N`$ is
$$k_p=[\frac{\mathrm{ln}N}{\mathrm{ln}p}]$$
Thus we find
$$\mathrm{ln}(N)=\underset{p_iN}{}[\frac{\mathrm{ln}N}{\mathrm{ln}p}]\mathrm{ln}p$$
(A5)
To estimate $`\mathrm{ln}(N)`$, break the previous sum into two parts, one $`Q_1`$ coming from primes $`p\sqrt{N}`$ and the second $`Q_2`$ from primes $`\sqrt{N}pN`$ :
$$Q_1=\underset{p_i\sqrt{N}}{}[\frac{\mathrm{ln}N}{\mathrm{ln}p}]\mathrm{ln}p,Q_2=\underset{\sqrt{N}pN}{}[\frac{\mathrm{ln}N}{\mathrm{ln}p}]\mathrm{ln}p$$
For estimating $`Q_1`$, use $`[x]x`$ and so we find
$$Q_1\underset{p_i\sqrt{N}}{}\frac{\mathrm{ln}N}{\mathrm{ln}p}\mathrm{ln}p=\mathrm{ln}N\pi (\sqrt{N})2\sqrt{N}$$
by the PNT. For the second sum $`Q_2`$ note that if $`\sqrt{N}pN`$ then
$$1\frac{\mathrm{ln}N}{\mathrm{ln}p}<2$$
and hence its integer part is identically 1. Thus
$$Q_2=\underset{\sqrt{N}pN}{}1\mathrm{ln}p=\theta (N)\theta (\sqrt{N})$$
Since $`\theta (\sqrt{N})\sqrt{N}`$ , we obtain finally
$$\mathrm{ln}(N)=\theta (N)+\theta (\sqrt{N})$$
Our $`\mathrm{𝖫𝖾𝗆𝗆𝖺}𝖠`$ follows immediately from F2, F3. $`\mathrm{}`$
## Appendix B Derivation of Sylvester waves $`V(s,𝒮_4)`$ and $`V(s,𝒮_5)`$.
We will illustrate how do the formulas (37-50) work in the case of the symmetric groups $`𝒮_4`$ and $`𝒮_5`$.
We start with Sylvester wave $`V(s,𝒮_3)`$ taken from (51)
$$V(s,𝒮_3)=\frac{s^2}{12}\frac{7}{72}\frac{1}{8}\mathrm{cos}\pi s+\frac{2}{9}\mathrm{cos}\frac{2\pi s}{3}$$
(B1)
and with successive usage of the formulas (37) and (46) one can obtain
$`R_1^4(s)={\displaystyle \frac{1}{144}},R_2^4(s)=0,R_3^4(s)={\displaystyle \frac{1}{96}}(5+3\mathrm{cos}\pi s),_4^4(s)={\displaystyle \frac{2}{9\sqrt{3}}}\mathrm{sin}{\displaystyle \frac{2\pi s}{3}}.`$ (B2)
Now we will use the representation (49)
$$V(s,𝒮_4)=\underset{j=1}{\overset{3}{}}R_j^4(s)s^{4j}+_4^4(s)+\rho _1^4\mathrm{sin}\frac{\pi }{2}s+\rho _2^4\mathrm{sin}\pi s.$$
(B3)
Since $`V(s,𝒮_4)=W(s5,𝒮_4)`$ the variable $`s`$ takes only integer values what makes the last contribution in (B3) into the $`V(s,𝒮_4)`$ irrelevant. The unknown coefficient $`\rho _1^4`$ is determined with help of zeroes (28) of $`W(s,𝒮_4)`$
$$0=V(1,𝒮_4)=\underset{j=1}{\overset{3}{}}R_j^4(1)+_4^4(1)+\rho _1^4,\text{or}\rho _1^4=\frac{1}{8}$$
(B4)
Thus we arrive at the Sylvester wave $`V(s,𝒮_4)`$ presented in (51).
Repeating the same procedure with symmetric group $`𝒮_5`$ we find
$`R_1^5(s)`$ $`=`$ $`{\displaystyle \frac{1}{2880}},R_2^5(s)=0,R_3^5(s)={\displaystyle \frac{11}{1152}},R_4^5(s)={\displaystyle \frac{1}{64}}\mathrm{sin}\pi s,`$ (B5)
$`_5^5(s)`$ $`=`$ $`{\displaystyle \frac{475}{27648}}{\displaystyle \frac{2}{27}}\mathrm{cos}{\displaystyle \frac{2\pi s}{3}}+{\displaystyle \frac{1}{8\sqrt{2}}}\mathrm{cos}{\displaystyle \frac{\pi s}{2}}.`$
The representation (50) produces
$$V(s,𝒮_5)=\underset{j=1}{\overset{4}{}}R_j^5(s)s^{5j}+_5^5(s)+\rho _0^5+\rho _1^5\mathrm{cos}\frac{2\pi s}{5}+\rho _2^5\mathrm{cos}\frac{4\pi s}{5}.$$
(B6)
Since $`V(s,𝒮_5)=W(s\frac{15}{2},𝒮_5)`$ the variable $`s`$ has only half-integer values. By solving three linear equations $`V(\frac{1}{2},𝒮_5)=V(\frac{3}{2},𝒮_5)=V(\frac{5}{2},𝒮_5)=0`$ we find
$$\rho _0^5=\frac{217}{28800},\rho _1^5=\frac{2}{25},\rho _2^5=\frac{2}{25},$$
(B7)
which together with (B6) produces the Sylvester wave $`V(s,𝒮_5)`$ from (51).
## Appendix C Table of restricted partition numbers $`W(s,𝒮_m)`$.
In this Appendix we give the Table of the restricted partition numbers $`𝒫_m(s)=W(s,𝒮_m)m10`$ for $`s`$ running in the different ranges. One can verify that the content of this Table can be obtained with the help of the formulas (51).
| $`s`$ | | $`S_1`$ | $`S_2`$ | $`S_3`$ | $`S_4`$ | $`S_5`$ | $`S_6`$ | $`S_7`$ | $`S_8`$ | $`S_9`$ | $`S_{10}`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- | --- |
| 1 | | 1 | 1 | 1 | 1 | 1 | 1 | 1 | 1 | 1 | 1 |
| 2 | | 1 | 2 | 2 | 2 | 2 | 2 | 2 | 2 | 2 | 2 |
| 3 | | 1 | 2 | 3 | 3 | 3 | 3 | 3 | 3 | 3 | 3 |
| 4 | | 1 | 3 | 4 | 5 | 5 | 5 | 5 | 5 | 5 | 5 |
| 5 | | 1 | 3 | 5 | 6 | 7 | 7 | 7 | 7 | 7 | 7 |
| 6 | | 1 | 4 | 7 | 9 | 10 | 11 | 11 | 11 | 11 | 11 |
| 7 | | 1 | 4 | 8 | 11 | 13 | 14 | 15 | 15 | 15 | 15 |
| 8 | | 1 | 5 | 10 | 15 | 18 | 20 | 21 | 22 | 22 | 22 |
| 9 | | 1 | 5 | 12 | 18 | 23 | 26 | 28 | 29 | 30 | 30 |
| 10 | | 1 | 6 | 14 | 23 | 30 | 35 | 38 | 40 | 41 | 42 |
| 51 | | 1 | 26 | 243 | 1215 | 4033 | 9975 | 19928 | 33940 | 51294 | 70760 |
| 52 | | 1 | 27 | 252 | 1285 | 4319 | 10829 | 21873 | 37638 | 57358 | 79725 |
| 53 | | 1 | 27 | 261 | 1350 | 4616 | 11720 | 23961 | 41635 | 64015 | 89623 |
| 54 | | 1 | 28 | 271 | 1425 | 4932 | 12692 | 26226 | 46031 | 71362 | 100654 |
| 55 | | 1 | 28 | 280 | 1495 | 5260 | 13702 | 28652 | 50774 | 79403 | 112804 |
| 56 | | 1 | 29 | 290 | 1575 | 5608 | 14800 | 31275 | 55974 | 88252 | 126299 |
| 57 | | 1 | 29 | 300 | 1650 | 5969 | 15944 | 34082 | 61575 | 97922 | 141136 |
| 58 | | 1 | 30 | 310 | 1735 | 6351 | 17180 | 37108 | 67696 | 108527 | 157564 |
| 59 | | 1 | 30 | 320 | 1815 | 6747 | 18467 | 40340 | 74280 | 120092 | 175586 |
| 60 | | 1 | 31 | 331 | 1906 | 7166 | 19858 | 43819 | 81457 | 132751 | 195491 |
| 101 | | 1 | 51 | 901 | 8262 | 48006 | 198230 | 628998 | 1621248 | 3539452 | 6757864 |
| 102 | | 1 | 52 | 919 | 8505 | 49806 | 207338 | 662708 | 1719877 | 3778074 | 7254388 |
| 103 | | 1 | 52 | 936 | 8739 | 51649 | 216705 | 697870 | 1823402 | 4030512 | 7782608 |
| 104 | | 1 | 53 | 954 | 8991 | 53550 | 226479 | 734609 | 1932418 | 4297682 | 8345084 |
| 105 | | 1 | 53 | 972 | 9234 | 55496 | 236534 | 772909 | 2046761 | 4580087 | 8942920 |
| 106 | | 1 | 54 | 990 | 9495 | 57501 | 247010 | 812893 | 2167057 | 4878678 | 9578879 |
| 107 | | 1 | 54 | 1008 | 9747 | 59553 | 257783 | 854546 | 2293142 | 5194025 | 10254199 |
| 108 | | 1 | 55 | 1027 | 10018 | 61667 | 269005 | 898003 | 2425678 | 5527168 | 10971900 |
| 109 | | 1 | 55 | 1045 | 10279 | 63829 | 280534 | 943242 | 2564490 | 5878693 | 11733342 |
| 110 | | 1 | 56 | 1064 | 10559 | 66055 | 292534 | 990404 | 2710281 | 6249733 | 12541802 |
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# The contribution of Narrow-Line Seyfert 1 galaxies to the soft X-ray background
## 1 The ROSAT Deep Surveys in the Lockman Hole
The most sensitive ROSAT surveys consist of a 207 ksec ROSAT PSPC exposure, a 205 ksec HRI raster scan and a total 1112 ksec HRI expsoure of a $`0.3deg^2`$ area in the Lockman Hole region. The HRI images are the basis for the Ultradeep HRI Survey (Hasinger at al. 1999). At a flux limit of 10<sup>-15</sup> erg s<sup>-1</sup> in the 0.5-2.0 keV energy band, the HRI survey has resolved about 70-80% of the soft X-ray background into discrete sources. The ROSAT Deep Survey (RDS), based on the PSPC image, includes a statistically complete sample of 50 X-ray sources with fluxes in the 0.5-2.0 keV band greater than 5.5$``$10<sup>-15</sup> erg s<sup>-1</sup> (Hasinger et al. 1998). The spectroscopic identification of the RDS using the Keck and Palomar telescopes have shown that about 75% of the sources are quasars and Seyfert galaxies (Schmidt et al. 1998, Lehmann et al. 2000).
Both surveys contain 91 X-ray sources, where the faintest sources reach a flux of 1.2$``$10<sup>-15</sup> erg s<sup>-1</sup> in the 0.5-2.0 keV band. Recent optical/infrared work has led to an identification of 88 of the 91 X-ray sources, confirming a high fraction of AGNs (72 objects). This is the largest fraction of AGNs found in any previous X-ray survey (Boyle et al. 1995, Georgantopolous et al. 1996, Bower et a. 1996 and McHardy et al. 1998). Among our AGNs is the most distant X-ray selected quasar at a redshift of 4.45 (Schneider et al. 1998). Groups and clusters of galaxies ($``$10%) form the second most abundant class of objects. One X-ray source has been classified as a narrow emission line galaxy (NELG), whereas some deep ROSAT PSPC surveys found a significantly larger fraction of NELGs (Boyle et al. 1995, McHardy et al. 1998). We see no evidence that NELGs or other classes of objects dominate the soft X-ray counts at faint fluxes.
## 2 Optical spectral properties of RDS quasars and Seyfert galaxies
In the following we discuss the emission line properties of the 41 spectroscopically confirmed RDS quasars and Seyfert galaxies. 33 X-ray sources have been identified with quasars and Seyfert 1 galaxies in the redshift range between 0.08 and 2.83. Their optical spectra show broad emission lines (FWHM $`>`$ 1500 km s<sup>-1</sup>) of Ly$`\alpha `$ $`\lambda `$1216, C IV $`\lambda `$1548, C III\] $`\lambda `$1908 and Mg II $`\lambda `$2798 at medium and large redshifts, or of Balmer lines (H$`\alpha `$ $`\lambda `$6563, H$`\beta `$4861) at lower redshifts. Most of them can be classified as type I AGN, whereas the large Balmer decrement of two Seyfert galaxies having broad H$`\alpha `$ or H$`\beta `$ emission lines indicates type II AGN (see object 59A in Fig. 1).
The optical spectra of six RDS AGNs show only narrow emission lines with FWHM $`<`$ 1500 km s<sup>-1</sup> indicating a type II AGN. When their redshifts are above 0.5, the spectra do not allow a classification using the diagnostic diagrams of Osterbrock (1981). But the existence of high excitation \[Ne V\] $`\lambda `$3426 (see 26A in Fig. 2) or strong \[Ne III\] $`\lambda `$3868 narrow emission lines, together with an X-ray luminosity above 10<sup>43</sup> erg s<sup>-1</sup> in the 0.5-2.0 keV band, reveal an AGN (Schmidt et al. 1998). Two further sources have been classified as type II AGNs, because of their high X-ray luminosity (log L$`{}_{X}{}^{}>43`$). One of them is a radio galaxy at $`z=0.708`$ showing only typical galaxy absorption lines. Several RDS AGNs (type I and II) below $`z=1`$ show a significant continuum emission originating in the host galaxy (Fig. 1/Fig. 2). The RDS sample contains in total 35 type I and 6 type II AGNs. The large $`RK^{}`$ colour of two spectroscopically unidentified sources indicates either obscured AGNs or high-redshift clusters of galaxies (Lehmann et al. 2000).
## 3 Spectroscopic discrimination of Seyfert galaxies and field NELGs
In the course of the spectrocopic identification of our X-ray sources we have taken optical spectra of 83 field galaxies. 67 of them are narrow emission line galaxies (NELGs) with no physical connection to the X-ray sources. There is no soft X–ray emission (0.5-2.0 keV band) detected above a limiting flux of 10<sup>-15</sup> erg cm<sup>-2</sup> s<sup>-1</sup> associated with these sources. The optical spectra of the field NELG (cf. 14C in Fig. 2) show only narrow emission lines, but no strong \[Ne III\] $`\lambda `$3868 or high ionization \[Ne V\] $`\lambda `$3426 emission lines. The possible misidentification of faint X-ray sources with NELGs is discussed by Schmidt et al. (1998) and Lehmann et al. (2000).
To study the emission line properties we have derived the FWHM and the rest frame EW for the most prominent emission lines of the RDS quasars/Seyfert galaxies and of the field NELGs. The mean FWHM/EW values for the broad emission lines of RDS quasars/Seyfert 1 galaxies are consistent with those values found for other X-ray selected AGN samples at lower mean redshift, eg., the RIXOS-sample (Puchnarewicz et al. 1997) and the CRSS-sample (Boyle et al. 1997), and the mean values derived from several emisssion line and UV/optical selected AGN samples (Schmidt et al. 1986).
The mean RDS EW values of the forbidden narrow emission lines, found for Seyfert galaxies at redshifts below 1, are slightly smaller than those from the AGN comparison samples (Lehmann et al. 2000). The RDS EW of \[Ne III\] $`\lambda `$3868 and the \[O III\] $`\lambda \lambda `$4959/5007 lines seem to be in better agreement with those from field NELGs (Tab. 1) than with those from other AGN samples. The optical spectra of most RDS AGNs, which contain forbidden narrow emission lines, show several typical galaxy absorption lines, eg., Ca H$`+`$K $`\lambda \lambda `$3934/3968, CH G $`\lambda `$4301 (see 26A in Fig. 2). A strong continuum contribution produced by the host galaxy could account for these less stronger forbidden lines in our RDS AGNs.
## 4 Contribution of NLS1 to the soft X-ray background
Although the majority of our faint X-ray sources have been spectroscopically identified with quasars and Seyfert galaxies, only one out of 69 AGNs matches the optical criteria of narrow line Seyfert 1 galaxies taken from Osterbrock & Pogge (1985) and more precisely defined by Goodrich (1989). These somewhat subjective criteria are as follows. 1) NLS1 have slightly broader Balmer lines ($`<`$ 2000 km s<sup>-1</sup>) in comparison to the forbidden lines such as \[O II\] $`\lambda `$3727 or \[O III\] $`\lambda \lambda `$4959/5007. 2) The flux ratio \[O III\] $`\lambda `$5007 to H$`\beta `$ $`\lambda `$4861 is $`<`$ 3, a level which allows to discriminate between NLS1 and Seyfert 2 galaxies (Shuder & Osterbrock 1981). 3) Emission lines of Fe II or of high ionization \[Fe VII\] $`\lambda `$6087 and \[Fe X\] $`\lambda `$6375 are often present.
The spectrum of the NLS1 37A (Fig. 3) shows prominent Fe II bumps at 4500–4600 Å and at 5250–5350 Å (in the rest frame). The FWHM of the broad component of the H$`\beta `$ $`\lambda `$4861 emission line is slightly above the limit for NLS1 (2070$`\pm `$20 km s<sup>-1</sup>). The flux ratio \[O III\] $`\lambda `$5007/H$`\beta `$ $`\lambda `$4861 is clearly below 3. The NLS1 37A at $`z=0.462`$ belongs to the highest-redshift objects of this class. The bright soft X–ray selected sample of Grupe et al. (1999) contains for example only 4 objects (12%) above $`z=0.3`$.
The fraction of NLS1 in X-ray selected AGN samples (see for an overview Grupe in these proceedings) has to be considered with caution. NLS1 at redshifts above 0.65 could be missed due to the limited wavelength coverage of the optical spectra ($`\lambda <`$ 8200 Å). The H$`\beta `$/\[O III\] region is covered in 13 of 69 spectroscopically identified AGNs from the ROSAT Deep Surveys. Most of our AGNs show only broad UV emission lines, eg., C III\] $`\lambda `$1908 or Mg II $`\lambda `$2798, at high redshifts. Recently, Rodriguez-Pascual et al. (1997) have detected broad components of strong UV emission lines (eg., L$`\alpha `$ $`\lambda `$1216 and C IV $`\lambda `$1548) in several NLS1. If this is a common property we can not exclude a significantly larger fraction of NLS1 in high-redshift AGN samples.
But even the fraction of NLS1 in some low-redshift AGN samples, eg. 7% of the ROSAT Bright Survey (Schwope et al. 2000) or 22% of the RIXOS AGN sample (Puchnarewicz et al. 1997), could result from a different classification scheme of NLS1 objects (upper limit for FWHM$`{}_{H\beta }{}^{}<1500`$ or 2000 km s<sup>-1</sup>). Without well-controlled samples at low and high redshifts it is not possible to obtain reliable statements about the cosmological evolution of NLS1.
The fraction of NLS1 found in the ROSAT Deep Surveys ($``$1%) may indicate a marginal contribution of NLS1 to the soft X-ray background. Nevertheless we have to consider this as a lower limit so far.
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# Formation of Few-Body Clusters in Nuclear Matter
To describe the formation of clusters – in the general case a nonequilibrium process – in an interacting many-body system constitutes a new challenge for few-body methods. This may happen when the residual interaction between the quasi-particles leads to correlations. An example for such a system is nuclear matter. In the laboratory finite pieces of nuclear matter can be produced in heavy ion collisions. In astrophysics nuclear matter occurs e.g. during the supernova collapse and the formation of a neutron star.
A microscopic approach to treat the formation of clusters uses a generalized quantum Boltzmann equation. This coupled equation has been numerically solved for nucleon $`f_N`$, deuteron $`f_d`$, triton $`f_t`$, and helium-3 $`f_h`$ distributions utilizing the Boltzmann-Uehling-Uhlenbeck (BUU) approach . The coupling between the different species is through the collision integrals $`𝒦[f_N,f_d,\mathrm{}]`$. For the deuteron loss $`𝒦_d^{\mathrm{out}}(P,t)`$, e.g., it is
$`𝒦_d^{\mathrm{out}}(P,t)`$ $`=`$ $`{\displaystyle d^3kd^3k_1d^3k_2d^3k_3|k_1k_2k_3|U_0|kP|_{dNpnN}^2}`$ (1)
$`\times \overline{f}_N(k_1,t)\overline{f}_N(k_2,t)\overline{f}_N(k_3,t)f_N(k,t)+\mathrm{}`$
where $`\overline{f}_N=(1f_N)`$. The ellipsis denote further possible contributions, e.g. $`ddtp`$, $`ddhp`$ or processes like $`\gamma dnp`$, etc. The quantity $`U_0`$ is the $`NdNNN`$ break-up transition operator that in general depends on the medium. This dependence is neglected, if experimental cross sections are used to replace $`U_0`$ in Eq. (1) – a standard technique and in many cases very successful. To calculate $`U_0`$ including the self energy shift and the proper Pauli blocking and study the influence of the medium on different observablesa generalized Alt-Grassberger-Sandhas (AGS) equation has been derived earlier . The effective few-body problem in matter arises in the Green function approach along with a cluster mean-field expansion or the a Dyson equation approach .
Using the effective equations derived elsewhere we may study two important effects of the medium on the effective few-body-systems embedded in nuclear matter: (1) The change of binding energy, i.e. the self energy shift, (2) the changes in the reactions rates. Both effects are important and have consequences for the simulation of heavy ion collisions.
The change of binding energy eventually leads to the Mott effect (where $`E_{\mathrm{bound}}0^{}`$). Not as dramatic as in Coulombic systems – here the Mott effect leads to the transition from isolating to conducting phase – it, however, influences the number of clusters and the energy spectrum produced in a heavy ion collision. The Mott density depends on the momentum of the cluster as for higher momenta blocking of the constituents of the cluster is less effective. The Mott momenta for deuteron and triton are shown in Fig. 2 for two different temperatures. Tritons are more stable than deuterons and at lower densities both clusters a more stable for higher temperatures.
For a typical temperature of the heavy ion collision (in the final stage) we have calculated the in-medium cross section. The result is shown in Fig. 2. The threshold shift is because of the change of the deuteron’s binding energy. A strong enhancement of the maximum of the cross section appears, however less change at higher energies.
As a consequence the reaction time scales become much faster, when in-medium rates are used in the calculation instead of isolated ones. Fig. 4 shown the deuteron break-up time $`\tau _{\mathrm{bu}}(P,n_N)`$ evaluated in linear response, where the life time of deuteron fluctuations $`\delta f_d(P,t)`$ depends on the deuteron momentum $`P`$ and the nuclear density . Similar, a chemical relaxation time $`\tau _{\mathrm{rel}}(n_N)`$ can be defined, which results from a linearization of the respective rate equations . The relaxation times are shown in Fig. 4.
Finally considering a specific heavy ion collision it is possible to calculate the total number of deuterons coming out of a central collision of <sup>129</sup>Xe on <sup>119</sup>Sn at 50 MeV/A . Fig. 6 shows the integrated number of deuterons taking into account gain and loss terms induced by the reactions in the medium. The net effect is a significant enhancement of the number of deuterons. Fig. 6 shows the influence of using in-medium rates on the spectrum of the proton to deuteron ratio along wuth experimental data.
From the analysis is becomes clear that the medium modified elementary cross section and the proper self energy correction (binding energy shift, Mott effect) of the clusters should be included in the simulation of heavy ion collisions at that energies.
Acknowledgment: It a pleasure for me to thank P. Danielewicz, C. Kuhrts, G. Röpke, W. Schadow, and P. Schuck for substantial contributions. This work and presentation has been supported by the Deutsche Forschungsgemeinschaft.
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# A non trivial extension of the two-dimensional Ising model: the 𝑑-dimensional “molecular” model.
## I Introduction
The molecular model has been first proposed as a very simple $`d`$-dimensional lattice model which incorporates some degree of frustration and thus describes some aspects of molecular orientation in covalently bound molecular solids. Most molecular liquids retain their molecular structure even in the solid phase, where some long range order usually shows up as a consequence of inter-molecular interaction. However in the solid the orientational order of the molecules may change according to the thermodynamic conditions giving rise to quite rich phase diagrams as recently observed for hydrogen under high pressure. Moreover orientational ordering is responsible for several phase transitions occurring even in the liquid phase (liquid crystals), and lattice models have been found to be useful for describing such transitions.
The most studied models describe a molecular interaction which arises from dipole fluctuation, is weak and gives rise to the observed three-dimensional ordering of most molecular Van der Waals solids. The thermodynamic behaviour of such weakly interacting systems can be analyzed in terms of $`O(3)`$ symmetric vectorial models.
Conversely the molecular model was motivated by a description of the almost-covalent molecular solids where the interaction has a covalent main component and is characterized by some level of frustration (since the coordination number for the covalent bond is quite low). In such solids each molecule must choose a few partners and cannot accept any further invitation. The lower is the allowed coordination number, the higher the frustration which gives rise to the low-dimensional structures observed in polymers (one-dimensional) or in iodine and hydrogen halides (two-dimensional). In iodine, where the covalent nature of the interaction is out of doubt, even one-dimensional zigzag chains have been reported under pressure. Moreover we expect that a covalent interaction should show up for all the molecular solids under high pressure, as the inter-molecular distance approaches the intra-molecular length, provided that some important structural transition does not occur first (like dissociation). Quite recently the existence of charge transfer between $`O_2`$ molecules has also been reported under pressure
The molecular model is a simple frustrated lattice model which can describe some aspects of a covalently bound molecular solid. It consists of a $`d`$-dimensional hypercubic lattice with a randomly oriented linear molecule at each site. In its simplest version each molecule is only allowed to be oriented towards one of its nearest neighbours. There is an energy gain for any pair of neighbours which are oriented along their common bond (a covalent bond). The existence of preferred orientational axes breaks the rotational invariance of the single molecule as it is likely to occur for any real molecular system under pressure. In fact even in hydrogen, the broken symmetry phase transition which is observed under pressure has been recently shown to be affected by the presence of a crystal field which breaks the isotropy.
Quite recently, similar lattice models have been used for describing the diffusion of particles and molecules inside a polymer, and the growth of one-dimensional islands (polymeric chains). The molecular model has already stimulated some recent work on molecular orientation in nitrogen which goes back to the phenomenology laid down by Pople and Karasz. It had been argued that the weak intermolecular bonds between two $`N_2`$ molecules should be favourable for the formation of an orientationally disordered “plastic crystal” solid phase, and should lead to freezing into an orientationally ordered phase. More recent experimental data on nitrogen confirm the existence of an orientational disordering temperature in the solid below the melting temperature. However, as far as we know, the molecular systems which are more closely described by the molecular model are the hydrogen halides $`HX(X=F,Cl,Br,I)`$. Their low-temperature structures are known to consist of planar chains of molecules in the condensed state while a totally disordered structure has been observed with increasing temperature at ambient pressure. Moreover the opposite transition, from orientational disorder to an ordered chain structure, has also been reported by increasing pressure.
The molecular model undergoes a transition from an high-temperature (or weakly interacting) fully isotropic disordered system, to a low-temperature (or strongly interacting) anisotropic low dimensional broken-symmetry phase. As a consequence of frustration the breaking of symmetry is accompanied by a sort of decomposition of the system in low-dimensional almost independent parts, as observed in solid iodine and hydrogen halides. Such remarkable behaviour requires a space dimension $`d>2`$, while for $`d=2`$ the model is shown to be equivalent to the exactly solvable two-dimensional Ising model. As shown by Monte Carlo calculations, in the broken-symmetry phase the system displays the presence of correlated chains of molecules (polymers) which point towards a common direction inside each two-dimensional sub-set of the lattice (plane). Such planes are weakly correlated in the low-temperature phase, and the system has a two-dimensional behaviour even for $`d>3`$.
In this paper the relevance of the molecular model as a non-trivial extension of the Ising model is pointed out. Thus, apart from the physical motivations, the model is fully examined and the phase transition is described by several methods: mean-field, real space renormalization group and numerical simulations. Exactly solvable models are important for our understanding of more complex systems, and provide a test for approximate techniques. The $`d`$-dimensional molecular model shares with the Ising model the $`d=2`$ realization, since their equivalence for $`d=2`$ has been proven to be exact. In this paper we will focus on the $`d=3`$ model, but we will take advantage of the existence of an exactly solvable realization for $`d=2`$. For $`d>2`$, as the frustration increases, the model shows a very different behaviour compared to the Ising or Potts models. These last show a fully $`d`$-dimensional broken-symmetry phase while the molecular model is characterized by a low-dimensional ordering inside the planes with negligible correlation among different planes. Moreover for $`d=3`$ the molecular model is shown to belong to a different universality class, since its critical exponent $`\nu `$ turns out to be $`\nu =0.44\pm 0.02`$ by finite size scaling. We expect that such new universality class should describe a broad group of isotropic physical systems characterized by a low-dimensional ordering in their low-temperature phase. Such broad class of phase transitions should be explored by experiments in order to compare with the theoretical predictions for the critical universal properties. In such respect the driving parameter does not need to be the temperature, as the bond strength can be directly modified by a change of pressure in several systems.
This paper is organized as follows. Section II contains a formal definition of the $`d`$-dimensional molecular model, and a proof of its equivalence to the Ising model for $`d=2`$. In Section III the mean-field solution is discussed for the generic $`d`$-dimensional model. In Section IV a modified variational Migdal-Kadanoff method is presented and its application is discussed for $`d=2`$ and $`d=3`$. At variance with a previous calculation which yielded a quite poor result, the variational method is shown to work very well provided that some assumptions are made on the nature of the broken phase. Section V contains the results of a numerical simulation by Monte Carlo sampling, and the numerical estimation of both critical temperature and exponent by finite size scaling. In Section VI the main findings are summarized and discussed.
## II The molecular model
The molecular model was first described in Ref.. We briefly describe its formal definition in order to fix the notation. Let us consider a $`d`$-dimensional hypercubic lattice, with a randomly oriented linear molecule at each site. The molecules are supposed to be symmetric with respect to their centre of mass which is fixed at the lattice site. Only a discrete number of space orientations are allowed for each molecule: we assume that each of them must point towards one of its $`2d`$ first neighbour sites. This choice can be justified by the existence of covalent interactions along preferred axes. Then each molecule has $`d`$ different states corresponding to molecular orientation along the hypercube axes (molecules are symmetric). Finally, each couple of first neighbour molecules, when pointing the one towards the other, are assumed to gain a bonding energy for their directional covalent bond (they touch each other). As shown in Fig.1 for $`d=2`$, bonding in a direction excludes any possible bond along the other $`(d1)`$ directions. The coordination number is $`2`$ for any value of $`d`$, and the frustration increases with increasing $`d`$.
According to such description we introduce a versor variable $`\widehat{w}_𝐫`$ for each of the $`N`$ sites $`𝐫`$ of the lattice, with $`\widehat{w}_𝐫\{\widehat{x}_1,\widehat{x}_2,\mathrm{}\widehat{x}_d\}`$ pointing towards one of the $`d`$ hypercube axes $`x_\alpha `$. The versors $`\widehat{x}_\alpha `$ are assumed to be orthonormal: $`\widehat{x}_\alpha \widehat{x}_\gamma =\delta _{\alpha \gamma }`$. The partition function follows
$$Z=\underset{\{\widehat{w}\}}{}e^S=\underset{\{\widehat{w}\}}{}\mathrm{exp}\left[4\beta \underset{𝐫,\alpha }{}(\widehat{w}_𝐫\widehat{x}_\alpha )(\widehat{w}_{𝐫+\widehat{x}_\alpha }\widehat{x}_\alpha )\right]$$
(1)
where $`\{\widehat{w}\}`$ indicates a sum over all configurations, $`\alpha `$ runs from $`1`$ to $`d`$, and the lattice spacing is set to unity. The inverse temperature $`\beta `$ (in units of binding energy) can be negative for a repulsive model, but is assumed positive in the molecular context.
The model may be generalized by introducing an external $`d`$-dimensional vectorial field $`𝐡(\alpha )`$ at each link. The dependence on $`\alpha `$ means that the field differs according to the space direction $`\alpha `$ of the lattice link which joins the sites. The modified partition function reads
$$Z_h=\underset{\{\widehat{w}\}}{}e^{S_h}=\underset{\{\widehat{w}\}}{}\mathrm{exp}\left\{4\beta \underset{𝐫,\alpha }{}\left[\left(\widehat{w}_𝐫\widehat{x}_\alpha \right)\left(\widehat{w}_{𝐫+\widehat{x}_\alpha }\widehat{x}_\alpha \right)+𝐡(\alpha )\widehat{w}_𝐫+𝐡(\alpha )\widehat{w}_{𝐫+\widehat{x}_\alpha }\right]\right\}$$
(2)
It is evident that if the field satisfies the condition
$$\underset{\alpha }{}𝐡(\alpha )=0$$
(3)
then $`S_h`$ does not depend on $`𝐡`$ and $`S_hS`$. In such case the extra degree of freedom provided by $`𝐡`$ can be regarded as a sort of internal symmetry of the model. This global symmetry can be made local by allowing the field $`𝐡`$ to depend on the site position $`𝐫`$. We will only take advantage of the global symmetry in this paper. We notice that such symmetry cannot be seen as a gauge invariance, since in lattice gauge models any gauge change leaves the energy gain unchanged at any link. Here the field $`𝐡`$ changes the energy gain of all the links while the whole action is invariant.
Adopting a more compact notation, the partition function reads
$$Z_h=\underset{\{\widehat{w}\}}{}e^{S_h}=\underset{\{\widehat{w}\}}{}e^{_{𝐫,\alpha }(𝐫,\alpha )}$$
(4)
where the Lagrangian density $``$ follows as
$$(𝐫,\alpha )=\widehat{w}_𝐫^{}M_\alpha (\beta ,𝐡)\widehat{w}_{𝐫+\widehat{x}_\alpha }$$
(5)
Here the canonical $`d`$-dimensional column vector representation of $`R^d`$ is employed with $`\widehat{x}_1(1,0,0\mathrm{})`$, $`\widehat{x}_2(0,1,0\mathrm{})`$, etc. The $`d\times d`$ matrix $`M_\alpha `$ does not depend on the configurations of the system, and entirely characterizes the model.
The global symmetry of the action provides a simple way to show the equivalence between molecular and Ising models for $`d=2`$. For the two-dimensional lattice the condition (3) is satisfied by the field $`𝐡(1)=h(\widehat{x}_1\widehat{x}_2)`$, $`𝐡(2)=𝐡(1)`$. The matrix $`M_\alpha `$ follows
$`M_1`$ $`=`$ $`\left(\begin{array}{cc}4\beta (1+2h)& 0\\ 0& 8\beta h\end{array}\right),`$ (6)
$`M_2`$ $`=`$ $`\left(\begin{array}{cc}8\beta h& 0\\ 0& 4\beta (1+2h)\end{array}\right).`$ (7)
Then for $`h=1/4`$, $`M_1M_2`$, and $``$ reads
$$L(𝐫,\alpha )=\beta +\widehat{w}_𝐫^{}\left(\begin{array}{cc}\beta & \beta \\ \beta & \beta \end{array}\right)\widehat{w}_{𝐫+\widehat{x}_\alpha }.$$
(8)
Identifying the two-dimensional column versors $`\widehat{w}`$ with spin variables, apart from an inessential factor, $`Z`$ reduces to the partition function of a two-dimensional Ising model
$$Z=e^{2\beta N}Z_{Ising}$$
(9)
and is exactly solvable. For $`\beta +\mathrm{}`$ a ground state is approached with all the molecules oriented along the same direction, and with formation of one-dimensional polymeric chains (Fig.2a); for $`\beta \mathrm{}`$ the repulsive model approaches a zero-energy (no bonds) ground state analogous to the antiferromagnetic configuration of the Ising model (Fig.2b).
For $`d3`$ the analogy with the Ising model breaks down, and this is evident from a simple analysis of the ground state configuration. Due to frustration the model has an infinitely degenerate ground state in the thermodynamic limit $`N\mathrm{}`$. For instance, in the case $`d=3`$, the minimum energy is obtained by orienting all the molecules along a common direction, as for $`d=2`$. However the ground state configuration is not unique: the number of molecular bonds does not change if we rotate together all the molecules belonging to an entire layer which is parallel to the original direction of orientation. As a consequence of frustration the total degeneration is $`3(2^{(N^{1/3})})`$, and the system could even behave like a glass for the large energy barriers which separate each minimum from the other. The phase diagram is expected to be quite rich, with at least a transition point between the high temperature disordered phase and an ordered broken-symmetry low temperature phase.
## III Mean Field Approximation
For the generic $`d`$-dimensional model, some analytical results can be obtained in Mean-Field (MF) approximation: neglecting second order fluctuation terms
$$(\widehat{w}_𝐫\widehat{x}_\alpha )(\widehat{w}_{𝐫+\widehat{x}_\alpha }\widehat{x}_\alpha )\mathrm{\Delta }_\alpha (\widehat{w}_𝐫\widehat{x}_\alpha )+\mathrm{\Delta }_\alpha (\widehat{w}_{𝐫+\widehat{x}_\alpha }\widehat{x}_\alpha )\mathrm{\Delta }_\alpha ^2$$
(10)
where $`\mathrm{\Delta }_\alpha =\widehat{w}_𝐫\widehat{x}_\alpha `$ is an average over the configurations, and $`_\alpha \mathrm{\Delta }_\alpha =1`$ (with the obvious bounds $`0\mathrm{\Delta }_\alpha 1`$). Here the order parameter $`\mathrm{\Delta }_\alpha `$ gives the probability of finding a molecule oriented along the direction of $`\widehat{x}_\alpha `$. The partition function factorizes as
$$Z_{MF}=\left(\underset{\alpha }{}e^{8\beta \mathrm{\Delta }_\alpha }\right)^N\mathrm{exp}\left(4N\beta \underset{\alpha }{}\mathrm{\Delta }_\alpha ^2\right)$$
(11)
and the free energy follows
$$F_{MF}=\frac{1}{N\beta }\mathrm{log}Z_{MF}=4\underset{\alpha }{}\mathrm{\Delta }_\alpha ^2\frac{1}{\beta }\mathrm{log}\left(\underset{\alpha }{}e^{8\beta \mathrm{\Delta }_\alpha }\right).$$
(12)
The derivative with respect to $`\mathrm{\Delta }_\mu `$ yields, for the stationary points,
$$\mathrm{\Delta }_\mu =\frac{e^{8\beta \mathrm{\Delta }_\mu }}{_\alpha e^{8\beta \mathrm{\Delta }_\alpha }}$$
(13)
which satisfies the condition $`_\alpha \mathrm{\Delta }_\alpha =1`$.
In the high temperature limit $`\beta 0`$ eq.(13) has the unique solution $`\mathrm{\Delta }_\mu =1/d`$ which reflects the complete random orientation of molecules. In the opposite limit $`\beta \mathrm{}`$, apart from such solution, eq.(13) is satisfied by the broken-symmetry field $`\mathrm{\Delta }_\mu =1`$, $`\mathrm{\Delta }_\alpha =0`$ for $`\alpha \mu `$, which obviously corresponds to a minimum for $`F_{MF}`$. Then at a critical point $`\beta =\beta _c`$ the high temperature solution must become unstable towards a multivalued minimum configuration. The Hessian matrix is easily evaluated at the stationary points by using eqs.(13) and (12):
$$H_{\mu \nu }=\frac{1}{8}\frac{^2F_{MF}}{\mathrm{\Delta }_\mu \mathrm{\Delta }_\nu }=\delta _{\mu \nu }\left(18\beta \mathrm{\Delta }_\mu \right)+8\beta \mathrm{\Delta }_\mu \mathrm{\Delta }_\nu $$
(14)
In the high temperature phase ($`\beta <\beta _c`$), inserting $`\mathrm{\Delta }_\mu =1/d`$, the eigenvalue problem
$$det|H_{\mu \nu }\lambda \delta _{\mu \nu }|=0$$
(15)
yields
$$\left(1\frac{8\beta }{d}\lambda \right)^{d1}\left(1\lambda \right)=0.$$
(16)
Thus the Hessian matrix is positive defined if and only if $`\lambda =(18\beta /d)>0`$. Beyond the critical point $`\beta =\beta _c=(d/8)`$ the solution $`\mathrm{\Delta }_\mu =1/d`$ is not a minimum, and a multivalued minimum configuration shows up. Such result obviously agrees with the MF prediction for the Ising model, $`\beta _{Ising}=1/(2d)`$, only for the special dimension $`d=2`$. For $`d>2`$ we observe an increase of $`\beta _c`$ with $`d`$, to be compared to the opposite trend shown by the Ising model. Such behaviour may be interpreted in terms of the low dimensionality of the ordered phase. Due to frustration the ordering may only occur on a low dimensional scale: for instance in three dimensions each layer has an independent internal ordering. Thus we expect a larger $`\beta _c`$ for $`d>2`$ since the increasing of $`d`$ only introduces larger fluctuations, with each molecule having $`(d2)`$ allowed out-of-plane orientations. For $`d=3`$ the low temperature phase can be regarded as a quenched disordered superposition of layers which are internally ordered along different in-plane directions. As a consequence of frustration the system shows a two-dimensional character below the critical point while behaving as truly three-dimensional in the high temperature domain. In MF the neglecting of some fluctuations usually leads to a critical temperature which overestimate the exact value (i.e. the critical inverse temperature $`\beta _c`$ is underestimated). For $`d=2`$ the MF prediction is $`\beta _c=0.25`$ to be compared with the exact value $`\beta _c=0.4407`$. For $`d=3`$ the MF prediction $`\beta _c=d/8=0.375`$ should provide a lower bound to the unknown exact value.
## IV Variational Migdal-Kadanoff Approximation
According to the Migdal-Kadanoff method, a link displacement may be introduced by considering that the configurational average of the Lagrangian density $``$ in eq.(5) must be translationally invariant
$$(𝐫,\alpha )=(𝐫^{},\alpha )$$
(17)
then defining
$$\mathrm{\Gamma }_\alpha (𝐫,𝐫^{})=(𝐫,\alpha )(𝐫^{},\alpha )$$
(18)
we can state that $`\mathrm{\Gamma }_\alpha (𝐫,𝐫^{})=0`$ and the same holds for any sum $`\mathrm{\Gamma }`$ over an arbitrary set of such terms
$$\mathrm{\Gamma }=\mathrm{\Gamma }_\alpha (𝐫,𝐫^{}).$$
(19)
Replacing the action $`S_h`$ by the sum $`S_h+\mathrm{\Gamma }`$, and assuming that the condition (3) is verified (so that we can drop the $`h`$ in $`S_h`$ and $`Z_h`$ which are invariant), the modified partition function $`Z_\mathrm{\Gamma }`$ can be approximated by cumulant expansion as
$$Z_\mathrm{\Gamma }=\underset{\{\widehat{w}\}}{}e^{S+\mathrm{\Gamma }}=Ze^\mathrm{\Gamma }Z\left[e^\mathrm{\Gamma }e^{\frac{1}{2}\left(\mathrm{\Gamma }^2\mathrm{\Gamma }^2\right)}\right]$$
(20)
then, since $`\mathrm{\Gamma }=0`$,
$$Z_\mathrm{\Gamma }Ze^{\frac{1}{2}\mathrm{\Gamma }^2}$$
(21)
For instance, the sum in equation (19) could run over all $`\alpha 1`$, and for appropriate values of the vectors $`𝐫,𝐫^{}`$, in order to yield a displacement of links which are orthogonal to $`\widehat{x}_1`$. To second order in $`\mathrm{\Gamma }`$, the error introduced by link displacement is controlled by the exponential factor in equation (21).
Link displacement breaks the internal symmetry of the model, so that $`Z_\mathrm{\Gamma }`$ is no longer invariant for any field change subject to the condition (3). Then we may improve the approximation by using the extra freedom on the choice of $`𝐡`$ for minimizing the difference between the approximate partition function $`Z_\mathrm{\Gamma }`$ and the exact $`Z`$.
If $`𝐡`$ satisfies the condition (3) then $`a𝐡`$ satisfies such condition as well for any choice of the scalar parameter $`a`$. Then a special class of invariance transformations can be described by a change of the strength parameter $`h`$, assuming the field $`𝐡`$ as proportional to $`h`$. The following discussion could be easily generalized to other classes of transformations described by more than one parameter. Since $`\mathrm{\Gamma }`$ is linear in the field $`𝐡`$, then in general
$$\mathrm{\Gamma }^2=\left[A+hB\right]^2$$
(22)
where $`A`$ and $`B`$ depend on the configuration of the system. For the average we have
$$\mathrm{\Gamma }^2=A^2+2hAB+h^2B^2$$
(23)
This last equation, inserted in eq.(21) leads to the following considerations: i) the coefficient $`B^2`$ is positive defined, then the average $`\mathrm{\Gamma }^2`$ always has a minimum for an appropriate value of $`h=h_0`$; ii) in general $`AB0`$ then $`h_00`$, and a direct use of the Migdal-Kadanoff method on the original model (with no field considered) would yield a larger error; iii) to the considered order of approximation $`Z_\mathrm{\Gamma }`$ is stationary at $`h=h_0`$, and is symmetric around that point, then all the physical properties described by such partition function must result symmetric with respect to $`h_0`$. Moreover, at the same order of approximation, any physical observable $`f`$ will acquire an unphysical dependence on $`h`$, and the symmetry around $`h_0`$ requires that $`\frac{df}{dh}=0`$ for $`h=h_0`$. Then we expect that all such observables should be stationary at $`h=h_0`$, and their best estimate should coincide with the extreme value.
As a consequence of the above statements, the Migdal-Kadanoff method can be improved by taking advantage of the global symmetry of the model. By use of the approximate partition function $`Z_\mathrm{\Gamma }`$ the critical temperature acquires a non-physical field dependence, but the best estimate of $`\beta _c`$ is its stationary value corresponding to $`h=h_0`$. The method can be seen as a variational method with the best approximation achieved by the minimum in the inverse temperature.
Such stationary condition resembles the principle of “minimum sensitivity” introduced by Stevenson for determining the best renormalization parameters whenever the physical amplitudes depend on them (and they should not). In our context, since the critical temperature should not depend on the choice of the field strength $`h`$, the best value for such field is the one which makes the critical temperature less sensitive i.e. the stationary point. However, according to equation (21) and (23), here we have got a formal proof of the stationary condition up to second order of the cumulant expansion.
The method may be used by performing a displacement of links that are orthogonal to $`\widehat{x}_1`$, and then a one-dimensional decimation along the $`\alpha =1`$ axis. According to such program let us define the alternative $`d\times d`$ matrix $`t_\alpha (\beta ,h)`$
$$e^{(𝐫,\alpha )}=\widehat{w}_𝐫^{}t_\alpha (\beta ,h)\widehat{w}_{𝐫+\widehat{x}_\alpha }$$
(24)
The partition function follows
$$Z_h=\underset{\{\widehat{w}\}}{}\underset{𝐫,\alpha }{}\left[\widehat{w}_𝐫^{}t_\alpha (\beta ,h)\widehat{w}_{𝐫+\widehat{x}_\alpha }\right]$$
(25)
After link displacement and decimation along the $`\alpha =1`$ axis, the modified partition function reads
$$Z_\mathrm{\Gamma }=\underset{\{\widehat{w}\}}{}\underset{𝐫,\alpha }{}\left[\widehat{w}_𝐫^{}\stackrel{~}{t}_\alpha (\beta ,h)\widehat{w}_{𝐫+\widehat{x}_\alpha }\right]$$
(26)
where the sum and the product run over the configurations and the sites of the new decimated lattice, and
$$\stackrel{~}{t}_1(\beta ,h)=\left[t_1(\beta ,h)\right]^\lambda $$
(27)
$$\stackrel{~}{t}_\alpha (\beta ,h)=t_\alpha (\lambda \beta ,h)\mathrm{for}\alpha 1$$
(28)
with $`\lambda `$ being the scale factor between the new and the old lattice. A renormalized inverse temperature $`\stackrel{~}{\beta }_\alpha `$ may be defined according to
$$\stackrel{~}{t}_1(\beta ,h)=t_1(\stackrel{~}{\beta }_1,h)$$
(29)
$$\stackrel{~}{\beta }_\alpha =\lambda \beta \mathrm{for}\alpha 1$$
(30)
Eventually, the same scaling operation should be performed consecutively for all the directions in order to obtain an hyper-cubic lattice again. For any finite scaling parameter $`\lambda >1`$ the renormalized inverse temperature is anisotropic, but an isotropic fixed-point can be recovered in the limit $`\lambda 1`$. The equations (29),(30) define the flow of the renormalized inverse temperature, which changes for any different value of the field strength $`h`$. Equation (29) has a more explicit aspect in the representation of the common eigenvectors of the matrices $`t_1`$ and $`\stackrel{~}{t}_1=[t_1]^\lambda `$. The rank of such matrices is 2 for any space dimension $`d`$, as can be expected from the definition of the model. Then both the matrices can be represented in terms of the two non-vanishing eigenvalues $`\eta _1`$, $`\eta _2`$, which are functions of $`\beta `$ and $`h`$. Assuming that $`\eta _20`$, and defining
$$f(\beta ,h)=\frac{\eta _1}{\eta _2}$$
(31)
apart from a regular multiplicative factor for the partition function, the scaling equation (29) reads
$$\left[f(\beta ,h)\right]^\lambda =f(\stackrel{~}{\beta }_1,h)$$
(32)
For any fixed $`h`$, the fixed points follow through the standard Migdal-Kadanoff equations
$$\left[f(\lambda ^{\alpha 1}\beta _\alpha ,h)\right]^\lambda =f(\lambda ^{\alpha d}\beta _\alpha ,h).$$
(33)
When $`\lambda `$ is analytically continued up to $`1`$ such equations give the same isotropic fixed point $`\beta _c`$. In fact, the expansion of equations (33) around $`\lambda =1`$ implies (up to first order in $`\lambda 1`$)
$$\mathrm{ln}f(\beta _c,h)=(d1)\beta _c\left[\frac{1}{f}\frac{df}{d\beta }\right]_{\beta _c}$$
(34)
which is an implicit equation for $`\beta _c`$. Such equations yield their best estimate of $`\beta _c`$ when the strength of the field $`h`$ is set to the stationary value $`h_0`$.
It is instructive to evaluate the stationary point $`h_0`$ for the case $`d=2`$ which is equivalent to the two-dimensional Ising model for the choice $`h=h_I=1/4`$, as shown in section II. The $`h`$ invariance of the exact partition function guarantees the equivalence of the two models for any choice of $`hh_I`$. However, the mere application of the Migdal-Kadanoff equations (33) to the simple $`h=0`$ molecular model fails to predict even the existence of the fixed point. On the other hand, for $`h=h_I`$, the very same recurrence equations (33) are known to predict the exact fixed point in the limit $`\lambda 1`$. That can also be checked by inserting in equation (34) the exact expression for the fixed point of the two-dimensional Ising model. Such contradictory results are not surprising since, as already discussed, link displacement breaks the $`h`$ invariance of the model, and the approximate solution is thus dependent on the choice of $`h`$. We would like to test the variational method on this exactly solvable model: we look for the stationary point of the function $`f(\beta ,h)`$. The matrices $`t_\alpha `$ follow from the equations (7)
$`t_1`$ $`=`$ $`\left(\begin{array}{cc}x^2b& 1\\ 1& x^2\end{array}\right)`$ (35)
$`t_2`$ $`=`$ $`\left(\begin{array}{cc}x^2& 1\\ 1& xb\end{array}\right)`$ (36)
where $`b=exp(4\beta )`$, and $`x=exp(4\beta h)`$. Then for the eigenvalues we obtain
$$f(\beta ,h)=\frac{\eta _1}{\eta _2}=\frac{(bx^4+1)\sqrt{(bx^41)^2+4x^4}}{(bx^4+1)+\sqrt{(bx^41)^2+4x^4}}$$
(37)
It can be easily shown that if the derivative of $`f`$ is zero at a given $`h`$ independent of $`\beta `$, then the solution $`\beta _c`$ of Eq.(34) is stationary at that $`h`$ value. Differentiating with respect to $`x`$, we find that the derivative of $`f`$ vanishes for $`x^4=1/b`$, which yields $`h=1/4=h_I`$ for any $`\beta `$. As expected, this is the required value in order to recover the Ising model. Thus the Migdal-Kadanoff approximation gives an improving estimate of the critical point as we move from the molecular towards the Ising representation (where the approximation yields the exact fixed point). We stress that all such representations are equivalent due to the $`h`$ invariance of the action.
For $`d>2`$ no equivalence to standard studied models has been found, and the behaviour seems to be dictated by the strong frustration which does not allow an higher coordination number than two, even for higher dimensions. We will focus on the $`d=3`$ model in order to compare the results with the Monte Carlo findings of the next section. First of all the fields $`𝐡(\alpha )`$ must be defined. An isotropic choice would be
$`𝐡(1)`$ $`=`$ $`h(\widehat{x}_1{\displaystyle \frac{1}{2}}\widehat{x}_2{\displaystyle \frac{1}{2}}\widehat{x}_3)`$ (38)
$`𝐡(2)`$ $`=`$ $`h(\widehat{x}_2{\displaystyle \frac{1}{2}}\widehat{x}_3{\displaystyle \frac{1}{2}}\widehat{x}_1)`$ (39)
$`𝐡(3)`$ $`=`$ $`h(\widehat{x}_3{\displaystyle \frac{1}{2}}\widehat{x}_1{\displaystyle \frac{1}{2}}\widehat{x}_2)`$ (40)
The matrix $`t_1`$ follows
$$t_1=\left(\begin{array}{ccc}e^{4\beta +8\beta h}& e^{2\beta h}& e^{2\beta h}\\ e^{2\beta h}& e^{4\beta h}& e^{4\beta h}\\ e^{2\beta h}& e^{4\beta h}& e^{4\beta h}\end{array}\right)=\left(\begin{array}{ccc}bx^2& \sqrt{x}& \sqrt{x}\\ \sqrt{x}& 1/x& 1/x\\ \sqrt{x}& 1/x& 1/x\end{array}\right)$$
(41)
Then from the eigenvalues we obtain
$$f(\beta ,h)=\frac{\eta _1}{\eta _2}=\frac{(bx^3+2)\sqrt{(bx^32)^2+8x^3}}{(bx^3+2)+\sqrt{(bx^32)^2+8x^3}}$$
(42)
Differentiating with respect to $`x`$ we find that the derivative of $`f`$ only vanishes for $`x^3=2/b`$. This is equivalent to say
$$h=h_m=\frac{1}{3}+\frac{1}{12\beta }\mathrm{ln}2$$
(43)
which depends on $`\beta `$. For such field strength $`h_m(\beta )`$ the ratio between the eigenvalues reduces to
$$f(\beta ,h_m)=\mathrm{tanh}(\beta )$$
(44)
which is exactly the same expression holding for the Ising model. However we must point out that in such case $`h_m`$ is not the stationary point $`h_0`$. Since $`h_m`$ depends on $`\beta `$, the vanishing of the derivative of $`f`$ does not imply that the solution of eq.(34) is stationary. In fact, for $`d=3`$, the choice $`h=h_m`$ yields the known poor result $`\beta _c=0.1398`$ by insertion of Eq.(44) in Eq.(34). On the other hand, by insertion of the general expression for $`f`$ Eq.(42), the scaling equation (34) can be numerically solved for $`\beta _c`$ as a function of $`h`$. At the stationary point $`\beta _c`$ has a minimum, and thus the variational method yields an even worse prediction ($`\beta _c0.12`$ at the stationary point). These shortcomings show that the isotropic $`d=3`$ variational method does not suite the molecular model. Actually both MF and Monte Carlo methods predict a larger $`\beta _c`$ and, as pointed out at the end of the previous section, the exact $`\beta _c`$ should be larger than the MF prediction $`\beta _{MF}=0.375`$.
We could have guessed such disagreement since we are using an isotropic version of the variational Migdal-Kadanoff method for a system which is not isotropic in its ordered phase. At the transition point the system choices a direction, as is usual for any symmetry breaking mechanism. However, at variance with usual models, in the ordered phase the correlation length cannot be isotropic: order occurs inside all layers which are orthogonal to the chosen direction, while there is a negligible correlation along such direction. It would be more sensible to describe the ordering which takes place inside a single layer, thus neglecting any correlation among different layers. Inside each layer the correlation length is isotropic, and the $`d=2`$ variational Migdal-Kadanoff method should give a better description of the transition. The same argument should hold for the generic $`d`$-dimensional molecular model. Moreover, despite the cost of this further approximation, the Migdal-Kadanoff method is known to work better for the lower dimensions, and a $`d=2`$ variational method could provide a tool for describing the generic $`d`$-dimensional molecular model even for $`d>3`$.
A $`d=2`$ version of the variational method requires a different choice for the fields $`𝐡(\alpha )`$ which do not need to be isotropic any more. Let us take the same field we used in section II, namely $`𝐡(1)=h(\widehat{x}_1\widehat{x}_2)`$, $`𝐡(2)=𝐡(1)`$ and $`𝐡(3)=0`$. The matrix $`t_1`$ follows
$$t_1=\left(\begin{array}{ccc}e^{4\beta +8\beta h}& 1& e^{4\beta h}\\ 1& e^{8\beta h}& e^{4\beta h}\\ e^{4\beta h}& e^{4\beta h}& 1\end{array}\right)=\left(\begin{array}{ccc}bx^2& 1& x\\ 1& 1/x^2& 1/x\\ x& 1/x& 1\end{array}\right)$$
(45)
Notice that this is a $`3\times 3`$ matrix since we are using the $`d=2`$ method but we are still dealing with a $`d=3`$ molecular model. The two matrices $`t_1`$ and $`t_2`$ share the same eigenvalues. Their ratio is
$$f(\beta ,h)=\frac{(bx^2+1+1/x^2)\sqrt{(bx^211/x^2)^2+4(1+x^2)}}{(bx^2+1+1/x^2)+\sqrt{(bx^211/x^2)^2+4(1+x^2)}}$$
(46)
Inserting this result in the scaling equation (34) evaluated at $`d=2`$ yields an implicit equation for $`\beta _c`$ versus $`h`$. The numerical solutions are reported in Fig.3. They share most of the features of the $`d=2`$ molecular model: (i) There are several solutions but there is no repulsive fixed point for $`h=0`$; (ii) The physical solution starts at a negative $`h`$ which in this case is $`h0.226`$; (iii) The physical solution has just one stationary point $`h_0`$ where $`\beta _c`$ reaches its minimum value. However in this case the stationary point is at $`h_0=0.2349`$ where $`\beta _c=0.6122`$. This best estimate of the critical point is not too far from the finite size scaling prediction of the next section $`\beta _c=0.53`$. The result is encouraging, and gives us more confidence in our understanding of the physics described by the molecular model. Strictly speaking, this $`d=2`$ variational method describes the transition occurring in a single layer of molecules. However, at variance with the $`d=2`$ molecular model, each molecule is now allowed to be oriented along three different axes (two in-plane and one out-of-plane orientations). Thus this reasonable prediction for $`\beta _c`$ could be regarded as an indirect proof that the correlation between two different layers is negligible, and that in the ordered phase the system behaves as a truly two-dimensional one.
## V Monte Carlo sampling
In order to check the prediction achieved by different approximate methods it would be desirable to have an accurate numerical estimate of the critical temperature. That can be easily obtained by finite size scaling. Moreover, according to the scaling hypothesis, the critical exponent $`\nu `$ can be extracted by the numerical data with a good accuracy.
Cubic samples $`N\times N\times N`$ with $`N=10,15,20,25,30`$ have been considered. All the averages have been evaluated by Monte Carlo sampling with no special boundary conditions.
In this model any ordering is characterized by the presence of some degree of correlation along one-dimensional chains of molecules. For $`d=3`$ there are $`3\times N\times N`$ different chains in each sample. Each chain may be labelled by its direction $`\alpha =1,2,3`$ and by a couple of integer coordinates $`I_1,I_2`$ running over a lattice layer orthogonal to the axis $`\widehat{x}_\alpha `$. For any chain we define an order parameter
$$m(\alpha ,I_1,I_2)=\frac{1}{N}\underset{J_\alpha =1}{\overset{N}{}}\left[\widehat{w}(J_\alpha ,I_1,I_2)\widehat{x}_\alpha \right]$$
(47)
where $`\widehat{w}(J_\alpha ,I_1,I_2)`$ is the versor $`\widehat{w}_𝐫`$ at the chain site $`𝐫`$ whose integer coordinates are determined by $`J_\alpha `$ along the chain and by the couple $`I_1,I_2`$ in the orthogonal directions. If there is no correlation at all ($`\beta 0`$) then $`m(\alpha ,I_1,I_2)1/3`$ for any chain in the sample. By averaging over all the chains of each sample and over all the configurations, we obtain $`m=1/3`$. For large $`N`$, according to the central limit theorem, in this statistical ensemble the variable $`m`$ follows a gaussian distribution centered at its average value. In the opposite limit ($`\beta \mathrm{}`$) a third of the chains in each sample have a large $`m1`$, while $`m0`$ for two thirds of them. Since any intermediate value of $`m`$ is unlikely, the statistical distribution of $`m`$ can be regarded as the superposition of two different peaked distributions centered at $`m=0`$ and $`m=1`$. If $`N`$ is large enough, and for a large number of configurations, such distributions are very peaked and their width is very small. Actually, just below the critical point the gaussian distribution already splits in a double-peak distribution. We can monitor the transition by use of the new variable $`\gamma `$
$$\gamma =\frac{mm}{\sqrt{m^2m^2}}$$
(48)
By its definition the configurational average of $`\gamma `$ is vanishing $`\gamma =0`$ and the second moment $`\gamma ^2=1`$. The variable $`\gamma `$ only differs from $`m`$ for a shift and a rescaling, thus the statistical distribution for $`\gamma `$ follows the same trend already discussed for $`m`$. However the fourth moment $`s=\gamma ^4`$ is now strongly dependent on the number of peaks characterizing the statistical distribution. For a single gaussian $`s=3`$ exactly. Below the critical temperature the distribution splits. In the thermodynamic limit $`N\mathrm{}`$ the width of each peak vanishes, while the two peaks separate by a finite quantity. For instance assume that just below the critical point a third of the chains yield $`m1/3+ϵ`$ where $`ϵ`$ is a very small increase in the chain correlation which breaks the symmetry of the sample. The other two thirds of chains must yield $`m1/3ϵ/2`$ since by its definition $`m=1/3`$ exactly. Neglecting the width of the peaks we may approximate the statistical distribution for $`m`$ as the superposition of two delta-functions with weight factors:
$$P(m)=\frac{2}{3}\delta (m\frac{1}{3}+\frac{ϵ}{2})+\frac{1}{3}\delta (m\frac{1}{3}ϵ)$$
(49)
By use of such approximate statistical distribution the calculation of the fourth moment $`s=\gamma ^4`$ is straightforward and gives $`s=1.5`$ for any $`ϵ`$, no matter how small. This is one half of the single gaussian value. Thus in the thermodynamic limit we expect that the fourth moment $`s`$ should behave like a step function with constant values $`s=3`$ and $`s=1.5`$ respectively above and below the critical temperature, and a sharp jump at the critical point. For finite-size samples the fourth moment is expected to be continuous across the transition, but according to the scaling hypothesis the critical value should not depend on the sample size if we assume a one-parameter scaling law across the critical point:
$$s=s\left(L/\xi (\beta )\right)$$
(50)
where $`L`$ is here the sample length, and $`\xi `$ is the correlation length which is a function of temperature. According to such scaling law $`s=s(0)`$ at the critical point for any $`L`$.
We have checked this prediction by standard Monte Carlo sampling. For any fixed sample size, we have taken a completely random initial configuration, and thermalized it at a very high temperature ($`\beta 0.02`$) by $`510^4`$ complete sweeps. The temperature is then decreased by steps of $`\mathrm{\Delta }\beta =0.02`$. At each step a good thermalization is achieved by $`810^3`$ complete sweeps, and then the averages are evaluated over the successive $`210^3`$ sweeps. Once a sufficiently low temperature is reached ($`\beta 1`$), the process is reversed and the temperature increased up to the initial value. We have checked that the hysteresis is small in all the considered range of temperature. Moreover the small differences observed going up and down give a measure of the errors on the configurational averages which have been approximated by the mean values. The fourth moment $`s`$ is reported in Fig.4 for $`N=15,20,25,30`$. All the curves cross at the same point $`\beta _c=0.53\pm 0.01`$ as predicted by the one-parameter scaling hypothesis. Moreover for very large or very small temperatures the correlation length becomes very small and the fourth moment $`s`$ should approach its thermodynamic-limit value $`ss(\mathrm{})`$ which is expected to be $`s(\mathrm{})=3`$ at high temperature and $`s(\mathrm{})=1.5`$ at low temperature. As shown in Fig.4 the measured $`s`$ approaches such limits far away from the critical point.
According to the usual definition of critical exponent
$$\xi \frac{1}{(\beta \beta _c)^\nu }$$
(51)
the scaling equation (50) allows for an accurate estimate of its value: linearizing $`s`$ around the critical point yields
$$\mathrm{ln}L=\nu \mathrm{ln}|s^{}(\beta _c)|+\mathrm{const}.$$
(52)
where $`s^{}(\beta _c)`$ is the derivative of $`s`$ as a function of $`\beta `$. In Fig.5 a best fit by least squares method is reported yielding $`\nu =0.44\pm 0.02`$. Here the error is the statistical one coming out from the linear fit.
Of course this Monte Carlo calculation is far from being the best numerical simulation which can be achieved by modern computing machines. Our sample sizes are relatively small and a slight shift of the critical point cannot be ruled out. However the estimates for the critical temperature and exponent are accurate enough for a comparison with experimental findings and for a check of the analytical results of the previous sections, and that is just what we needed a the moment. More refined calculations are called for in order to establish more accurate predictions.
## VI Discussion
Here we summarize and discuss the main findings of the previous sections. According to mean-field and finite-size scaling the three-dimensional molecular model has a second order continuous transition from an isotropic disordered high-temperature phase to an anisotropic two-dimensional ordered low-temperature phase. The $`d=3`$ realization of the model is the one which more closely describes real molecular systems. For this reason the $`d=3`$ model has been studied by the variational Migdal-Kadanoff method and by numerical Monte Carlo simulation. The transition point is characterized by a diverging correlation length according to the one-parameter scaling hypothesis which seems to be fulfilled as shown by the data of the previous section. On the other hand the $`d=2`$ model is special by itself for its equivalence to the two-dimensional Ising model, and for the existence of exact analytical results. Thus the $`d=3`$ model can be seen as a non-trivial extension to higher dimension of the two-dimensional Ising model. Here “non-trivial” means that the $`d=3`$ molecular model does not belong to the universality classes of the standard $`d=3`$ extensions of the Ising model (three-dimensional Ising and Potts models). The difference is evident from a comparison of the ground state $`T=0`$ configurations: highly degenerate and anisotropic in the molecular model (with a two-dimensional character even for higher dimensions); with a small degeneration and fully isotropic in the Potts models (including the Ising one as a special case). By considering the two-dimensional character of the low-temperature phase, the molecular model could be thought to belong to the universality class of the simple two-dimensional Ising or three-states Potts models. However in the high temperature unbroken-symmetry phase the molecular model is fully isotropic and has a three-dimensional character.
A formal proof of such statements comes from a comparison of the critical exponents. For the three-dimensional molecular model the finite size scaling calculation of the previous section yields $`\nu =0.44`$ to be compared to the two-dimensional two-state (Ising) and three-state Potts models whose exponents are $`\nu =1`$ and $`\nu =0.83`$ respectively, to the three-dimensional Ising model whose exponent is $`\nu =0.64`$, and to the three-state three-dimensional Potts model which is known to undergo a first-order transition.
The molecular model belongs to a new universality class which is characterized by a sort of dimensional transmutation. In fact order takes place in chains which are arranged in layers, and the disorder-order transition requires a decrease of the effective dimensionality of the system. In the ordered phase the molecules are correlated inside layers, but there is no correlation between molecules which belong to different layers. This understanding of the ordered phase is in agreement with our finding that the two-dimensional Migdal-Kadanoff variational method for a single layer yields a better prediction for the critical point than the three-dimensional method applied to the whole lattice. On the other hand the very same two-dimensional variational method provides a convenient analytical tool for describing the generic $`d`$-dimensional molecular model by a straightforward generalization.
From such arguments the critical point has been given an upper bound by the variational method which yields $`\beta _c=0.61`$, while a lower bound is usually provided by mean-field that in this $`d=3`$ case gives $`\beta _c=0.375`$. The numerical estimate of the previous section $`\beta _c=0.53`$ fits nicely inside such bounds.
Having discussed some formal aspects of the molecular model and few approximate methods which throw some light on its phase transition, we would like to make contact with the phenomenology. Our main finding - that the order-disorder transition described by the model belongs to a new universality class - deserves some experimental test. Transitions of this kind have been observed in several systems, as discussed in the introduction. Since the critical properties should not depend on the microscopic details of the system we expect that the simple molecular model could predict the correct critical exponent of real orientational transitions occurring in complex real molecular systems expecially under pressure. New experiments are called for in order to test such ideas and explore this broad universality class.
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# New Insights Into The Narrow-Line Seyfert 1 Phenomenon
## 1 Introduction
The organizers of this workshop have asked me to talk about the X-ray/UV absorbers in AGN and discuss what they can tell us about the Narrow Line Seyfert 1 phenomenon. However, I am going to go well beyond my expertise and speculate. I think that this meeting is an ideal platform for discussing new ideas, so let me start with my conclusions first and get them out of the way. Conclusions:
1. The covering factor of warm absorbers is likely to be high in NLS1s.
2. Metallicities might be super solar.
3. NLS1s may be AGN in the making and as such, may be low redshift, low luminosity analogues of the high redshift quasars.
4. NLS1s reside in rejuvenated/ gas rich galaxies.
Having said that, let’s start from the beginning.
### 1.1 The X/UV Absorbers
The UV and X-ray intrinsic absorption systems observed in active galactic nuclei (AGN) were never thought to be related to one another because of their apparently different physical properties. The inferred total column density of UV absorption lines was believed to be about $`10^{20}`$ cm<sup>-2</sup>, and the ionization parameter U was thought to be such that CIV would be the dominant ionization phase of carbon. In contrast, the cold X-ray absorbers are neutral with column densities $`>10^{21}`$ cm<sup>-2</sup>. Models of warm X-ray absorbers also required larger column densities and ionization parameters than the UV absorbers. As was shown later, the ionization structure of the absorbers was what was lacking from our understanding.
The above situation changed with the quasi-simultaneous observations of 3C351 with the ROSAT PSPC and the HST FOS. The UV spectrum of 3C351 showed OVI absorption line doublets and the X-ray spectrum showed edges due to OVII/OVIII. These observations lead Mathur et al. (1994) to conclude that the UV and X-ray absorbers are in fact one and the same. The combined UV and X-ray analysis proved to be a powerful technique to probe the physical properties of the nuclear regions of AGN: high column density, highly ionized, outflowing material situated outside the broad emission line region. In addition, the density of the absorber could be constrained in variable systems. This was a previously unknown component with important consequences: an outflow that carries a significant amount of kinetic energy, and with a mass outflow rate comparable to the accretion rate needed to power the AGN engine.
Later work supported the unified X/UV absorber scenario (Mathur 1994, Mathur et al. 1995, 1997, 1998, 1999, Shields & Hamann 1997a, Crenshaw 1997). I will defer to Mathur (1997) for a discussion of controversies regarding the X/UV models. I conclude that all the observed data support this picture: the UV and X-ray absorbers are physically related to each other. At the very least, the X-ray absorbers make a substantial contribution to the absorption seen in the UV.
## 2 X-ray & UV Absorption in NLS1s
Table 1 lists some NLS1 galaxies showing absorption in X-rays and/or UV. I have compiled this list from the literature and it is by no means complete. There are a few more ROSAT selected NLS1s with X-ray warm absorbers for which UV spectra are not available.
From Table 1 we see that for some NLS1s with X-ray warm absorbers either (1) UV data do not exist, or (2) UV absorption is not expected because the X-ray absorber is very highly ionized. I cannot comment on the X/UV connection for these objects. Excluding these objects, we have warm absorbers with either (1) UV absorption, or (2) blue asymmetric emission lines. The first case is what is expected from the X/UV absorber models, while the second (blue asymmetric emission lines) is a new observation and is intriguing. I would like to propose that the blue asymmetry is caused by emission from the outflowing/inflowing warm absorber. (See also B. Peterson’s article on the broad HeII emission line in NGC4051). A preliminary investigation shows consistency with the predictions from photoionization models if the covering factor is large. This leads to my first conclusion: Warm absorbers in NLS1s have large covering factors.
Now consider the warm absorbers with UV absorbers. The existence of UV absorbers is in itself favored by the unified X/UV models. There are three such objects listed in Table 1: (1) IRAS 13349+2438: the ROSAT spectrum of this object is discussed in Brandt, Mathur, Reynolds & Elvis (1997). Using photoionization models and the parameters of warm absorbers, we can predict the column density of UV absorption lines. HST FOS observed the CIV absorption line in this object (Figure 1). Being a single orbit snapshot observation, the signal-to-noise ratio is poor. Also, with just one line, the result cannot be conclusive. Nevertheless, the observed CIV column density is as predicted for a reasonable value of the b parameter, b$`100`$ km/s. (2) I Zw 1: This object has a weak warm absorber, and as expected, a weak associated absorption line system of Ly$`\alpha `$, NV and CIV was observed with HST (Laor 1998a). (3) Mrk 1298 (or PG1126$``$041): CIV and NV absorption lines are clearly present in the low dispersion, IUE spectrum of this object (Figure 2). The S/N, however, is too low to develop detailed models. To conclude, the NLS1s also support the unified X/UV absorber models.
This now brings me to the last object in the Table, PG1404+266.
### 2.1 PG 1404+266
ASCA, ROSAT and HST observations of PG1404+266 are discussed in Ulrich et al. 1999. The X-ray spectrum of the object shows the signature of a $``$1 keV absorption feature which may be interpreted as a NeVII–NeX warm absorber (note that this would be different from the OVII/OVIII warm absorbers generally observed with ROSAT). PG 1404+226 also has an unusually strong Fe K$`\alpha `$ emission line. The HST spectrum shows the presence of the CIII$`\lambda `$1175.7 emission line and associated absorption lines of Ly$`\alpha `$, NV and CIV. Ulrich et al. found that, while the strength of the Ly$`\alpha `$ and CIV lines was generally consistent with that expected from the warm absorber, the strength of NV was not: the NV absorption line was stronger than expected.
This is reminiscent of the metallicity determinations in quasars (see review by Hamann & Ferland 1999). As shown by Hamann & Ferland, nitrogen serves as an metallicity indicator. Moreover, nitrogen is preferentially enhanced when metallicities are high, N$`/`$H $``$ Z<sup>2</sup>. The observations of PG1404+226 thus suggest that nitrogen is enhanced and so the metallicities are super solar in this object and likely to be in NLS1s in general. We are investigating whether this suggestion is correct by performing a detailed metallicity analysis (Mathur & Komossa 2000).
## 3 Are NLS1s Active Galaxies in the making?
In addition to the case of PG 1404+226 discussed above, there are other lines of evidence suggestive of high metallicities in NLS1s. Wills et al. (1999) found that the strength of the NV $`\lambda 1240`$ emission line was systematically larger while the strength of CIV $`\lambda 1549`$ was systematically smaller in AGN with narrow emission lines (see also B. Wills, these proceedings). The strength of the fluorescent Fe-K alpha line in some NLS1s is also indicative of a super-solar abundance (A. Fabian, these proceedings). Thus the observations of emission as well as absorption lines in NLS1s imply super-solar gas phase metallicities.
Such metal enrichment is possible when the initial mass function of star formation is flat, favorable for high mass star formation, and the evolution is fast. Such a star formation scenario is likely to be present in deep potential wells like galactic nuclei and protogalactic clumps (HF99). Moreover, high metallicities are achieved while consuming less gas (HF99). The NLS1s may then represent that early phase in galactic evolution when rapid star formation is taking place in the nucleus.
Note also that NLS1s have relatively smaller BH masses. As per the well known correlation of Magorrian et al. (1998), smaller mass BHs reside in galaxies with smaller spheroids. Since NLS1s have relatively smaller mass BHs compared to normal Seyferts, the spheroids of their host galaxies might be smaller (see also Laor 1998b). Indeed, in the compilation of Wandel (1999), the NLS1 galaxy NGC4051 has the smallest black hole to bulge mass ratio. An accreting BH would also grow in mass with time \[the Salpeter time scale of growth is determined by t$`{}_{s}{}^{}=3\times 10^7(L_{Edd}/L_B)\eta _{0.1}`$ yr. where $`\eta _{0.1}`$ is the radiative efficiency in units of 0.1 (see Fabian 1999)\]. Since NLS1s accrete at close to the Eddington limit, their BHs would grow faster. So, smaller BHs in NLS1s are likely to be younger as well.
These arguments support my proposal that NLS1s might be Active Galaxies in early phases of their evolution.
## 4 Similarity between NLS1s and the high redshift quasars.
Hamann & Ferland (1993) found high metallicities in high redshift quasars (Z$`\genfrac{}{}{0pt}{}{_>}{^{}}`$Z at z$`\genfrac{}{}{0pt}{}{_>}{^{}}`$4). Similarly we find that NLS1s may also have large metallicities.
Similarities between the observed properties of low ionization Broad Absorption Line Quasars (BALQSOs) and NLS1s have been reported in the literature (e.g. Lawrence et al. 1997, Leighly et al. 1997). Both these classes show strong Fe II$`\lambda 4570`$ and AlIII$`\lambda 1857`$ and weak CIV$`\lambda 1549`$ and \[OIII\]$`5007`$ emission lines. Their continua are red in the optical and strong in the IR. Evidence of relativistic outflow is also reported in three NLS1s (Leighly et al. 1997). If these two classes are indeed related (see also Brandt, these proceedings), then NLS1s, at least those with some evidence of outflow, might be low redshift, low luminosity cousins of BALQSOs. BALQSOs are tentatively identified with a phase in quasar evolution when the matter around the nuclear BH is being blown away, and a quasar emerges (see, e.g. Fabian 1999). NLS1s may then represent a similar early evolutionary phase at low redshift.
Optical spectra of a sample of z$`\genfrac{}{}{0pt}{}{_>}{^{}}`$4 quasars revealed that their emission lines are typically narrower than the low redshift quasars (FWHM $`\genfrac{}{}{0pt}{}{_<}{^{}}`$2000 km/s, Shields & Hamann 1997b). The normal explanation of this observation is that these are type 2 quasars, where the broad emission lines are obscured from our line of sight. Alternatively, these high z quasars might be true “narrow” broad line objects (see Mathur 2000 for a discussion of selection effects).
There is another interesting connection with high redshift. As discussed above, NLS1s have strong Fe II emission lines. Quasars Q0014+813 and Q0636+680 at redshifts z=3.398 and z=3.195 respectively, were observed to have very strong Fe II emission (Elston, Thompson, & Hill 1994). Are they also highly accreting objects at an early evolutionary phase? Note also the narrow UV emission lines (FWHM $`\genfrac{}{}{0pt}{}{_<}{^{}}`$2150 km s<sup>-1</sup>) in the ultra strong UV Fe II emitter Q2226-3905 (Graham, Clowes, & Campusano, 1996).
All these similarities point towards NLS1s being low redshift, low luminosity analogues of high redshift quasars.
## 5 Do NLS1s reside in rejuvenated galaxies?
I have argued that NLS1s may represent an early phase in AGN evolution. Whether they reside in young galaxies is a separate question and a step further. That young galaxies are gas rich is helpful; they would have the large reservoir of gas necessary to sustain the close to Eddington-rate accretion in NLS1s. But do we have any evidence that they indeed reside in young galaxies? There is no published systematic study of the properties of the host galaxies of NLS1s. However, some of the NLS1s are originally from the Zwicky (e.g. I Zw 1) and Markarian (e.g. Mrk 766) samples of galaxies implying that they are blue. While the blue color might be due to big blue bumps in the active nuclei, as in normal Seyfert galaxies, NLS1s have weak blue bumps and so the blue colors might be a result of actively star forming galaxies. Some NLS1s (e.g. IRAS 13349+2438), are infrared luminous, and star forming. Using the galaxy catalogs RC3 (de Vaucouleurs et al. 1991) and UGC (Nilson 1973), we looked into the morphology of a small sample of NLS1s listed in Table 1 and found information on seven of them. Three were found to be compact (I Zw 1, Mrk 507, and Mrk 1298), two showed signatures of an inner ring (NGC 4051 and Ark 564), and three have nuclear bars (NGC 4051, Mrk 766 and Ark 564). These are signatures of recent activity (though not necessarily), quite likely due to galaxy-galaxy interactions or mergers. In this scenario the galaxies are newly formed, or rejuvenated.
That NLS1s reside in young galaxies is also consistent with the hypothesis that the formation and evolution of galaxies and their active nuclei is intimately related (Rees 1997, Fabian 1999, Granato et al. 1999, Haehnelt & Kauffmann 1999). In this scenario, the process of formation of a massive BH and the active nucleus is the very process of galaxy formation. The active nucleus and the galaxy evolve together, with the BH accreting matter and the galaxy making stars. At one stage the winds from the active nucleus blow away the matter surrounding it and a quasar emerges. This is not only the end of the active evolution of the quasar, but of the galaxy as well, as it is evacuated of its interstellar medium. The quasar then shines as long as there is fuel in the accretion disk (Fabian 1999). In this scenario, high redshift quasars represent an early stage of galaxy evolution, BALQSOs at z$`2`$ represent the stage when the gas is being blown away, and z$`1`$ quasars would be the passively evolving population. Massive ellipticals found today might be the dead remnants of what were once quasars.
The quasar phenomenon may thus be a result of galaxy formation due to primordial density fluctuations. At low redshift, when new galaxies are formed due to interactions or mergers, similar evolution may take place. As argued above, the NLS1s may represent a crucial early phase. (In our scenario, the accretion rate ṁ = Ṁ/Ṁ<sub>Edd</sub> is large in the early stages of evolution and reduces later on. This is opposite to the proposal by Wandel (1999) in which ṁ increases with time.)
In fact, there might be some NLS1s with a starburst component. Soft X-ray spectra of NLS1 are steep and often variable. However, Leighly et al. (1996) and Page et al. (1999) report that while the power-law component in the NLS1 Mrk 766 varied, the thermal black-body component did not. This component might well be due to a nuclear starburst. Note also the strong CO emission in the prototype NLS1 I Zw 1 (Barvainis, Alloin & Antonucci 1989). Schinnerer, Eckart & Tacconi (1998) mapped I Zw 1 in CO and found a circumnuclear ring of diameter 1.8 kpc. The authors found strong evidence for a nuclear starburst. There is also a companion to I Zw 1, supporting an interpretation of the starburst being due to an interaction. Similarly, AGN activity is known to exist in starburst galaxies (see Heckman 1999 for a review). The poster by Dennefeld et al. reports observations of narrow optical emission lines in a sample of IR selected starburst galaxies. See Mathur (2000) for the details of the connection between NLS1s and the evolution of galaxies and active galaxies.
## 6 On strong Fe II emission in NLS1s.
There is a general consensus that large accretion rate, ṁ, is a possible driver of many of the observed properties of NLS1s. Strength of Fe II emission may also be linked to the large accretion rate. In a model by Kwan et al. (1995), Fe II line emission is produced in an accretion disk. The accretion disks with larger accretion rate may simply have more mass to produce stronger Fe II. Thus I would like to argue in favor of collisional ionization as the origin of the Fe II emission in NLS1s and AGNs in general (see also the article by S. Collin in these proceedings). Remember also the CIII$`\lambda `$1174.7 observed in PG 1404+266. This line is also seen in the HST spectrum of another NLS1, I Zw 1 (Laor et al. 1997). Laor et al. preferred the explanation that it is produced by resonance scattering of continuum photons. However, in this scenario the velocity gradient in each emitting cloud and the total covering factor would be exceptionally large. The presence of the CIII$`\lambda `$1174.7 line, together with strong Fe II emission, suggests a collisional origin.
## 7 Conclusions
That the AGN phenomenon was so much stronger at z$``$2–3 than today has long elicited the suspicion that there is a connection between the youth of a galaxy and the likelihood that an AGN forms inside it. The question then naturally arises, “what are the local counterparts to the young galaxies in the early universe in which local AGN may live?” (Krolik 1999). A standard answer to this question is “Starburst galaxies”. Heckman (1999) has argued that starburst galaxies are the low redshift analogues of Lyman break galaxies at high redshift. Similarly, we ask, what are the low redshift analogues of high redshift (z$`\genfrac{}{}{0pt}{}{_>}{^{}}`$4) quasars? I propose that they might be NLS1s.
It is my pleasure to thank Th. Boller and the organizing committee for inviting me to this wonderfully stimulating workshop. I thank the Wilhelm and Else Heraeus foundation for travel support and the delightful stay at the Physikzentrum, Bad Honnef. This work is supported in part through NASA grant NAG 5-3249 (LTSA). The figures were created using the archives at the Space Telescope Science Institute, operated by the Association of Universities for Research in Astronomy, Inc., from NASA contract NAS5-26555.
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# 1 Introduction
## 1 Introduction
Search for the Higgs bosons and the precise measurements on their properties, such as the masses, the decay widths and the decay branching ratios, are the most important subjects to study the mechanism of the electroweak symmetry breaking . In the Standard Model (SM), only one physical neutral Higgs boson appears. On the other hand, models with multiple Higgs doublets have CP-even and odd neutral Higgs bosons as well as charged Higgs bosons, if CP is a good symmetry. Otherwise, the neutral Higgs bosons do not have to carry any definite CP–parity.
CP violation was observed in the neutral kaon system and the violation in B–meson decays is strongly suggested by recent experiments . In addition, CP violation constitutes one of the crucial ingredients for an efficient non–SM generation of the cosmological baryon asymmetry at the electroweak scale . An appealing scheme of CP violation beyond the SM can be provided by models with an extended Higgs sector where the CP asymmetry is broken by the ground state of the Higgs potential. It has recently been realized that, even though the Higgs potential of the minimal supersymmetric SM (MSSM) is CP–invariant at the tree level, explicit CP violation in the mass matrices of the third generation squarks can induce sizable CP violation in the MSSM Higgs sector through loop corrections . The CP violating phases for the third generation sfermions can be quite large, since they contribute to the electric dipole moments (EDM’s) of the electron and neutron only at the two–loop level with no generation mixing <sup>*</sup><sup>*</sup>*The CP–violating phases associated with the sfermions of the first and second generations can be severely constrained by the EDM constraints. However, there have been several suggestions to evade these constraints without suppressing the CP–violating phases. in the sfermion sector . As a result, although a one–loop effect, the induced CP violation in the MSSM Higgs sector can be large enough to significantly affect Higgs phenomenology at present and future colliders . In this light, it is very important to examine the possibility of observing these Higgs bosons and investigating their properties in detail.
A muon collider is one of the ideal machines to look for the Higgs bosons and their properties, where $`\mu ^+\mu ^{}`$ pairs can be directly converted into neutral Higgs–boson resonances . Polarized muons can be used to detect the Higgs bosons and measure their properties such as masses, total widths, decay branching fractions and CP parities . In this paper, we present a general formalism and a detailed analysis for the Higgs–boson effects on the polarization observables of the production process $`\mu ^+\mu ^{}f\overline{f}`$ $`(f=\tau ^{},b`$ and $`t`$) using the initial muon beam polarization along with the unpolarized final fermions and with the final–fermion polarization configuration of equal helicity, respectively. It has been recently pointed out that the amplitudes for two Higgs states, a scalar and a pseudoscalar, can interfere with each other sizably if the helicities of the initial and final particles are properly fixed and if the mass difference of these Higgs bosons is at most of the same order as their decay widths . We extend the work comprehensively in order to consider Higgs bosons of no definite CP–parity in the MSSM with explicit CP violation induced from the third generation squark sectors.
The remainder of this article is organized as follows. Section 2 is devoted to a brief review of the explicit CP violation in the MSSM Higgs sector based on the work . In section 3 we present the helicity amplitudes of the production of a fermion pair in $`\mu ^+\mu ^{}`$ collisions and give a comprehensive classification of the observables according to the CP and CP$`\stackrel{~}{\mathrm{T}}`$ transformation properties that are constructed by the initial muon beam polarization along with the final fermions of no polarization and of equal helicity, respectively. In section 4 we make a detailed numerical analysis based on a given parameter set so as to get a concrete estimate of the usefulness of the observables. Section 5 is devoted to a brief summary of our findings and to a conclusion.
## 2 Explicit CP Violation in the MSSM Higgs Sector
The existence of the non–trivial CP–violating phases in the MSSM is due to the breakdown of supersymmetry so that the CP–violating phases appear in the soft–breaking parameters and the mixing among sparticles due to the electroweak gauge symmetry breaking. There are three well–known sources of CP violation in the MSSM . The first is related to the two Higgs–boson doublets present in the model since both the $`\mu `$ parameter in the superpotential and the soft breaking parameter $`m_{12}^2`$ can be complex. Secondly, there are three more phases of the complex masses of the U(1)<sub>Y</sub>, SU(2)<sub>L</sub> and SU(3)<sub>C</sub> gauginos of the SM gauge group. Thirdly, the other CP–violating phases originate from the flavor sector of the MSSM Lagrangian, either in the scalar soft mass matrices or the trilinear matrices.
The CP–violating phases associated with the off-diagonal terms of the trilinear matrices are, however, strongly suppressed by the same mechanism required to suppress the flavor changing neutral current effects. Therefore, all these flavor–changing CP–violating phases are neglected in the present work such that the scalar soft mass matrices and trilinear parameters are flavor diagonal and the complex trilinear terms are proportional to the corresponding fermion Yukawa couplings. Clearly, the Yukawa interactions of the third–generation quarks and squarks play the most significant role in radiative corrections to the Higgs sector. In this section, we give a brief review of the calculation of the Higgs–boson mass matrix based on the full one–loop effective potential, valid for all values of the relevant third–generation soft–breaking parameters.
The MSSM contains two Higgs doublets $`H_1,H_2`$, with hypercharges $`Y(H_1)=Y(H_2)=1/2`$. Here we are only interested in the neutral components, which we write as
$`H_1^0={\displaystyle \frac{1}{\sqrt{2}}}\left(\varphi _1+ia_1\right);H_2^0={\displaystyle \frac{\mathrm{e}^{i\xi }}{\sqrt{2}}}\left(\varphi _2+ia_2\right),`$ (1)
where $`\varphi _{1,2}`$ and $`a_{1,2}`$ are real fields. The constant phase $`\xi `$ can be set to zero at tree level, but will in general become non–zero once loop corrections are included.
The mass matrix of the neutral Higgs bosons can be computed from the effective potential
$`V_{\mathrm{Higgs}}`$ $`={\displaystyle \frac{1}{2}}m_1^2\left(\varphi _1^2+a_1^2\right)+{\displaystyle \frac{1}{2}}m_2^2\left(\varphi _2^2+a_2^2\right)\left|m_{12}^2\right|\left(\varphi _1\varphi _2a_1a_2\right)\mathrm{cos}(\xi +\theta _{12})`$ (2)
$`\left|m_{12}^2\right|\left(\varphi _1a_2+\varphi _2a_1\right)\mathrm{sin}(\xi +\theta _{12})+{\displaystyle \frac{\widehat{g}^2}{8}}𝒟^2+{\displaystyle \frac{1}{64\pi ^2}}\mathrm{Str}\left[^4\left(\mathrm{log}{\displaystyle \frac{^2}{Q^2}}{\displaystyle \frac{3}{2}}\right)\right],`$
where we have allowed the soft breaking parameter $`m_{12}^2=\left|m_{12}^2\right|\mathrm{e}^{i\theta _{12}}`$ to be complex, and we have introduced the quantities
$`𝒟=\varphi _2^2+a_2^2\varphi _1^2a_1^2;\widehat{g}^2={\displaystyle \frac{g^2+g^2}{4}},`$ (3)
where the symbols $`g`$ and $`g^{}`$ stand for the SU(2)<sub>L</sub> and U(1)<sub>Y</sub> gauge couplings, respectively. $`Q`$ in Eq. (2) is the renormalization scale; the parameters of the tree–level potential, in particular the mass parameters $`m_1^2,m_2^2`$ and $`m_{12}^2`$, are running parameters, taken at scale $`Q`$. The potential (2) is then independent of $`Q`$, up to two–loop corrections.
The matrix $``$ is the field–dependent mass matrix of all modes that couple to the Higgs bosons. The by far dominant contributions come from the third generation quarks and squarks. The (real) masses of the former are given by
$`m_b^2={\displaystyle \frac{1}{2}}|h_b|^2\left(\varphi _1^2+a_1^2\right);m_t^2={\displaystyle \frac{1}{2}}|h_t|^2\left(\varphi _2^2+a_2^2\right),`$ (4)
where $`h_b`$ and $`h_t`$ are the bottom and top Yukawa couplings. The corresponding squark mass matrices can be written as
$`_{\stackrel{~}{t}}^2`$ $`=\left(\begin{array}{cc}m_{\stackrel{~}{Q}}^2+m_t^2\frac{1}{8}\left(g^2\frac{g^2}{3}\right)𝒟& h_t^{}\left[A_t^{}\left(H_2^0\right)^{}+\mu H_1^0\right]\\ h_t\left[A_tH_2^0+\mu ^{}\left(H_1^0\right)^{}\right]& m_{\stackrel{~}{U}}^2+m_t^2\frac{g^2}{6}𝒟\end{array}\right);`$ (7)
$`_{\stackrel{~}{b}}^2`$ $`=\left(\begin{array}{cc}m_{\stackrel{~}{Q}}^2+m_b^2+\frac{1}{8}\left(g^2+\frac{g^2}{3}\right)𝒟& h_b^{}\left[A_b^{}\left(H_1^0\right)^{}+\mu H_2^0\right]\\ h_b\left[A_bH_1^0+\mu ^{}\left(H_2^0\right)^{}\right]& m_{\stackrel{~}{D}}^2+m_b^2+\frac{g^2}{12}𝒟\end{array}\right).`$ (10)
Here, $`H_1^0`$ and $`H_2^0`$ are given by Eq. (1) while $`m_t^2`$ and $`m_b^2`$ are as in Eq. (4) and $`𝒟`$ has been defined in Eq. (3). In Eq. (10) $`m_{\stackrel{~}{Q}}^2,m_{\stackrel{~}{U}}^2`$ and $`m_{\stackrel{~}{D}}^2`$ are real soft breaking parameters, $`A_b`$ and $`A_t`$ are complex soft breaking parameters, and $`\mu `$ is the complex supersymmetric Higgs(ino) mass parameter.
The calculation proceeds by plugging the field–dependent top/bottom–(s)quark mass eigenvalues into the potential (2). The mass matrix of the Higgs bosons (at vanishing external momentum) is then given by the matrix of second derivatives of this potential, computed at its minimum. In order to make sure that we are indeed in the minimum of the potential, we solve the stationarity relations, i.e. set the first derivatives of the potential to zero. This allows us to, e.g., express $`m_1^2,m_2^2`$ and $`m_{12}^2\mathrm{sin}(\xi +\theta _{12})`$ as functions of the vacuum expectation values (vevs) and the remaining parameters appearing in the loop–corrected Higgs potential. The equations $`V_{\mathrm{Higgs}}/a_1=0`$ and $`V_{\mathrm{Higgs}}/a_2=0`$ are linearly dependent, i.e. lead to only one constraint on parameters for $`a_1=a_2=0`$; the remaining vevs are defined through $`\varphi _1^2+\varphi _2^2=v^2(246\mathrm{GeV})^2`$ and $`\varphi _2/\varphi _1=\mathrm{tan}\beta `$. The re–phasing invariant quantity $`|m_{12}^2|\mathrm{sin}(\xi +\theta _{12})`$ is then determined by $`|m_{12}^2|`$, $`m_{\stackrel{~}{t}_i}^2`$ and $`m_{\stackrel{~}{b}_i}^2`$ as well as by the re–phasing invariant quantities
$`\mathrm{\Delta }_{\stackrel{~}{t}}={\displaystyle \frac{\mathrm{}\mathrm{m}(A_t\mu \mathrm{e}^{i\xi })}{m_{\stackrel{~}{t}_2}^2m_{\stackrel{~}{t}_1}^2}};\mathrm{\Delta }_{\stackrel{~}{b}}={\displaystyle \frac{\mathrm{}\mathrm{m}(A_b\mu \mathrm{e}^{i\xi })}{m_{\stackrel{~}{b}_2}^2m_{\stackrel{~}{b}_1}^2}},`$ (11)
which describe the amount of CP violation in the squark mass matrices.
The mass matrix of the neutral Higgs bosons can now be computed from the matrix of second derivatives of the potential (2), where (after taking the derivatives) $`m_1^2,m_2^2`$ and $`m_{12}^2\mathrm{sin}(\xi +\theta _{12})`$ are determined by the stationarity conditions. The massless state $`G^0=a_1\mathrm{cos}\beta a_2\mathrm{sin}\beta `$ is the would–be Goldstone mode “eaten” by the longitudinal $`Z`$ boson. We are thus left with a squared mass matrix $`_H^2`$ for the three states $`a=a_1\mathrm{sin}\beta +a_2\mathrm{cos}\beta ,\varphi _1`$ and $`\varphi _2`$. This matrix is real and symmetric, i.e. it has 6 independent entries. The diagonal entry for $`a`$ reads:
$`_H^2|_{aa}=m_A^2+{\displaystyle \frac{3}{8\pi ^2}}\left\{{\displaystyle \frac{|h_t|^2m_t^2}{\mathrm{sin}^2\beta }}g(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{t}_2}^2)\mathrm{\Delta }_{\stackrel{~}{t}}^2+{\displaystyle \frac{|h_b|^2m_b^2}{\mathrm{cos}^2\beta }}g(m_{\stackrel{~}{b}_1}^2,m_{\stackrel{~}{b}_2}^2)\mathrm{\Delta }_{\stackrel{~}{b}}^2\right\},`$ (12)
and the CP–violating entries of the mass matrix, which mix $`a`$ with $`\varphi _1`$ and $`\varphi _2`$ read:
$`_H^2|_{a\varphi _1}`$ $`=`$ $`{\displaystyle \frac{3}{16\pi ^2}}\{{\displaystyle \frac{m_t^2\mathrm{\Delta }_{\stackrel{~}{t}}}{\mathrm{sin}\beta }}[g(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{t}_2}^2)(X_t\mathrm{cot}\beta 2|h_t|^2R_t)\widehat{g}^2\mathrm{cot}\beta \mathrm{log}{\displaystyle \frac{m_{\stackrel{~}{t}_2}^2}{m_{\stackrel{~}{t}_1}^2}}]`$
$`+{\displaystyle \frac{m_b^2\mathrm{\Delta }_{\stackrel{~}{b}}}{\mathrm{cos}\beta }}[g(m_{\stackrel{~}{b}_1}^2,m_{\stackrel{~}{b}_2}^2)(X_b+2|h_b|^2R_b^{})+(\widehat{g}^22|h_b|^2)\mathrm{log}{\displaystyle \frac{m_{\stackrel{~}{b}_2}^2}{m_{\stackrel{~}{b}_1}^2}}]\},`$
$`_H^2|_{a\varphi _2}`$ $`=`$ $`{\displaystyle \frac{3}{16\pi ^2}}\{{\displaystyle \frac{m_t^2\mathrm{\Delta }_{\stackrel{~}{t}}}{\mathrm{sin}\beta }}[g(m_{\stackrel{~}{t}_1}^2,m_{\stackrel{~}{t}_2}^2)(X_t+2|h_t|^2R_t^{})+(\widehat{g}^22|h_t|^2)\mathrm{log}{\displaystyle \frac{m_{\stackrel{~}{t}_2}^2}{m_{\stackrel{~}{t}_1}^2}}]`$
$`+{\displaystyle \frac{m_b^2\mathrm{\Delta }_{\stackrel{~}{b}}}{\mathrm{cos}\beta }}[g(m_{\stackrel{~}{b}_1}^2,m_{\stackrel{~}{b}_2}^2)(X_b\mathrm{tan}\beta 2|h_b|^2R_b)\widehat{g}^2\mathrm{tan}\beta \mathrm{log}{\displaystyle \frac{m_{\stackrel{~}{b}_2}^2}{m_{\stackrel{~}{b}_1}^2}}]\}.`$
where $`\mathrm{\Delta }_{\stackrel{~}{t}}`$ and $`\mathrm{\Delta }_{\stackrel{~}{b}}`$ are as in Eq. (11) and the function $`g(m_1^2,m_2^2)`$ is given by
$`g(m_1^2,m_2^2)=2{\displaystyle \frac{m_1^2+m_2^2}{m_1^2m_2^2}}\mathrm{log}{\displaystyle \frac{m_1^2}{m_2^2}}.`$ (15)
The definition of the mass squared $`m_A^2`$ and the dimensionless quantities $`X_{t,b}`$, $`R_{t,b}`$ and $`R_{t,b}^{}`$ as well as the other CP–preserving entries of the mass matrix squared $`_H^2`$ can be found in Ref. . As noted earlier, the size of these CP–violating entries is controlled by $`\mathrm{\Delta }_{\stackrel{~}{t}}`$ and $`\mathrm{\Delta }_{\stackrel{~}{b}}`$.
The real, symmetric matrix $`_H^2`$ can be diagonalized with an orthogonal rotation $`O`$;
$`O^T_H^2O=\mathrm{diag}(m_{H_1}^2,m_{H_2}^2,m_{H_3}^2),`$ (16)
with the ordering of $`m_{H_1}m_{H_2}m_{H_3}`$ taken as a convention in the present work. Note that the loop–corrected neutral Higgs–boson sector is determined by fixing the values of various parameters; $`m_A`$, $`\mu `$, $`A_t`$, $`A_b`$, a renormalization scale $`Q`$, $`\mathrm{tan}\beta `$, and the soft–breaking third generation sfermion masses, $`m_{\stackrel{~}{Q}}`$, $`m_{\stackrel{~}{U}}`$, and $`m_{\stackrel{~}{D}}`$. The radiatively induced phase $`\xi `$ is no more an independent parameter and it can be absorbed into the definition of the $`\mu `$ parameter so that the physically meaningful CP phases in the Higgs sector are the phases of the re–phasing invariant combinations $`A_t\mu \mathrm{e}^{i\xi }`$ and $`A_b\mu \mathrm{e}^{i\xi }`$. This neutral Higgs–boson mixing changes the couplings of the Higgs fields to fermions, gauge bosons, and Higgs fields themselves so that the effects of CP violation in the Higgs sector can be probed through various processes .
## 3 Third Generation Fermion Production
### 3.1 Production amplitudes
The couplings of the $`\gamma `$ and the neutral gauge boson $`Z`$ to fermions in the MSSM is described by the same interaction Lagrangian as in the SM:
$`_{V\mu \mu }=eQ_f\overline{f}\gamma _\mu fA^\mu {\displaystyle \frac{e}{s_Wc_W}}\overline{f}\left[(T_{f3}Q_fs_W^2)P_{}Q_fs_W^2P_+\right]\gamma _\mu fZ^\mu ,`$ (17)
with the chirality projection operators $`P_\pm =(1\pm \gamma _5)/2`$, and those of the neutral Higgs–boson fields with leptons and quarks are described by the interaction Lagrangians
$`_{Hff}`$ $`=`$ $`{\displaystyle \frac{h_l}{\sqrt{2}}}{\displaystyle \underset{k=1}{\overset{3}{}}}\overline{\mathrm{}}\left[O_{2k}is_\beta O_{1k}\gamma _5\right]\mathrm{}H_k{\displaystyle \frac{h_d}{\sqrt{2}}}{\displaystyle \underset{k=1}{\overset{3}{}}}\overline{d}\left[O_{2k}is_\beta O_{1k}\gamma _5\right]dH_k`$ (18)
$`{\displaystyle \frac{h_u}{\sqrt{2}}}{\displaystyle \underset{k=1}{\overset{3}{}}}\overline{u}\left[O_{3k}ic_\beta O_{1k}\gamma _5\right]uH_k,`$
where $`h_l`$, $`h_d`$, and $`h_u`$ are the lepton and quark Yukawa couplings:
$`h_l={\displaystyle \frac{gm_l}{\sqrt{2}m_Wc_\beta }};h_d={\displaystyle \frac{gm_d}{\sqrt{2}m_Wc_\beta }};h_u={\displaystyle \frac{gm_u}{\sqrt{2}m_Ws_\beta }},`$ (19)
respectively. It is then clear that all the neutral Higgs bosons couple dominantly to the third generation fermions $`t`$, $`b`$ and $`\tau `$ and they couple to a muon about 200 times more strongly than to an electron – the primary reason for having a muon collider.
Figure 1: The mechanisms contributing to the process $`\mu ^+\mu ^{}f\overline{f}`$; three spin–0 neutral–Higgs–boson exchanges and spin–1 $`\gamma `$ and $`Z`$ exchanges. Here, the index $`k`$ is 1,2 or 3.
The fermion pair production process $`\mu ^+\mu ^{}f\overline{f}`$ is generated by five $`s`$–channel mechanisms: spin–1 $`\gamma `$ and $`Z`$ exchanges and three spin–0 neutral–Higgs–boson exchanges, cf. Fig. 1. The transition matrix element of the process
$`={\displaystyle \frac{S_{ab}}{s}}\left[\overline{v}(\mu ^+)P_au(\mu ^{})\right]\left[\overline{u}(f)P_bv(\overline{f})\right]+{\displaystyle \frac{V_{ab}}{s}}\left[\overline{v}(\mu ^+)\gamma _\mu P_au(\mu ^{})\right]\left[\overline{u}(f)\gamma ^\mu P_bv(\overline{f})\right],`$ (20)
can be expressed in terms of the scalar and vector bilinear charges, $`S_{ab}`$ and $`V_{ab}`$, classified according to the chiralities $`a,b=\pm `$ for the right–/left–handed chiralities of the associated muon and produced fermion currents, respectively. Here, the scalar bilinear charges are given by:
$`S_{ab}={\displaystyle \frac{h_\mu h_f}{2}}{\displaystyle \underset{k=1}{\overset{3}{}}}D_{H_k}(s)\left[O_{2k}ias_\beta O_{1k}\right]\left[S_k^fib\xi _\beta O_{1k}\right]{\displaystyle \frac{e^2m_\mu m_f}{4m_W^2s_W^2}}D_Z(s)T_{f3}ab,`$ (21)
where $`h_\mu `$ and $`h_f`$ are the initial muon and final fermion Yukawa couplings, respectively, and
$`D_{H_k}(s)={\displaystyle \frac{s}{sm_{H_k}^2+im_{H_k}\mathrm{\Gamma }_{H_k}}},D_Z(s)={\displaystyle \frac{s}{sm_Z^2+im_Z\mathrm{\Gamma }_Z}},`$
$`S_k^f=\{\begin{array}{c}O_{2k}\hfill \\ O_{3k}\hfill \end{array},\xi _\beta =\{\begin{array}{cc}s_\beta \hfill & \mathrm{for}f=l\mathrm{or}d\hfill \\ c_\beta \hfill & \mathrm{for}f=u\hfill \end{array}.`$ (26)
On the other hand, the vector bilinear charges $`V_{ab}`$ have the contributions only from the $`\gamma `$ and $`Z`$ boson exchanges:
$`V_{ab}=e^2\left[Q_f+g_a^\mu g_b^fD_Z(s)\right],`$ (27)
where $`Q_f`$ is the electric charge of the fermion, $`g_\pm ^f`$ denote the right/left–handed couplings of the $`Z`$ boson to the fermions, respectively:
$`g_+^f=Q_f\mathrm{tan}\theta _W,g_{}^f=Q_f\mathrm{tan}\theta _W+{\displaystyle \frac{T_{3f}}{s_Wc_W}}.`$ (28)
The vector bilinear charges $`V_{ab}`$ are real in the approximation of neglecting the $`Z`$ boson width $`\mathrm{\Gamma }_Z`$, which is valid for the c.m. energy away from the $`Z`$ boson pole.
Figure 2: The schematic description of the production plane with the scattering angles $`\mathrm{\Theta }`$ and $`\mathrm{\Phi }`$ as well as the transverse polarization vectors $`P_T`$ and $`\overline{P}_T`$ with the azimuthal angles $`\alpha `$ and $`\overline{\alpha }`$ with respect to the scattering angle, respectively.
Defining the polar angle of the flight direction of the fermion $`f`$ with respect to the $`\mu ^{}`$ beam direction by $`\mathrm{\Theta }`$ (See Fig. 2), the explicit form of the production amplitude (20) can be evaluated in the helicity basis by the 2–component spinor technique of Ref. . Denoting the $`\mu ^{}(\mu ^+)`$ helicity by the first (second) index, the $`f`$ and $`\overline{f}`$ helicities by the remaining two indices, then the general form of the helicity amplitude $`\sigma \overline{\sigma };\lambda \overline{\lambda }`$, consisting of a scalar helicity amplitude and a vector helicity amplitude, reads
$`\sigma \overline{\sigma };\lambda \overline{\lambda }=\sigma \overline{\sigma };\lambda \overline{\lambda }__S+\sigma \overline{\sigma };\lambda \overline{\lambda }__V,`$ (29)
where the scalar helicity amplitude $`\sigma \overline{\sigma };\lambda \overline{\lambda }__S`$ is given by
$`\sigma \overline{\sigma };\lambda \overline{\lambda }__S={\displaystyle \frac{1}{4}}{\displaystyle \underset{ab}{}}S_{ab}(a+\sigma \beta _\mu )(b\lambda \beta )\delta _{\sigma \overline{\sigma }}\delta _{\lambda \overline{\lambda }},`$ (30)
and the vector helicity amplitude $`\sigma \overline{\sigma };\lambda \overline{\lambda }__V`$ by
$`\sigma \overline{\sigma };\lambda \overline{\lambda }__V`$ $`=`$ $`{\displaystyle \frac{1}{4}}{\displaystyle \underset{ab}{}}V_{ab}\{(1+a\sigma \beta _\mu )(1+b\lambda \beta )(\lambda \sigma +\mathrm{cos}\mathrm{\Theta })\delta _{\sigma ,\overline{\sigma }}\delta _{\lambda ,\overline{\lambda }}`$ (31)
$`{\displaystyle \frac{4m_\mu m_f}{s}}(ab\sigma \lambda \mathrm{cos}\mathrm{\Theta })\delta _{\sigma \overline{\sigma }}\delta _{\lambda \overline{\lambda }}`$
$`{\displaystyle \frac{2m_\mu }{\sqrt{s}}}(1+b\lambda \beta )(\sigma \mathrm{sin}\mathrm{\Theta })\delta _{\sigma \overline{\sigma }}\delta _{\lambda ,\overline{\lambda }}`$
$`+{\displaystyle \frac{2m_f}{\sqrt{s}}}(1+a\sigma \beta _\mu )(\lambda \mathrm{sin}\mathrm{\Theta })\delta _{\sigma ,\overline{\sigma }}\delta _{\lambda \overline{\lambda }}\},`$
respectively, with $`a,b=\pm `$, $`\beta _\mu =\sqrt{14m_\mu ^2/s}`$ and $`\beta =\sqrt{14m_f^2/s}`$. However, the muon mass $`m_\mu `$ ($`=106`$ MeV) is extremely small compared to the present experimental mass bound of approximately 100 GeV on the lightest Higgs boson so that one can safely neglect all the terms involving the kinematical muon massFor consistency, the chirality–flipped $`Z`$–boson contributions to the scalar bilinear charges $`S_{ab}`$, which are also proportional to the muon mass kinematically, should be neglected in the massless muon limit.. In the approximation, the scalar and vector helicity amplitudes can be written as:
(i) Scalar helicity amplitudes:
$`++;++__S=\beta _+S_+\beta _{}S_{++},`$
$`++;__S=\beta _{}S_+\beta _+S_{++},`$
$`;++__S=\beta _{}S_+\beta _+S_{},`$
$`;__S=\beta _+S_+\beta _{}S_{},`$ (32)
where $`\beta _\pm =(1\pm \beta )/2`$. At asymptotically high energies $`\beta _+1`$ and $`\beta _{}0`$, and the other scalar helicity amplitudes vanish.
(ii) Vector helicity amplitudes:
$`+;++__V=\sqrt{\beta _+\beta _{}}(V_{++}+V_+)\mathrm{sin}\mathrm{\Theta },`$
$`+;+__V=[\beta _+V_{++}+\beta _{}V_+](1+\mathrm{cos}\mathrm{\Theta }),`$
$`+;+__V=+[\beta _{}V_{++}+\beta _+V_+](1\mathrm{cos}\mathrm{\Theta }),`$
$`+;__V=+\sqrt{\beta _+\beta _{}}(V_{++}+V_+)\mathrm{sin}\mathrm{\Theta },`$
$`+;++__V=\sqrt{\beta _+\beta _{}}(V_++V_{})\mathrm{sin}\mathrm{\Theta },`$
$`+;+__V=+[\beta _+V_++\beta _{}V_{}](1\mathrm{cos}\mathrm{\Theta }),`$
$`+;+__V=[\beta _{}V_++\beta _+V_{}](1+\mathrm{cos}\mathrm{\Theta }),`$
$`+;__V=+\sqrt{\beta _+\beta _{}}(V_++V_{})\mathrm{sin}\mathrm{\Theta },`$ (33)
and again the other vector helicity amplitudes vanish.
The CP transformation leads to the relation among the transition helicity amplitudes:
$`\sigma \overline{\sigma };\lambda \overline{\lambda }`$ $`\stackrel{\mathrm{CP}}{}`$ $`+(1)^{(\sigma \overline{\sigma })/2}(1)^{(\lambda \overline{\lambda })/2}\overline{\sigma },\sigma ;\overline{\lambda },\lambda ,`$ (34)
or equivalently for the scalar and vector helicity amplitudes:
$`\pm \pm ;\lambda \overline{\lambda }__S`$ $`\stackrel{\mathrm{CP}}{}`$ $`+(1)^{(\lambda \overline{\lambda })/2};\overline{\lambda },\lambda __S,`$
$`\pm ;\lambda \overline{\lambda }__V`$ $`\stackrel{\mathrm{CP}}{}`$ $`(1)^{(\lambda \overline{\lambda })/2}\pm ;\overline{\lambda },\lambda __V.`$ (35)
Only the simultaneous presence of the scalar and pseudoscalar couplings can lead to CP violation as can be explicitly checked in Eqs. (30) and (31). One additional useful classification is provided by the so–called “naive” time reversal $`\stackrel{~}{\mathrm{T}}`$; under the CP$`\stackrel{~}{\mathrm{T}}`$ transformations the helicity amplitudes are transformed as follows:
$`\pm \pm ;\lambda \overline{\lambda }__S`$ $`\stackrel{\mathrm{CP}\stackrel{~}{\mathrm{T}}}{}`$ $`+(1)^{(\lambda \overline{\lambda })/2};\overline{\lambda },\lambda __S^{},`$
$`\pm ;\lambda \overline{\lambda }__V`$ $`\stackrel{\mathrm{CP}\stackrel{~}{\mathrm{T}}}{}`$ $`(1)^{(\lambda \overline{\lambda })/2}\pm ;\overline{\lambda },\lambda __V^{}.`$ (36)
We note that it is crucial to have finite $`Z`$ or Higgs–boson widths for CP$`\stackrel{~}{\mathrm{T}}`$ violation. In this light, it is very useful to take into account the CP and CP$`\stackrel{~}{\mathrm{T}}`$ properties of any physical observable simultaneously so as to investigate not only CP violation itself but also the overlapping of any pair of three neutral Higgs boson resonances effectively.
### 3.2 Polarized production cross section
The matrix element squared for general (longitudinal or transverse) beam polarization can be computed either using standard trace techniques (employing general spin projection operators), or from the helicity amplitudes by a suitable rotation from the helicity basis to a general spin basis. In the former case, neglecting the muon mass in the spin projection operators we can obtain the following approximated form for the $`\mu ^{}`$ projection operators
$`{\displaystyle \frac{1}{2}}(\overline{)}p+m)(1+\gamma _5\overline{)}s){\displaystyle \frac{1}{2}}(1+P_L\gamma _5)\overline{)}p+{\displaystyle \frac{1}{2}}\gamma _5P_T(\mathrm{cos}\alpha \overline{)}n_1+\mathrm{sin}\alpha \overline{)}n_2)\overline{)}p,`$
$`{\displaystyle \frac{1}{2}}(\overline{)}\overline{p}m)(1+\gamma _5\overline{)}\overline{s}){\displaystyle \frac{1}{2}}(1\overline{P}_L\gamma _5)\overline{)}\overline{p}+{\displaystyle \frac{1}{2}}\gamma _5\overline{P}_T(\mathrm{cos}\overline{\alpha }\overline{)}n_1+\mathrm{sin}\overline{\alpha }\overline{)}n_2)\overline{)}\overline{p}.`$ (37)
Equivalently in the helicity basis the polarization weighted matrix element squared is given by
$`\overline{{\displaystyle }}||^2={\displaystyle \underset{\sigma \sigma ^{}\overline{\sigma }\overline{\sigma }^{}}{}}_{\sigma \overline{\sigma }}_{\sigma ^{}\overline{\sigma }^{}}^{}\rho _{\sigma \sigma ^{}}^{}\rho _{\overline{\sigma }\overline{\sigma }^{}}^+=\mathrm{Tr}\left[\rho ^+^{}\rho ^T\right],`$ (38)
where $`_{\sigma \overline{\sigma }}`$ ($`\sigma ,\overline{\sigma }=\pm `$) denotes the helicity amplitude for any given production process $`\mu ^{}(\sigma )\mu ^+(\overline{\sigma })X`$ and the $`2\times 2`$ matrices $`\rho ^{}`$ are the polarization density matrices for the initial $`\mu ^{}`$ beams:
$`\rho ^{}={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1+P_L& P_T\mathrm{e}^{i\alpha }\\ P_T\mathrm{e}^{i\alpha }& 1P_L\end{array}\right),\rho ^+={\displaystyle \frac{1}{2}}\left(\begin{array}{cc}1+\overline{P}_L& \overline{P}_T\mathrm{e}^{i\overline{\alpha }}\\ \overline{P}_T\mathrm{e}^{i\overline{\alpha }}& 1\overline{P}_L\end{array}\right).`$ (43)
Here, $`P_L`$ and $`\overline{P}_L`$ are the longitudinal polarizations of the $`\mu ^{}`$ and $`\mu ^+`$ beams, while $`P_T`$ and $`\overline{P}_T`$ are the degrees of transverse polarization with $`\alpha `$ and $`\overline{\alpha }`$ being the azimuthal angles between the transverse polarization vectors and the momentum vector of $`f`$ as shown in Fig. 2.
Applying the projection operators (37) or/and evaluating the trace (38) leads to the polarized matrix element squared of the form
$`\mathrm{\Sigma }{\displaystyle \underset{\lambda \overline{\lambda }}{}}{\displaystyle \underset{\sigma \sigma ^{}\overline{\sigma }\overline{\sigma }^{}}{}}\sigma \overline{\sigma };\lambda \overline{\lambda }\sigma ^{}\overline{\sigma ^{}};\lambda \overline{\lambda }^{}\rho _{\sigma \sigma ^{}}^{}\rho _{\overline{\sigma }\overline{\sigma }^{}}^+`$
$`=\left(1P_L\overline{P}_L\right)C_1+\left(P_L\overline{P}_L\right)C_2`$
$`+\left(1+P_L\overline{P}_L\right)C_3+\left(P_L+\overline{P}_L\right)C_4`$
$`+(P_T\mathrm{cos}\alpha +\overline{P}_T\mathrm{cos}\overline{\alpha })C_5+(P_T\mathrm{sin}\alpha +\overline{P}_T\mathrm{sin}\overline{\alpha })C_6`$
$`+(P_T\mathrm{cos}\alpha \overline{P}_T\mathrm{cos}\overline{\alpha })C_7+(P_T\mathrm{sin}\alpha \overline{P}_T\mathrm{sin}\overline{\alpha })C_8`$
$`+(P_L\overline{P}_T\mathrm{cos}\overline{\alpha }+\overline{P}_LP_T\mathrm{cos}\alpha )C_9+(P_L\overline{P}_T\mathrm{sin}\overline{\alpha }+\overline{P}_LP_T\mathrm{sin}\alpha )C_{10}`$
$`+(P_L\overline{P}_T\mathrm{cos}\overline{\alpha }\overline{P}_LP_T\mathrm{cos}\alpha )C_{11}+(P_L\overline{P}_T\mathrm{sin}\overline{\alpha }\overline{P}_LP_T\mathrm{sin}\alpha )C_{12}`$
$`+P_T\overline{P}_T\left[\mathrm{cos}(\alpha +\overline{\alpha })C_{13}+\mathrm{sin}(\alpha +\overline{\alpha })C_{14}\right]`$
$`+P_T\overline{P}_T\left[\mathrm{cos}(\alpha \overline{\alpha })C_{15}+\mathrm{sin}(\alpha \overline{\alpha })C_{16}\right],`$ (44)
where the coefficients $`C_n`$ ($`n=1`$ \- $`16`$) are defined in terms of the helicity amplitudes by
$`C_1={\displaystyle \frac{1}{4}}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}\left[\right|+;\lambda \overline{\lambda }|^2+\left|+;\lambda \overline{\lambda }|^2\right],C_2={\displaystyle \frac{1}{4}}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}\left[\right|+;\lambda \overline{\lambda }|^2\left|+;\lambda \overline{\lambda }|^2\right],`$
$`C_3={\displaystyle \frac{1}{4}}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}\left[\right|++;\lambda \overline{\lambda }|^2+\left|;\lambda \overline{\lambda }|^2\right],C_4={\displaystyle \frac{1}{4}}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}\left[\right|++;\lambda \overline{\lambda }|^2\left|;\lambda \overline{\lambda }|^2\right],`$
$`C_5={\displaystyle \frac{1}{4}}\mathrm{}\mathrm{e}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}(++;\lambda \overline{\lambda };\lambda \overline{\lambda })(+;\lambda \overline{\lambda }+;\lambda \overline{\lambda })^{},`$
$`C_6={\displaystyle \frac{1}{4}}\mathrm{}\mathrm{m}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}(++;\lambda \overline{\lambda };\lambda \overline{\lambda })(+;\lambda \overline{\lambda }++;\lambda \overline{\lambda })^{},`$
$`C_7={\displaystyle \frac{1}{4}}\mathrm{}\mathrm{e}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}(++;\lambda \overline{\lambda }+;\lambda \overline{\lambda })(+;\lambda \overline{\lambda }++;\lambda \overline{\lambda })^{},`$
$`C_8={\displaystyle \frac{1}{4}}\mathrm{}\mathrm{m}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}(++;\lambda \overline{\lambda }+;\lambda \overline{\lambda })(+;\lambda \overline{\lambda }+;\lambda \overline{\lambda })^{},`$
$`C_9={\displaystyle \frac{1}{4}}\mathrm{}\mathrm{e}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}(++;\lambda \overline{\lambda }+;\lambda \overline{\lambda })(+;\lambda \overline{\lambda }+;\lambda \overline{\lambda })^{},`$
$`C_{10}={\displaystyle \frac{1}{4}}\mathrm{}\mathrm{m}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}(++;\lambda \overline{\lambda }+;\lambda \overline{\lambda })(+;\lambda \overline{\lambda }++;\lambda \overline{\lambda })^{},`$
$`C_{11}={\displaystyle \frac{1}{4}}\mathrm{}\mathrm{e}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}(++;\lambda \overline{\lambda };\lambda \overline{\lambda })(+;\lambda \overline{\lambda }++;\lambda \overline{\lambda })^{},`$
$`C_{12}={\displaystyle \frac{1}{4}}\mathrm{}\mathrm{m}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}(++;\lambda \overline{\lambda };\lambda \overline{\lambda })(+;\lambda \overline{\lambda }+;\lambda \overline{\lambda })^{},`$
$`C_{13}={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{e}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}\left[+;\lambda \overline{\lambda }+;\lambda \overline{\lambda }^{}\right],C_{14}={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{m}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}\left[+;\lambda \overline{\lambda }+;\lambda \overline{\lambda }^{}\right],`$
$`C_{15}={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{e}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}\left[;\lambda \overline{\lambda }++;\lambda \overline{\lambda }^{}\right],C_{16}={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{m}{\displaystyle \underset{\lambda ,\overline{\lambda }=\pm }{}}\left[;\lambda \overline{\lambda }++;\lambda \overline{\lambda }^{}\right].`$
The production cross section is then given in terms of the distribution $`\mathrm{\Sigma }`$ in Eq. (44) by
$`{\displaystyle \frac{\mathrm{d}\sigma }{\mathrm{d}\mathrm{cos}\mathrm{\Theta }\mathrm{d}\mathrm{\Phi }}}={\displaystyle \frac{N_C\beta }{64\pi ^2s}}\mathrm{\Sigma },`$ (46)
where $`N_C`$ is the color factor of the final fermion; 3 for the top or bottom quark and 1 for the tau lepton. The dependence of the distribution $`\mathrm{\Sigma }`$ on the azimuthal angle $`\mathrm{\Phi }`$ of the production plane is encoded in the angles $`\alpha `$ and $`\overline{\alpha }`$. If the azimuthal angle $`\mathrm{\Phi }`$ is measured with respect to the direction of the $`\mu ^{}`$ transverse polarization vector, the $`\mathrm{\Phi }`$ dependence can be exhibited explicitly by taking
$`\alpha =\mathrm{\Phi },\overline{\alpha }=\eta \mathrm{\Phi },`$ (47)
where $`\eta `$ is the rotational invariant difference $`\overline{\alpha }\alpha `$ of the azimuthal angles of the $`\mu ^+`$ and $`\mu ^{}`$ transverse polarization vectors with respect to the production plane. Clearly, only the six observables $`\{C_1,C_2,C_3,C_4,C_{15},C_{16}\}`$ can be measured independently of the azimuthal angle, but the other ten observables, in particular, the observables involving the SV correlations, require the reconstruction of the production plane.
The CP transformation on the polarization vectors of the initial muon beams corresponds to the simultaneous exchanges:
$`P_L\overline{P}_L,P_T\overline{P}_T,\alpha \overline{\alpha }.`$ (48)
The CP relation (48) of the polarization vectors leads to the classification that the first 3 terms and the distributions $`\{C_5,C_6,C_{11},C_{12},C_{13},C_{14},C_{15}\}`$ in Eq. (44) are CP–even while the other six distributions $`\{C_4,C_7,C_8,C_9,C_{10},C_{16}\}`$ are CP–odd. Among the CP–odd observables, the three observables $`\{C_4,C_7,C_9\}`$ are CP$`\stackrel{~}{\mathrm{T}}`$–odd and the other three observables $`\{C_8,C_{10},C_{16}\}`$ are CP$`\stackrel{~}{\mathrm{T}}`$–even. Here, it will be noteworthy to emphasize again that one crucial requirement for having a large CP$`\stackrel{~}{\mathrm{T}}`$–odd observable is the presence of so–called CP–preserving re-scattering phases which can be provided by the Higgs boson propagators with relatively large widths in the present work.
### 3.3 Initial spin correlations
Let us now express the coefficients $`C_i`$ ($`i=1`$ to 16) in terms of the scalar and vector bilinear and quartic charges as well as the newly–introduced notations:
$`R_1={\displaystyle \underset{ab}{}}S_{ab},R_2={\displaystyle \underset{ab}{}}abS_{ab},R_3={\displaystyle \underset{ab}{}}aS_{ab},R_4={\displaystyle \underset{ab}{}}bS_{ab},`$
$`W_1={\displaystyle \underset{ab}{}}V_{ab},W_2={\displaystyle \underset{ab}{}}abV_{ab},W_3={\displaystyle \underset{ab}{}}aV_{ab},W_4={\displaystyle \underset{ab}{}}bV_{ab},`$ (49)
which are to be employed for the interference between the scalar and vector contributions. It is worthwhile to note that both $`R_3`$ and $`R_4`$ vanish in the CP–invariant theory. The sixteen polarization distributions can be then classified as follows:
(i) Scalar–scalar (SS) correlations:
$`C_{\mathrm{\hspace{0.17em}3}}[++]={\displaystyle \frac{1}{2}}[(1+\beta ^2)S_1(1\beta ^2)S_2],`$
$`C_{\mathrm{\hspace{0.17em}4}}[]={\displaystyle \frac{1}{2}}[(1+\beta ^2)S_1^{}(1\beta ^2)S_2^{}],`$
$`C_{15}[++]={\displaystyle \frac{1}{2}}[(1+\beta ^2)S_4(1\beta ^2)S_5],`$
$`C_{16}[+]={\displaystyle \frac{1}{2}}[(1+\beta ^2)S_4^{}(1\beta ^2)S_5^{}],`$ (50)
where the first and second signatures in the square brackets are for the CP and CP$`\stackrel{~}{\mathrm{T}}`$ parities of the corresponding observable, respectively, and the relevant scalar quartic charges are defined as
$`S_1={\displaystyle \frac{1}{4}}\left[|S_{++}|^2+|S_{}|^2+|S_+|^2+|S_+|^2\right],`$
$`S_1^{}={\displaystyle \frac{1}{4}}\left[|S_{++}|^2+|S_+|^2|S_+|^2|S_{}|^2\right],`$
$`S_2={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{e}\left[S_{++}S_+^{}+S_{}S_+^{}\right],S_2^{}={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{e}\left[S_{++}S_+^{}S_{}S_+^{}\right],`$
$`S_4={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{e}\left[S_{++}S_+^{}+S_{}S_+^{}\right],S_4^{}={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{m}\left[S_{++}S_+^{}S_{}S_+^{}\right],`$
$`S_5={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{e}\left[S_{++}S_{}^{}+S_+S_+^{}\right],S_5^{}={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{m}\left[S_{++}S_{}^{}+S_+S_+^{}\right],`$ (51)
(ii) Vector–vector (VV) correlations:
$`C_{\mathrm{\hspace{0.17em}1}}[++]=(1+\beta ^2\mathrm{cos}^2\mathrm{\Theta })V_1+(1\beta ^2)V_2+2\beta V_3\mathrm{cos}\mathrm{\Theta },`$
$`C_{\mathrm{\hspace{0.17em}2}}[++]=(1+\beta ^2\mathrm{cos}^2\mathrm{\Theta })V_1^{}+(1\beta ^2)V_2^{}+2\beta V_3^{}\mathrm{cos}\mathrm{\Theta },`$
$`C_{13}[++]=\beta ^2\mathrm{sin}^2\mathrm{\Theta }V_4,`$
$`C_{14}[+]=\beta ^2\mathrm{sin}^2\mathrm{\Theta }V_4^{},`$ (52)
where the quartic charges $`V_3`$ and $`V_3^{}`$ are given by
$`V_3={\displaystyle \frac{1}{4}}\left[|V_{++}|^2+|V_{}|^2|V_+|^2|V_+|^2\right],`$
$`V_3^{}={\displaystyle \frac{1}{4}}\left[|V_{++}|^2|V_+|^2+|V_+|^2|V_{}|^2\right],`$ (53)
and the other vector quartic charges are defined in the same way as the scalar quartic charges with the notation $`S`$ replaced by $`V`$ everywhere. We note that the scalar as well as vector quartic charges defined as an imaginary part of the bilinear–charge correlations might be non-vanishing only when there are complex CP–violating couplings or/and CP–preserving phases like re-scattering phases or finite widths of the intermediate particles. So, if there are no CP–preserving phases, non–vanishing values of these quartic charges signal CP violation in the given process.
(iii) Scalar–vector (SV) correlations:
$`C_{\mathrm{\hspace{0.17em}5}}[++]=+{\displaystyle \frac{m_f}{4\sqrt{s}}}\beta \mathrm{sin}\mathrm{\Theta }\mathrm{}\mathrm{e}(W_3R_1^{}),C_{\mathrm{\hspace{0.17em}6}}[+]=+{\displaystyle \frac{m_f}{4\sqrt{s}}}\beta \mathrm{sin}\mathrm{\Theta }\mathrm{}\mathrm{m}(W_1R_1^{}),`$
$`C_{\mathrm{\hspace{0.17em}7}}[]={\displaystyle \frac{m_f}{4\sqrt{s}}}\beta \mathrm{sin}\mathrm{\Theta }\mathrm{}\mathrm{e}(W_1R_3^{}),C_{\mathrm{\hspace{0.17em}8}}[+]={\displaystyle \frac{m_f}{4\sqrt{s}}}\beta \mathrm{sin}\mathrm{\Theta }\mathrm{}\mathrm{m}(W_3R_3^{}),`$
$`C_{\mathrm{\hspace{0.17em}9}}[]=+{\displaystyle \frac{m_f}{4\sqrt{s}}}\beta \mathrm{sin}\mathrm{\Theta }\mathrm{}\mathrm{e}(W_3R_3^{}),C_{10}[+]=+{\displaystyle \frac{m_f}{4\sqrt{s}}}\beta \mathrm{sin}\mathrm{\Theta }\mathrm{}\mathrm{m}(W_1R_3^{}),`$
$`C_{11}[++]=+{\displaystyle \frac{m_f}{4\sqrt{s}}}\beta \mathrm{sin}\mathrm{\Theta }\mathrm{}\mathrm{e}(W_1R_1^{}),C_{12}[+]=+{\displaystyle \frac{m_f}{4\sqrt{s}}}\beta \mathrm{sin}\mathrm{\Theta }\mathrm{}\mathrm{m}(W_3R_1^{}).`$ (54)
All the SV correlations are proportional to the mass of the final–state fermion divided by the c.m. energy so that they are strongly suppressed at high energies. In this light, they are sizable only in the production of the top–quark pair because of the largest fermionic mass.
### 3.4 Final spin correlations for Higgs interference effects
Various types of experimental observables by use of the initial muon polarization have been considered in the previous section for determining the CP character of the neutral MSSM Higgs bosons. In addition, spin correlations of $`t\overline{t}`$ or $`\tau ^+\tau ^{}`$ in the final state can also probe the CP nature of the Higgs boson, independently of or combined with initial beam polarization. If the production plane or directly the momentum directions of two final fermions are determined, the secondary decays of the primary final state fermions can allow a complete analysis of their spin or helicity directions. However, even if the production plane is not reconstructed, one can obtain the fermion–pair polarization combination by statistically studying decay products from the correlated and polarized $`f\overline{f}`$. It is then natural that the production distribution needs to be integrated over the azimuthal angle $`\mathrm{\Phi }`$.
The $`s`$–channel Higgs boson (spin–0) exchange populates the equal $`\mu ^{}\mu ^+`$ helicity combinations, i.e. $`\sigma =\overline{\sigma }`$. This results in the correlations of $`f\overline{f}`$ polarization with $`\lambda =\overline{\lambda }`$ by angular momentum conservation, while the SM background channels yield dominantly the $`(+)`$ or $`(+)`$ polarization combination. In addition, the CP transformation interchanges the $`(++)`$ and $`()`$ helicity configurations. In the light of these two arguments combined with the one in the previous paragraph, it will be valuable to consider the $`\mathrm{\Phi }`$–independent polarization observableIn principle, there can exist the $`\gamma `$ and $`Z`$–exchange terms corresponding to the polarization combinations $`(1P_L\overline{P}_L)`$ and $`(P_L\overline{P}_L)`$ of the initial $`\mu ^+\mu ^{}`$ beams. However, those terms turn out to be zero for the real vector and axial–vector couplings in the approximation of neglecting the $`Z`$–boson width.:
$`\mathrm{\Delta }{\displaystyle \underset{\lambda }{}}{\displaystyle \underset{\sigma \sigma ^{}\overline{\sigma }\overline{\sigma }^{}}{}}\lambda \sigma \overline{\sigma };\lambda \lambda \sigma ^{}\overline{\sigma }^{};\lambda \lambda ^{}\rho _{\sigma \sigma ^{}}^{}\rho _{\overline{\sigma }\overline{\sigma }^{}}^+`$
$`=\left(1+P_L\overline{P}_L\right)D_1+\left(P_L+\overline{P}_L\right)D_2+P_T\overline{P}_T\left[\mathrm{cos}\eta D_3+\mathrm{sin}\eta D_4\right],`$ (55)
where the distributions $`D_i`$ ($`i=1`$ to 4) are defined in terms of the helicity amplitudes as
$`D_1[]={\displaystyle \frac{1}{4}}{\displaystyle \underset{\lambda }{}}\lambda [|++;\lambda \lambda |^2+|;\lambda \lambda |^2],`$
$`D_2[++]={\displaystyle \frac{1}{4}}{\displaystyle \underset{\lambda }{}}\lambda [|++;\lambda \lambda |^2|;\lambda \lambda |^2],`$
$`D_3[]={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{e}{\displaystyle \underset{\lambda }{}}\lambda \left[;\lambda \lambda ++;\lambda \lambda ^{}\right],`$
$`D_4[+]={\displaystyle \frac{1}{2}}\mathrm{}\mathrm{m}{\displaystyle \underset{\lambda }{}}\lambda \left[;\lambda \lambda ++;\lambda \lambda ^{}\right],`$ (56)
where the first and second signatures in the square brackets are for the CP and CP$`\stackrel{~}{\mathrm{T}}`$ parities of the corresponding observable.
### 3.5 CP$`\stackrel{~}{\mathrm{T}}`$–even and odd combinations of Higgs–boson propagators
The CP$`\stackrel{~}{\mathrm{T}}`$ parity of each observable plays a crucial role in determining the interference pattern of the Higgs bosons among themselves and with the $`\gamma `$ and $`Z`$ bosons appearing in the observable; every CP$`\stackrel{~}{\mathrm{T}}`$–even (odd) observable involving only the scalar contributions depends on the real (imaginary) part of the combination $`D_{H_k}D_{H_l}^{}`$ of each pair of two Higgs–boson propagators. The CP$`\stackrel{~}{\mathrm{T}}`$–odd observable $`C_{14}`$ involving only the vector contributions is negligible because the $`Z`$–boson width effect is significantly small for the energy far from the $`Z`$–boson pole. When the $`Z`$–boson width is neglected, all the vector couplings are real. As a result, the interference terms between the scalar and vector contributions depends on the real or imaginary parts of the Higgs–boson propagators $`D_{H_k}`$ itself.
For the sake of discussion in the following, let us introduce for the real and imaginary parts the abbreviated notations:
$`𝒮_{kl}\mathrm{}\mathrm{e}[D_{H_k}D_{H_l}^{}]={\displaystyle \frac{(sm_{H_k}^2)(sm_{H_l}^2)+m_{H_k}m_{H_l}\mathrm{\Gamma }_{H_k}\mathrm{\Gamma }_{H_l}}{[(sm_{H_k}^2)^2+m_{H_k}^2\mathrm{\Gamma }_{H_k}^2][(sm_{H_l}^2)^2+m_{H_l}^2\mathrm{\Gamma }_{H_l}^2]}},`$
$`𝒟_{kl}\mathrm{}\mathrm{m}[D_{H_k}D_{H_l}^{}]={\displaystyle \frac{(sm_{H_k}^2)m_{H_l}\mathrm{\Gamma }_{H_l}(sm_{H_l}^2)m_{H_k}\mathrm{\Gamma }_{H_k}}{[(sm_{H_k}^2)^2+m_{H_k}^2\mathrm{\Gamma }_{H_k}^2][(sm_{H_l}^2)^2+m_{H_l}^2\mathrm{\Gamma }_{H_l}^2]}}.`$ (57)
It is worthwhile to note two points; (a) the denominator reveals a typical two–pole structure so that the Higgs–boson contributions are greatly enhanced at the poles; (b) the numerator of $`𝒮_{kl}`$ is negative in the middle of two resonances with a mass splitting larger than their typical widths but positive otherwise, while the numerator of $`𝒟_{kl}`$ is always positive (negative) if $`m_{H_l}m_{H_k}`$ ($`m_{H_k}m_{H_l}`$) and linearly increasing (decreasing ) if $`m_{H_l}\mathrm{\Gamma }_{H_l}m_{H_k}\mathrm{\Gamma }_{H_k}`$ ($`m_{H_k}\mathrm{\Gamma }_{H_k}m_{H_l}\mathrm{\Gamma }_{H_l}`$), respectively.
On the other hand, the SV interference terms depend on the real and imaginary parts of each Higgs–boson propagators $`D_{H_k}`$, the form of which is given by
$`\mathrm{}\mathrm{e}[D_{H_k}]={\displaystyle \frac{sm_{H_K}^2}{(sm_{H_k}^2)^2+m_{H_k}^2\mathrm{\Gamma }_{H_k}^2}},\mathrm{}\mathrm{m}[D_{H_k}]={\displaystyle \frac{m_{H_K}\mathrm{\Gamma }_{H_k}}{(sm_{H_k}^2)^2+m_{H_k}^2\mathrm{\Gamma }_{H_k}^2}}.`$ (58)
Note that the real part changes its sign whenever the c.m. energy crosses the pole, but the imaginary part is always negative.
Combining the coefficients from the mixing matrix elements with those propagator–dependent parts enables us to make a straightforward qualitative understanding of the $`\sqrt{s}`$ dependence of each observable. This will be demonstrated in detail in the following section with a concrete numerical example.
## 4 Numerical Results
We are now ready to present some numerical results. It is known that loop–induced CP violation in the Higgs sector can only be large if both $`|\mu |`$ and $`|A_t|`$ (or $`|A_b|`$, if $`\mathrm{tan}\beta 1`$) are sizable . We therefore choose $`|A_t|=|A_b|=2m_{\stackrel{~}{Q}}`$. For definiteness we will present results only for the $`t\overline{t}`$ mode by taking $`\mathrm{tan}\beta =3,10`$ when probing the region around two heavy Higgs–boson resonances, but we will give a qualitative description of probing the CP property of the lightest Higgs boson. We take the running Yukawa couplings by including the QCD loop effects depending on the gluino mass properly. Since for moderate values of $`\mathrm{tan}\beta `$ the contributions from the (s)bottom sector are still quite small, our results are not sensitive to $`m_{\stackrel{~}{D}}`$ and $`A_b`$; we therefore fix $`m_{\stackrel{~}{D}}=m_{\stackrel{~}{U}}=m_{\stackrel{~}{Q}}`$, although different values for the SU(2) doublet and singlet soft breaking squark masses, $`m_{\stackrel{~}{Q}}m_{\stackrel{~}{U}}`$ are allowed, and also take equal phases for $`A_t`$ and $`A_b`$. Since we are basically interested in distinguishing the CP non–invariant Higgs sector from the CP invariant one, we take for the re-phasing–invariant phase $`\mathrm{\Phi }_{A\mu }`$ of $`A_{t,b}\mu \mathrm{e}^{i\xi }`$ two values; 0 (CP invariant case) and $`\pi /2`$ (maximally CP violating case). On the other hand, we fix for simplicity the other real mass parameters and couplings except for $`\mathrm{tan}\beta `$ as follows:
$`m_A=0.4\mathrm{TeV},|A_{t,b}|=1.0\mathrm{TeV},|\mu |=1.0\mathrm{TeV},`$
$`m_{\stackrel{~}{g}}=0.5\mathrm{TeV},m_{\stackrel{~}{Q}}=m_{\stackrel{~}{U}}=m_{\stackrel{~}{D}}=0.5\mathrm{TeV}.`$ (59)
Finally, we take for the top and bottom quark running masses $`\overline{m}_t(m_t)=165`$ GeV and $`\overline{m}_b(m_b)=4.2`$ GeV. For the given $`m_A`$ much larger than $`m_Z`$, two neutral Higgs bosons have almost degenerate masses:
$`\mathrm{tan}\beta =\mathrm{\hspace{0.17em}\hspace{0.17em}\hspace{0.17em}3},\mathrm{\Phi }_{A\mu }=\mathrm{\hspace{0.17em}0}`$ $`:`$ $`m_{H_2}=400.0\mathrm{GeV},m_{H_3}=400.4\mathrm{GeV},`$
$`\mathrm{tan}\beta =\mathrm{\hspace{0.17em}\hspace{0.17em}3},\mathrm{\Phi }_{A\mu }={\displaystyle \frac{\pi }{2}}`$ $`:`$ $`m_{H_2}=396.6\mathrm{GeV},m_{H_3}=404.5\mathrm{GeV},`$
$`\mathrm{tan}\beta =10,\mathrm{\Phi }_{A\mu }=\mathrm{\hspace{0.17em}0}`$ $`:`$ $`m_{H_2}=397.5\mathrm{GeV},m_{H_3}=400.0\mathrm{GeV},`$
$`\mathrm{tan}\beta =10,\mathrm{\Phi }_{A\mu }={\displaystyle \frac{\pi }{2}}`$ $`:`$ $`m_{H_2}=397.3\mathrm{GeV},m_{H_3}=400.5\mathrm{GeV}.`$ (60)
As a result, a significant overlapping, i.e. interference between two Higgs–boson resonances is expected.
In the previous section, we have listed 16 polarization observables constructed solely by initial muon polarizations under the assumption that the production plane is (at least statistically) reconstructed, and 4 additional polarization observables by combining initial muon polarizations and final equal helicity configurations only for the case when the production plane does not have to be explicitly reconstructed. For some polarization observables it is also important to experimentally identify the electric charges of the produced fermion pair from their decay products. Taking into account all these requirements naturally leads us to the following classification; for the energy range around the heavy Higgs boson resonances the production of a top quark pair is recommended to be used (when the heavy Higgs–boson masses are larger than twice the top–quark mass), and for the energy close to the lightest Higgs mass, the tau lepton or bottom quark pairs are recommended to be considered<sup>§</sup><sup>§</sup>§The $`b\overline{b}`$ decay mode of the lightest Higgs boson, even if it has the largest branching fraction, is not useful for certain observables because of the difficulty in determining the polarization and electric charge of the produced bottom quark and anti–quark experimentally.
### 4.1 Two heavy Higgs bosons for top–pair production
In the MSSM two heavy Higgs bosons are (almost) degenerate with a mass splitting comparable to or less than their widths for $`M_Am_Z`$ as in the case of the parameter set (59), and one of them is CP–even and the other CP–odd in the CP invariant theory. However, the CP–violating neutral Higgs–boson mixing can cause a mass splitting larger than their typical widths and it can lead to several significant non–vanishing CP–odd observables, in particular, in the $`t\overline{t}`$ mode. In this section, we investigate in detail the possibility of probing those aspects through the observables classified in the previous section by considering the cases without and with direct reconstruction of the production plane separately.
#### 4.1.1 Without direct reconstruction of the production plane
If the production plane is not directly reconstructed but only the helicities of the produced top quark and anti-quark are statistically determined, all the interference between the scalar and vector contributions is averaged away over the azimuthal angle $`\mathrm{\Phi }`$ and only the observables, $`\{C_1,C_2,C_3,C_4,C_{15},C_{16};D_1,D_2,D_3,D_4\}`$ can be reconstructed.
Figures 3 and 4 exhibit six independent observables of definite CP and CP$`\stackrel{~}{\mathrm{T}}`$ parities near the region of the heavy Higgs–boson resonances;
$`(\mathrm{a})\sigma _{_{RL}}[++],(\mathrm{b})\sigma _{_{LR}}[++],(\mathrm{c})/(\mathrm{d}){\displaystyle \frac{\sigma _{_{RR}}\pm \sigma _{_{LL}}}{2}}[\pm \pm ],(\mathrm{e})\sigma _{}[++],(\mathrm{f})\sigma _{}[+],`$
for $`\mathrm{tan}\beta =3,10`$ in Figs. 3 and 4, respectively. The observables obtained by controlling solely the muon and anti–muon polarizations are given by
$`\sigma _{_{RL/LR}}`$ $`=`$ $`{\displaystyle \frac{N_C\beta }{16\pi s}}{\displaystyle d\mathrm{cos}\mathrm{\Theta }\left[C_1\pm C_2\right]},`$
$`\sigma _{_{RR/LL}}`$ $`=`$ $`{\displaystyle \frac{N_C\beta }{8\pi s}}\left[C_3\pm C_4\right],`$
$`\sigma _/`$ $`=`$ $`{\displaystyle \frac{N_C\beta }{16\pi s}}(\pm )C_{15/16}.`$ (61)
As shown in the frames (a) and (b) the spin–1 (LR) and (RL) observables $`\sigma _{LR}`$ and $`\sigma _{RL}`$ in both figures are almost constant and the former is almost twice larger than the latter, showing the large left–right asymmetry. On the other hand, the average of the observables $`\sigma _{LL}`$ and $`\sigma _{RR}`$ showing the Higgs–boson contributions are peaked on each heavy Higgs–boson resonance and the size at each pole is comparable to that of the spin–1 observable. In the CP invariant case ($`\mathrm{\Phi }_{A\mu }=0`$) with $`\mathrm{tan}\beta =3`$ the parameter set (59) with $`\mathrm{tan}\beta =3`$ accidentally generates two extremely degenerate Higgs bosons as denoted by a single resonance line in the frame (c) as well as a single oscillating pattern in the frame (e) for $`\sigma _{}`$ of Fig. 3 unlike the case for $`\mathrm{tan}\beta =10`$ with a finite mass splitting of about 3 GeV. The CP–odd observables $`(\sigma _{RR}\sigma _{LL})/2`$ and $`\sigma _{}`$ are identically zero for both $`\mathrm{tan}\beta =3`$ and $`10`$ as they ought to be. On the other hand, in the CP non–invariant case a large mass splitting between two Higgs–boson resonances is developed for $`\mathrm{tan}\beta =3`$, implying a large CP–violating mixing between two Higgs bosons, while the splitting in the case of $`\mathrm{tan}\beta =10`$ is not so much enlarged. A more quantitative understanding about the mass splitting can be obtained from Eq. (60). In this case, the CP–odd observables are non–vanishing. Note that observable $`\sigma _{}`$ are quite sizable on the poles, especially, in the case of $`\mathrm{tan}\beta =3`$.
The observable $`(\sigma _{_{RR}}\sigma _{_{LL}})/2`$ in the frame (d) is CP$`\stackrel{~}{\mathrm{T}}`$–odd, involving the propagator combination $`𝒟_{kl}`$, but the observables $`\sigma _{}`$ and $`\sigma _{}`$ in the frames (e) and (f) are CP$`\stackrel{~}{\mathrm{T}}`$–even, involving the propagator combination $`𝒮_{kl}`$. Examining the analytic structure of those observables and taking into account the parameter set (59), we find that (i) the coefficient related with the CP–odd (CP–even) resonance pole is negative (positive), respectively; (ii) for $`\mathrm{tan}\beta =3`$ the lighter (heavier) resonance of the heavy Higgs bosons is dominantly CP–odd (CP–even) in both the CP invariant and non–invariant cases, respectively; (iii) for $`\mathrm{tan}\beta =10`$ the lighter resonance is dominantly CP–even (CP–odd) and the heavier one CP–odd (CP–even) in the CP invariant (non–invariant) case for $`\mathrm{\Phi }_{A\mu }=0`$ ($`\pi /2`$), respectively. Combining all these features leads to a symmetric resonance pattern in the frame (d) and anti–symmetric patterns in the frames (e) and (f) in both figures and a remarkable sign flipping at each resonance pole in the frame (e) of Fig. 4.
We present in Fig. 5 three observablesThe sum $`\mathrm{\Delta }_{_{RR}}+\mathrm{\Delta }_{_{LL}}`$ is strongly suppressed in the present case due to the orthogonality of the neutral–Higgs mixing matrix $`O`$: $`\mathrm{sin}\beta `$ is very close to unity even for $`\mathrm{tan}\beta =3`$ such that the product sum $`O_{2k}O_{2l}+s_\beta ^2O_{1k}O_{1l}`$ is approximately $`\delta _{kl}`$, i.e. vanishing in the present case because of $`kl`$ forced by the antisymmetric combination $`O_{1k}S_l^fS_k^fO_{1l}`$ with $`S_l^f=O_{3l}`$ for $`f=t`$. related with the difference between the cross sections with the $`(++)`$ and $`()`$ helicity configurations:
$`{\displaystyle \frac{\mathrm{\Delta }_{_{RR}}\mathrm{\Delta }_{_{LL}}}{2}}[++],\mathrm{\Delta }_{}[+],\mathrm{\Delta }_{}[],`$
near the heavy Higgs–boson resonances, the explicit form of which is obtained by integrating the distributions $`\{D_2,D_3,D_4\}`$ over the production angles as
$`\mathrm{\Delta }_{_{RR/LL}}={\displaystyle \frac{N_C\beta }{8\pi s}}\left[D_1\pm D_2\right],\mathrm{\Delta }_/={\displaystyle \frac{N_C\beta }{16\pi s}}D_{3/4}.`$ (62)
It should be noted that every observable $`D_i`$ ($`i=1`$ to 4) involves the antisymmetric combination $`(O_{1k}O_{3l}O_{3k}O_{1l})`$ forcing $`kl`$ ($`k,l=1,2,3`$). Because two heavy Higgs bosons exhibit almost a typical two–state mixing for the parameter set (59), there exist simply one coupling combination formed by the mixing matrix elements for each observable. As a result, the $`\sqrt{s}`$ dependence of each observable is (almost) completely determined by that of the propagator combinations $`𝒮_{23}`$ or $`𝒟_{23}`$, depending on whether the observable is CP$`\stackrel{~}{\mathrm{T}}`$–even or CP$`\stackrel{~}{\mathrm{T}}`$–odd. As two Higgs bosons are extremely degenerate in the CP invariant case with $`\mathrm{tan}\beta =3`$ a single resonance peak (solid line) is shown in the first two upper frames of Fig. 5. However, we note in the upper part of Fig. 5 that (i) in the CP non–invariant case ($`\mathrm{\Phi }_{A\mu }=\pi /2`$) a large mass splitting is developed; and (ii) the CP$`\stackrel{~}{\mathrm{T}}`$–even observable $`(\mathrm{\Delta }_{RR}\mathrm{\Delta }_{LL})/2`$ changes its sign at each resonance pole as $`𝒮_{23}`$ does, while the other two CP$`\stackrel{~}{\mathrm{T}}`$–odd observables have a two–pole structure of the same sign. On the other hand, the mass splitting between two heavy Higgs bosons for $`\mathrm{tan}\beta =10`$ is not so different in the CP invariant and non–invariant cases. As explained in the case of $`\mathrm{tan}\beta =3`$, one can see a sign change at each resonance pole in the CP$`\stackrel{~}{\mathrm{T}}`$–even observable and a typical two–pole structure in the CP$`\stackrel{~}{\mathrm{T}}`$–odd observables. It is also noteworthy that the observable $`\mathrm{\Delta }_{}`$ is negative (positive) for the CP invariant (non–invariant) case. Finally, the typical size of the CP–odd observables is very small compared to that of the CP–even observables.
#### 4.1.2 With direct reconstruction of the production plane
All the terms involving interference between the scalar and vector contributions requires a reasonable identification of the production plane and they are proportional to the mass of the final fermion. The $`\tau ^+\tau ^{}`$ mode with two final neutrinos escaping detection can be used at high energies where the produced tau leptons are very energetic and therefore their decay products fly along the original tau direction to a very good approximation. Nevertheless, the small mass leads to very small interference effects. On the other hand, the production plane in the $`t\overline{t}`$ mode can be determined without big difficulty and the top–quark mass is almost 100 times larger than the tau–lepton mass. In this light, we will concentrate on the $`t\overline{t}`$ mode which could generate significant interference effects.
Among the eight interference terms the four terms $`\{C_7,C_8,C_9,C_{10}\}`$ are CP–odd and the other four terms $`\{C_5,C_6,C_{11},C_{12}\}`$ are CP–even. It is noteworthy that all the CP–odd terms are proportional to $`R_3`$ while all the CP–even observables to $`R_1`$; as a matter of fact the CP properties of the interference terms originates from those of $`R_1`$ and $`R_3`$ as can be worked out by their explicit form for the $`t\overline{t}`$ mode:
$`R_1=2h_\mu h_f{\displaystyle \underset{k=1}{\overset{3}{}}}D_{H_k}O_{2k}O_{3k},R_3=2ih_\mu h_fs_\beta {\displaystyle \underset{k=1}{\overset{3}{}}}D_{H_k}O_{1k}O_{3k}.`$ (63)
We note that (i) $`R_1`$ describes the CP–preserving mixing between two CP–even states, $`\varphi _1`$ and $`\varphi _2`$, so that the relevant observables can be used to select the dominantly CP–even Higgs–boson states; (ii) the coefficients $`O_{22}O_{32}`$ and $`O_{23}O_{33}`$ have an equal sign while the coefficients $`O_{12}O_{32}`$ and $`O_{13}O_{33}`$ have an opposite sign; (iii) $`R_3`$ the CP–violating mixing between $`\varphi _2`$ and the CP–odd state $`a`$. In addition, the coefficients $`W_{1,3}`$ involving vector bilinear charges are (almost) real to a very good approximation so that the SV interference terms depends on the real and imaginary parts of each Higgs–boson propagator $`D_{H_k}`$ itself.
Certainly we need to apply an appropriate $`\mathrm{\Phi }`$–dependent weight function to extract each interference term; for example, ($`\mathrm{cos}\alpha \mathrm{cos}\overline{\alpha }`$) with $`P_T=\overline{P}_T=1`$ taken can be used as a weight function to extract $`C_7`$. The extraction efficiency should be determined with the detailed information on various experimental machine parameters. So, we will not provide any further detailed procedure to extract those observables, but concentrate ourselves on the observables $`C_i`$ themselves. Let us define the observables $`\sigma _i`$ ($`i=5`$ to $`12`$) as:
$`\sigma _i={\displaystyle \frac{N_C\beta }{32\pi s}}{\displaystyle \mathrm{d}\mathrm{cos}\mathrm{\Theta }C_i},`$ (64)
so that $`\sigma _i`$ has the same CP property as $`C_i`$.
Taking the maximal CP–violating phase $`\mathrm{\Phi }_{A\mu }=\pi /2`$, we present in Fig. 6 the CP–odd observables $`\{\sigma _7[],\sigma _8[+],\sigma _9[],\sigma _{10}[+]\}`$ for the $`t\overline{t}`$ mode near the heavy Higgs–boson resonances for $`\mathrm{tan}\beta =3`$ (solid line) and $`\mathrm{tan}\beta =10`$ (dashed line) with the SUSY parameter set (59). The CP$`\stackrel{~}{\mathrm{T}}`$–even observables $`\sigma _8`$ and $`\sigma _{10}`$ changes their sign at each pole as the real part of each Higgs–boson propagator changes its sign. On the other hand, the fact that the coefficients $`O_{12}O_{32}`$ and $`O_{13}O_{33}`$ have an opposite sign is responsible for the sign–flipping pattern at the poles, appearing in the left frames of Fig. 6. On the whole, among those CP–odd observables $`\sigma _7`$ and $`\sigma _{10}`$ are sizable while the other two observables are small in size.
Figure 7 shows the CP–even observables $`\{\sigma _5[++],\sigma _6[+],\sigma _{11}[++],\sigma _{12}[+]\}`$ for the $`t\overline{t}`$ mode near the heavy Higgs–boson resonances for $`\mathrm{\Phi }_{A\mu }=0`$ (solid line) and $`\mathrm{\Phi }_{A\mu }=\pi /2`$ (dashed line) with the SUSY parameter set (59). In this case, we take $`\mathrm{tan}\beta =3`$ causing a larger CP–violating Higgs–boson mixing than $`\mathrm{tan}\beta =10`$. As explained before, the observables dependent on the coefficients $`O_{2k}O_{3k}`$ single the dominantly CP–even states out. As can be seen in every frame of Fig. 7 the heavier Higgs boson of two heavy Higgs bosons is (dominantly) CP–even in the case of $`\mathrm{\Phi }_{A\mu }=0`$ and $`\pi /2`$, respectively. The CP$`\stackrel{~}{\mathrm{T}}`$–even observables $`\sigma _5`$ and $`\sigma _{11}`$ show a sign change at the Higgs–boson resonance pole, while the CP$`\stackrel{~}{\mathrm{T}}`$–odd observables show a typical peak.
### 4.2 The lightest Higgs boson
According to a detailed analysis of Higgs boson decays, the width of the lightest Higgs boson $`H_1`$ is of the order of MeV for intermediate $`\mathrm{tan}\beta `$, which is smaller than a typical muon energy resolution. The lightest Higgs boson decays dominantly into $`b\overline{b}`$ and sub-dominantly into $`\tau ^+\tau ^{}`$ with the (almost) fixed branching fractions and the interference between the Higgs–exchange and the $`\gamma `$ and $`Z`$–boson exchanges is negligible because of the small $`b`$ and $`\tau `$ masses compared to the Higgs–boson mass, even if the muon beam polarization is employed. Furthermore, $`\mathrm{\Gamma }(H_1f_R\overline{f}_R)=\mathrm{\Gamma }(H_1f_L\overline{f}_L)`$ at the tree level so that any further information on the CP property of the lightest Higgs boson will not be obtained by measuring the difference between the final $`(++)`$ and $`()`$ helicity configurationsThe identification of the $`\tau ^+\tau ^{}`$ helicities is, however, quite useful to reduce the continuum backgrounds from the $`\gamma `$ and $`Z`$ exchanges ..
A convolution of the Higgs–boson production cross section with the beam energy distribution tends to reduce the production rate as given by the expression
$`\mathrm{\Sigma }_{H_1eff}{\displaystyle \frac{\pi \mathrm{\Gamma }_{H_1}}{2\sqrt{2\pi }\sigma _E}}\mathrm{\Sigma }_{H_1}(m_{H_1}),`$ (65)
where $`\sigma _E`$ is the muon beam energy resolution and the distribution $`\mathrm{\Sigma }_{H_1}`$ is given by
$`\mathrm{\Sigma }_{H_1}(m_{H_1})=\left(1+P_L\overline{P}_L\right)C_3+\left(P_L+\overline{P}_L\right)C_4+P_T\overline{P}_T\left[\mathrm{cos}\eta C_{15}\mathrm{sin}\eta C_{16}\right].`$ (66)
Every $`C_i`$ is proportional to the second power of $`m_{H_1}/\mathrm{\Gamma }_{H_1}`$ and factored into the production and decay parts on the lightest Higgs–boson resonance pole. So, apart from the factor depending on the beam energy resolution, the CP property of the lightest Higgs boson can be directly investigated through the Higgs–boson resonance production by polarized muon and anti–muon collisions as well as through final fermion spin–spin correlations with high efficiencies. The reduction factor in Eq. (65) clearly shows the importance of having a very good beam energy resolution in order to obtain the spin–0 Higgs–exchange signal events clearly distinguished from the spin–1 $`\gamma /Z`$–exchange background events. Since several works along this line have been already done, we close this section to referring to the relevant recent works.
## 5 Summary and Conclusion
We have performed a systematic investigation of the production of a third–generation fermion–pair in polarized $`\mu ^+\mu ^{}`$ collisions so as to probe the explicit CP violation in the MSSM Higgs sector, induced radiatively by soft trilinear interactions related to squarks of the third generation. We have classified all the observables for probing the CP property of the Higgs bosons constructed by the initial muon beam polarization along with the final fermions of no polarization and of equal helicity, respectively. The polarization observables have turned out to allow for complete determination of the CP property of the Higgs bosons. Furthermore, we have found that the interference between the spin–0 Higgs–boson and spin–1 $`\gamma /Z`$ contributions can provide a powerful and independent means for the determination of the CP property of two heavy Higgs boson in the top–quark pair production with the c.m. energy near the (almost) degenerate heavy Higgs–boson resonances. On the other hand, there is no sizable interference between the lightest Higgs–boson and spin–1 $`\gamma /Z`$ contributions so that the CP property of the lightest Higgs boson can be optimally measured on its pole by using the initial muon beam polarization or the final fermion spin–spin correlations with high efficiencies while keeping a very good beam energy resolution.
In conclusion, the third–generation fermion–pair production in $`\mu ^+\mu ^{}`$ collisions equipped with initial muon beam polarization and final–fermion spin correlations provides a powerful probe of the CP property of the Higgs bosons in the MSSM with explicit CP violation.
## Acknowledgements
S.Y.C wishes to acknowledge financial support of the 1997 Sughak program of the Korea Research Foundation. E.A thanks KIAS for the great hospitality extended to her while this work was being performed.
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# Direct sampling of exponential phase moments of smoothed Wigner functions
## I Introduction
Quantum-state tomography is a powerful tool allowing us to reconstruct the quantum state of a traveling optical mode, provided that many identical copies of the state can be prepared . The idea of homodyne tomography stimulated research in the field of quantum-state reconstruction of other simple quantum-mechanical systems. Recently, reconstructions of the quantum state of a molecular vibrational mode and the motional quantum state of a trapped ion have been reported.
Optical homodyne tomography relies on balanced homodyne detection. The signal field is mixed with a strong coherent local oscillator (LO) at a lossless $`50/50`$ beam splitter. Both the LO and the signal are derived from a common master oscillator to ensure a stable phase difference $`\theta `$ between them. Two photodetectors are placed at the output ports of the beam splitter and the measured photocurrents are subtracted. The resulting signal is proportional to the rotated quadrature of the signal mode $`x_\theta `$. The measurement, which yields the probability distribution $`w(x_\theta ,\theta )`$ of the quadrature $`x_\theta `$, is repeated for many different phase shifts $`\theta `$ from interval $`[0,2\pi ]`$.
The Wigner function of the signal mode can be recovered from the measured statistics $`w(x_\theta ,\theta )`$ by means of inverse Radon transform . Numerical implementation of this inversion is not simple and a filtering algorithm has to be applied to achieve the desired reconstruction. To avoid these complications, it was suggested to directly get quantities of interest from the measured data by averaging appropriate kernels over the distributions $`w(x_\theta ,\theta )`$. This approach proved to be very fruitful, and kernels for the direct sampling of density-matrix elements in the Fock basis $`\rho _{mn}`$ , the moments $`a^ja^k`$ , Fourier coefficients of the canonical phase distribution , and for smoothed Wigner functions have been found. A different approach to the quantum-state reconstruction employs a maximum likelihood estimation . It was demonstrated recently that this technique can be used to estimate photon number distribution and even a whole density matrix . For a review, see .
In recent years, great attention has been devoted to the quantum phase. Canonical phase distribution introduced by London represents a limit of Pegg-Barnett phase formalism . Recently, an approximate measurement of the canonical phase distribution, using the phase-coherent states, has been proposed . One can also construct phase distributions from the phase-space quasidistributions . The phase distribution obtained from the $`Q`$ function (or smoothed $`Q`$ function in the case of imperfect detection) can be directly measured . An operational approach to the quantum phase, based on the description of a given experimental setup, has been proposed by Noh et al. . The relation between canonical and measured phase distributions was discussed in . For a recent review, see .
Canonical phase distribution as well as phase distributions obtained from quasidistributions cannot be directly sampled from the homodyne data. One has to reconstruct the Wigner function or the density matrix and then use the definition of the phase distribution to calculate it . This detour via the Wigner function or the density matrix complicates numerical data processing and increases error in the final result. However, the exponential phase moments (Fourier coefficients) of the canonical phase distribution can be directly sampled with the use of appropriate kernels . Phase-number uncertainty relations can be verified by sampling the first exponential moment of the canonical phase distribution and the photon-number variance . It was also pointed out in that the exponential phase moments of the Wigner function can be directly sampled.
© 2000 The American Physical Society
But we do not have to restrict ourselves to the exponential phase moments of canonical phase distribution or the Wigner function. In this paper, we consider direct sampling of the exponential phase moments of general $`s`$-parametrized phase distributions. We show that it is possible to directly sample the exponential phase moments of any $`s`$-parametrized quasidistribution for $`s<(1\eta )/\eta `$, where $`\eta `$ is the overall detection efficiency. Namely, we find the expressions for the kernels whose average over data recorded in balanced homodyne detection yields the exponential phase moments. We show that a knowledge of these moments as functions of $`s`$ and the photon-number distribution provides complete characteristics of a given quantum state. The phase moments are Fourier coefficients of the phase distributions defined as radial integrals of the $`s`$-parametrized quasidistributions in the polar coordinates. We demonstrate that these phase distributions can be successfully reconstructed from the sampled phase moments.
The paper is organized as follows. In Sec. II the exponential phase moments are introduced and discussed. In Sec. III simple analytical expressions for the sampling kernels are derived and the influence of imperfect detection is addressed. In Sec. IV the results of Monte Carlo simulations are presented. Section V contains conclusions. Some mathematical issues are linked to the Appendix.
## II Exponential phase moments
The quasidistributions related to various $`s`$ orderings of creation and annihilation operators can be expressed in terms of the density matrix $`\rho `$ ,
$$W(\alpha ,s)=\frac{1}{\pi ^2}e^{s|\beta |^2/2}\mathrm{Tr}\left[\rho e^{(a^{}\alpha ^{})\beta (a\alpha )\beta ^{}}\right]d^2\beta ,$$
(1)
where $`a`$,$`a^{}`$ are annihilation and creation operators. One gets the $`P`$ representation for $`s=1`$, the Wigner function for $`s=0`$ and the $`Q`$ function for $`s=1`$. The $`s`$-parametrized quasidistributions are mutually related through the convolution
$`W(q,p,s_2)`$ $`=`$ $`{\displaystyle \frac{1}{\pi (s_1s_2)}}`$ (4)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\mathrm{exp}[{\displaystyle \frac{(qq^{})^2+(pp^{})^2}{s_1s_2}}]`$
$`\times W(q^{},p^{},s_1)dq^{}dp^{},`$
where $`q=(\alpha +\alpha ^{})/\sqrt{2}`$ and $`p=i(\alpha \alpha ^{})/\sqrt{2}`$ are the usual quadratures, and $`s_1>s_2`$ must hold.
It is convenient to introduce polar coordinates $`q=r\mathrm{cos}\theta `$, $`p=r\mathrm{sin}\theta `$. The phase distribution $`P_s(\theta )`$ related to $`s`$-parametrized quasidistribution is defined as
$$P_s(\theta )=_0^{\mathrm{}}W(r,\theta ,s)r𝑑r.$$
(5)
It should be noted that the phase distributions $`P_s(\theta )`$ can be negative for $`s>1`$. Only the phase distributions obtained from the $`Q`$ function (or the smoothed $`Q`$ function) are positively defined for every quantum state. Moreover, for $`s>0`$, the distributions can be highly singular generalized functions. Thus we restrict ourselves to the negative $`s`$ in the following.
The exponential phase moments are defined as
$`\mathrm{\Psi }_l(s)`$ $`=`$ $`\mathrm{exp}(il\theta )_s={\displaystyle P_s(\theta )e^{il\theta }𝑑\theta }`$ (6)
$`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle _0^{\mathrm{}}}W(r,\theta ,s)e^{il\theta }r𝑑r𝑑\theta .`$ (7)
The moments $`\mathrm{\Psi }_l(s)`$ can be simply determined from any quasidistribution $`W(q,p,s_0)`$ provided that $`s_0>s`$. Indeed, inserting the relation (4) into Eq. (7), we have
$$\mathrm{\Psi }_l(s)=_0^{\mathrm{}}_0^{2\pi }F_l\left(\frac{r}{\sqrt{s_0s}}\right)W(r,\theta ,s_0)e^{il\theta }r𝑑r𝑑\theta .$$
(8)
The filtering functions $`F_l(u)`$ are given by
$$F_l(u)=\frac{1}{\pi }_0^{\mathrm{}}_0^{2\pi }e^{il\varphi }e^{u^2\rho ^2+2u\rho \mathrm{cos}\varphi }\rho 𝑑\rho 𝑑\varphi .$$
(9)
Integration over the angle variable $`\varphi `$ yields the modified Bessel function $`I_l(2u\rho )`$. The resulting integral over radial variable $`\rho `$ can be found in the tables of integrals (Ref. , p. 306, Eq. 2.15.5.4) and we have
$$F_l(u)=\sqrt{\pi }\frac{u}{2}\mathrm{exp}\left(\frac{u^2}{2}\right)\left[I_{\frac{|l|1}{2}}\left(\frac{u^2}{2}\right)+I_{\frac{|l|+1}{2}}\left(\frac{u^2}{2}\right)\right].$$
(10)
The first four filtering functions are plotted in Fig. 1. They start from zero and asymptotically reach unity. The interval, where the functions $`F_l(r/\sqrt{s_0s})`$ are significantly lower than $`1`$, increases with decreasing $`s`$. This implies that the absolute values of the phase moments $`\mathrm{\Psi }_l(s)`$ decrease with decreasing $`s`$ because the modulation of the phase distribution $`P_s(\theta )`$ is suppressed by the smoothing convolution (4).
It is remarkable that the functions $`F_l(u)`$ are closely related to the exponential phase moments of the coherent state $`|\xi `$,
$$\mathrm{\Psi }_l(\xi ;s)=F_l\left(\sqrt{\frac{2}{1s}}|\xi |\right)e^{il\psi },\psi =\mathrm{arg}\xi .$$
(11)
To prove this, we notice that the quasidistributions $`W_s(\alpha )`$ of the coherent state $`|\xi `$ are shifted Gaussians,
$$W(\alpha ,s)=\frac{2}{(1s)\pi }\mathrm{exp}\left(\frac{2|\alpha \xi |^2}{1s}\right).$$
(12)
Inserting this into Eq. (7), we immediately obtain Eq. (11).
The filtering functions $`F_l(u)`$ can be expanded in Taylor series,
$$F_l(u)=\underset{n=0}{\overset{\mathrm{}}{}}f_{n,l}u^{2n+|l|},$$
(13)
where
$$f_{n,l}=\frac{|l|}{2}(1)^n\frac{\mathrm{\Gamma }(n+|l|/2)}{n!(n+|l|)!}.$$
(14)
It is convenient to introduce the parameter $`t`$, $`s_0s=1/t^2`$. With the help of the expansion (13), we can rewrite Eq. (8) as
$`\mathrm{\Psi }_l\left(s_0{\displaystyle \frac{1}{t^2}}\right)=`$ (15)
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}f_{n,l}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{2\pi }}(rt)^{2n+|l|}W(r,\theta ,s_0)e^{il\theta }r𝑑r𝑑\theta .`$ (16)
It follows from this formula that $`\mathrm{\Psi }_l(s)`$ are generating functions of the $`s_0`$-ordered moments,
$$r^{2n+|l|}e^{il\theta }_{s_0}=\frac{1}{(2n+|l|)!f_{n,l}}\frac{d^{2n+|l|}}{dt^{2n+|l|}}\mathrm{\Psi }_l\left(s_0\frac{1}{t^2}\right)|_{t=0}.$$
(17)
The limit $`t0`$ should be taken only after the derivative is performed. The generating functions $`\mathrm{\Psi }_l(s)`$ can be used to determine the moments $`r^{2n+|l|}e^{il\theta }_{s_0}`$ for any ordering parameter $`s_0`$. Notice, however, that the formula (17) fails for $`l=0`$. The exponential phase moments do not allow us to determine the moments $`r^{2n}`$ which are related to photon-number statistics. As an example, consider the Fock state $`|n`$. This state is phase insensitive, $`\mathrm{\Psi }_l(s)=0`$ for $`l0`$, and the phase is uniformly distributed over the $`2\pi `$ interval, $`P_s(\theta )=1/2\pi `$. Note also that the photon-number distribution $`p(n)`$ can be recovered from the phase-averaged quadrature distributions .
The $`s`$-ordered moments (17) are simply related to more familiar moments of creation and annihilation operators. With the help of $`\alpha =2^{1/2}r\mathrm{exp}(i\theta )`$ we find that
$$a^na^{n+l}_s=2^{(n+l/2)}r^{2n+l}e^{il\theta }_s$$
(18)
and a similar expression holds for $`a^{n+l}a^n_s`$. The formula (17) allows us to find any moments $`a^ma^n`$ provided that $`mn`$. Complementarily, the moments $`a^ka^k=:n^k:`$ can be determined from the photon-number distribution.
The phase moments $`\mathrm{\Psi }_l(s)`$ are linear combinations of density-matrix elements $`\rho _{n+l,n}`$,
$$\mathrm{\Psi }_l(s)=\underset{n=0}{\overset{\mathrm{}}{}}c_{n,l}(s)\rho _{n+l,n},$$
(19)
where
$`c_{n,l}(s)`$ $`=`$ $`\left({\displaystyle \frac{2}{1s}}\right)^{n+l/2}[n!(n+l)!]^{1/2}`$ (21)
$`\times {\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{\mathrm{\Gamma }(nk+l/2+1)}{k!(nk)!(n+lk)!}}({\displaystyle \frac{1+s}{2}})^k.`$
If $`l0`$, the relation (19) can be inverted and $`\rho _{n+l,n}`$ can be found from $`\mathrm{\Psi }_l(s)`$. In principle, the knowledge of $`\mathrm{\Psi }_l(s)`$ at an infinite but countable number of points $`s_j`$ can be sufficient for determination of all $`\rho _{n+l,n}`$ from Eq. (19). Diagonal matrix elements appear only in
$$\mathrm{\Psi }_0(s)=\underset{n=0}{\overset{\mathrm{}}{}}\rho _{nn}\mathrm{Tr}\rho =1,$$
(22)
and this relation cannot be inverted. Only when we know both the phase moments $`\mathrm{\Psi }_l(s)`$ and the photon-number distribution $`p(n)=\rho _{nn}`$ can we determine all density-matrix elements $`\rho _{mn}`$ or, equivalently, all moments $`a^na^m_s`$. Thus the simultaneous knowledge of the functions $`\mathrm{\Psi }_l(s)`$ and $`p(n)`$ provides complete information on the quantum state and it is equivalent to the knowledge of the Wigner function or the density matrix.
## III Sampling kernels for the exponential phase moments
Balanced homodyne detection provides statistics $`w(x_\theta ,\theta )`$ of rotated quadratures,
$$x_\theta =\frac{1}{\sqrt{2}}\left(ae^{i\theta }+a^{}e^{i\theta }\right),$$
(23)
where $`\theta `$ is the relative phase between the LO and the signal mode. The probability distribution $`w(x_\theta ,\theta )`$ can be obtained from the Wigner function $`W(q,p)`$ as a marginal distribution ,
$`w(x_\theta ,\theta )`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\delta (x_\theta q\mathrm{cos}\theta p\mathrm{sin}\theta )`$ (25)
$`\times W(q,p)dqdp.`$
We would like to sample the moments $`\mathrm{\Psi }_l(s)`$ directly from the homodyne data $`w(x_\theta ,\theta )`$ with the use of the kernels $`𝒦_l(x_\theta ,\theta ;s)`$:
$`\mathrm{\Psi }_l(s)={\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _0^{2\pi }}𝒦_l(x_\theta ,\theta ;s)w(x_\theta ,\theta )𝑑x_\theta 𝑑\theta .`$ (26)
The $`\theta `$ dependence of the kernels must be of the form $`\mathrm{exp}(il\theta )`$ . Thus we look for the kernels in the form
$`𝒦_l(x_\theta ,\theta ;s)=K_l(x_\theta ,s)e^{il\theta }.`$ (27)
In what follows we restrict ourselves to positive $`l`$. For negative $`l`$, the exponential moments can be obtained by complex conjugation, $`\mathrm{\Psi }_l(s)=\mathrm{\Psi }_l^{}(s)`$. Now we substitute Eq. (25) into Eq. (26), perform integration over $`x_\theta `$, and rewrite the remaining integral in polar coordinates. After some algebra, we arrive at
$`\mathrm{\Psi }_l(s)`$ $`=`$ $`{\displaystyle _0^{2\pi }}{\displaystyle _0^{\mathrm{}}}\left[{\displaystyle _0^{2\pi }}K_l(r\mathrm{cos}\varphi ,s)e^{il\varphi }𝑑\varphi \right]`$ (29)
$`\times W(r,\theta )e^{il\theta }rdrd\theta .`$
Comparing the formulas (29) and (8), where we set $`s_0=0`$, we conclude that the kernel $`K_l(x_\theta ,s)`$ must fulfill the integral equation
$$_0^{2\pi }K_l(r\mathrm{cos}\theta ,s)e^{il\theta }𝑑\theta =F_l\left(r/|s|^{1/2}\right).$$
(30)
In order to solve this equation, we expand the kernel $`K_l(x_\theta ,s)`$ in Taylor series,
$$K_l(x_\theta ,s)=\underset{n=0}{\overset{\mathrm{}}{}}a_n(l,s)x_\theta ^n.$$
(31)
This expansion is inserted into Eq. (30) and the integration over $`\theta `$ is carried out, using the formula
$`{\displaystyle _0^{2\pi }}(\mathrm{cos}\theta )^{2n+l}e^{il\theta }𝑑\theta ={\displaystyle \frac{2\pi }{2^{2n+l}}}\left({\displaystyle \genfrac{}{}{0pt}{}{2n+l}{n}}\right).`$ (32)
Comparing the Taylor series on the left-hand side of Eq. (30) with the series (13), we find the coefficients $`a_n(l,s)`$. Inserting them back into the series (31), we arrive at
$`K_l(x_\theta ,s)={\displaystyle \frac{l}{4\pi }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{\mathrm{\Gamma }(n+l/2)}{(2n+l)!}}\left({\displaystyle \frac{2x_\theta }{|s|^{1/2}}}\right)^{2n+l}.`$ (33)
Notice that the kernel is a function of a specific combination of $`x_\theta `$ and $`s`$, $`u=x_\theta /|s|^{1/2}`$. In the following we use $`u`$ for simplicity. Let us discuss the relation between the kernels $`K_l(u)`$ and $`K_{l+2}(u)`$. We have
$`K_{l+2}(u)`$ $`=`$ $`{\displaystyle \frac{l+2}{4\pi }}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{\mathrm{\Gamma }(n+l/2)}{(2n+l)!}}\left(2u\right)^{2n+l}`$ (34)
$`=`$ $`{\displaystyle \frac{l+2}{l}}K_l(u)+{\displaystyle \frac{l+2}{4\pi }}{\displaystyle \frac{\mathrm{\Gamma }(l/2)}{l!}}\left(2u\right)^l.`$ (35)
However, the kernels $`K_l`$ are not uniquely determined. Any polynomial of order lower than $`l`$ can be added to kernel $`K_l`$, because all such polynomials are solutions of the homogeneous integral equation
$$_0^{2\pi }f(r\mathrm{cos}\theta )e^{il\theta }=0.$$
(36)
Thus we can neglect the last term in the formula (35) and we can define the kernels for which
$`K_{l+2}(u)={\displaystyle \frac{l+2}{l}}K_l(u)`$ (37)
holds. It remains to find out the kernels $`K_1`$ and $`K_2`$. The summation of the series can be found in the Appendix. The results are
$`K_1(u)={\displaystyle \frac{1}{4}}\mathrm{erf}\left(u\right),`$ (38)
$`K_2(u)={\displaystyle \frac{1}{\sqrt{\pi }}}{\displaystyle _0^u}e^{y^2}\mathrm{erfi}(y)𝑑y.`$ (39)
The kernels are plotted in Fig. 2 and Fig. 3, respectively.
Combining this result with the recurrence formula (37), we finally have
$`𝒦_{2l+1}(x_\theta ,\theta ;s)`$ $`=`$ $`(1)^l(2l+1)K_1\left(x_\theta /|s|^{\frac{1}{2}}\right)e^{i(2l+1)\theta },`$ (40)
$`𝒦_{2l}(x_\theta ,\theta ;s)`$ $`=`$ $`(1)^{l1}lK_2\left(x_\theta /|s|^{\frac{1}{2}}\right)e^{i2l\theta }.`$ (41)
For large $`x_\theta `$, all the kernels tend to the same limit because we move to the strong classical field domain and the differences between various $`s`$ orderings vanish. The limit for odd kernels is straightforward. We simply notice that
$$\underset{x\pm \mathrm{}}{lim}\mathrm{erf}(x)=\pm 1.$$
(42)
The limit for even kernels can be found if we take into account that for large $`x`$,
$$e^{x^2}\mathrm{erfi}(x)\frac{1}{\sqrt{\pi }}\frac{1}{x}.$$
(43)
Inserting this into Eq. (41) we have for large $`x_\theta `$
$$𝒦_{2l}(x_\theta ,\theta ;s)\frac{1}{\pi }l(1)^{l1}\mathrm{ln}|x_\theta |e^{i2l\theta }+C_{l,s}e^{i2l\theta }.$$
(44)
Here $`C_{l,s}`$ is some constant. The superfluous term containing this constant can be omitted for reasons discussed above and we can see that as a limit all kernels approach those for the phase moments of the Wigner function :
$`𝒦_{2l+1}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(2l+1)(1)^l\mathrm{sgn}(x_\theta )e^{i(2l+1)\theta },`$ (45)
$`𝒦_{2l}`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}l(1)^{l1}\mathrm{ln}|x_\theta |e^{i2l\theta }.`$ (46)
Up to now, we have considered ideal detectors having unit quantum efficiency. In a realistic experiment, the detection efficiency $`\eta `$ is lower than $`100\%`$ and the smoothed quadrature distributions $`w(x_\theta ,\theta ;\eta )`$ are recorded ,
$`w(x_\theta ,\theta ;\eta )`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{\pi (1\eta )}}}`$ (48)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}w(x_\theta ^{},\theta )\mathrm{exp}[{\displaystyle \frac{(x_\theta \sqrt{\eta }x_\theta ^{})^2}{1\eta }}]dx_\theta ^{}.`$
The smoothed quadrature distributions $`w(x_\theta ,\theta ;\eta )`$ can be obtained from the scaled and smoothed Wigner function,
$`w(x_\theta ,\theta ;\eta )`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}\delta (x_\theta q\mathrm{cos}\theta p\mathrm{sin}\theta )`$ (51)
$`\times {\displaystyle \frac{1}{\eta }}W({\displaystyle \frac{q}{\sqrt{\eta }}},{\displaystyle \frac{p}{\sqrt{\eta }}},{\displaystyle \frac{1\eta }{\eta }})dqdp.`$
The scaling and smoothing are two factors which must be included in the kernels $`𝒦_l(x_\theta ,\theta ;s,\eta )`$. The scaling means that we must replace $`x_\theta `$ by $`x_\theta /\sqrt{\eta }`$. The smoothing tells us that the kernels $`𝒦(x_\theta ,\theta ,s)`$ would provide us with exponential phase moments $`\mathrm{\Psi }_l(s+s_\eta )`$, $`s_\eta =(1\eta )/\eta `$. Thus we must replace $`s`$ with $`ss_\eta `$ in all the expressions (41). It is obvious that the losses impose a new limit. We can reconstruct only exponential phase moments for the phase distributions corresponding to $`s<s_\eta `$. The modified kernels are
$$𝒦_l(x_\theta ,\theta ;s,\eta )=𝒦_l(x_\theta /\sqrt{\eta },\theta ;s+(1\eta )/\eta ),$$
(52)
and the condition $`s<(1\eta )/\eta `$ must be fulfilled.
## IV Monte Carlo simulations
In order to test the kernels, we performed Monte Carlo simulations of the homodyne detection and we present here the results of simulations for the squeezed vacuum state $`|\zeta `$,
$$|\zeta =\mathrm{exp}\left(\frac{1}{2}\zeta a^2\frac{1}{2}\zeta ^{}a^2\right)|0,$$
(53)
where $`|0`$ is the vacuum state. The squeezed vacuum state belongs to the class of Gaussian states, i.e. states whose quasidistributions $`W(q,p,s)`$ have Gaussian form. The phase distribution $`P_s(\theta )`$ for the general Gaussian mixed state was determined in . In particular, it holds that $`P_s(\theta )`$ of the squeezed vacuum state can be expressed as
$$P_s(\theta )=\frac{1}{2\pi }\frac{(B_s^2C^2)^{1/2}}{B_sC\mathrm{cos}(2\theta \psi )},$$
(54)
where
$`B_s`$ $`=`$ $`\mathrm{sinh}^2|\zeta |+(1s)/2,`$ (55)
$`C`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{sinh}(2|\zeta |),`$ (56)
and $`\psi =\mathrm{arg}\zeta `$. The phase moments can be calculated with the help of the residue theorem. One arrives at
$`\mathrm{\Psi }_{2l}(s)`$ $`=`$ $`\left(B_s/C\sqrt{B_s^2/C^21}\right)^le^{il\psi },`$ (57)
$`\mathrm{\Psi }_{2l1}(s)`$ $`=`$ $`0.`$ (58)
In our simulations, the sampling was performed for 120 values of $`\theta `$ equidistantly placed at the interval $`[0,2\pi ]`$ and 5000 samples have been made for each $`\theta `$. We assumed that the overall detection efficiency is $`\eta =80\%`$ and we used the loss compensating kernels (52).
Figure 4 shows the reconstructed phase moments of the $`Q`$ function, $`\mathrm{\Psi }_l(1)`$. The results are in very good agreement with the exact values following from Eq. (58). Statistical errors were calculated in a manner described in . As a rule, error increases with increasing $`l`$ and this uncertainty is responsible for the fast oscillations in the reconstructed probability distribution $`P_1(\theta )`$, see Fig. 5.
The reconstructed moments $`\mathrm{\Psi }_l(s)`$, considered as functions of the ordering parameter $`s`$, are plotted in Fig. 6. Again, we found that the curves are in good agreement with their theoretical counterparts. Notice that, due to the assumed $`80\%`$ efficiency of the detection, we were able to sample only moments for $`s<0.25`$.
We repeated our simulations also for other types of quantum states such as coherent states and displaced Fock states. In all cases, the reconstruction procedure worked well and the sampled moments were in good agreement with the theoretical values. We emphasize that we have used only $`6\times 10^5`$ samples in our simulations and such an amount of data can be routinely recorded in the experiment.
## V Conclusions
We have shown that the exponential phase moments of the $`s`$-parametrized quasidistributions are generating functions of the moments of creation and annihilation operators. A simultaneous knowledge of photon-number distribution and the functions $`\mathrm{\Psi }_l(s)`$ provides a complete description of the quantum state. We have found kernels for direct sampling of the moments $`\mathrm{\Psi }_l(s)`$ from quadrature distributions measured in optical homodyne detection. The detection efficiency $`\eta `$ imposes a bound on the ordering parameter, we can sample only phase moments for $`s<(1\eta )/\eta `$. In the ideal case $`\eta =1`$ and the Wigner function represents the limit; for $`\eta =0.5`$ the limit is formed by a $`Q`$ function. We performed numerical Monte Carlo simulations of homodyne detection, thereby demonstrating the feasibility of direct sampling of the exponential phase moments from experimental data.
###### Acknowledgements.
I would like to thank T. Opatrný, J. Peřina, and I.Sh. Averbukh for helpful discussions. I am pleased to acknowledge support of the U.S.-Israel Binational Science Foundation (Grant No. 96–00432).
## Summation of the series for the kernels $`K_1`$ and $`K_2`$
Here we sum the Taylor series for kernels $`K_1(u)`$ and $`K_2(u)`$. We start with $`K_1(u)`$. Using the formula for the Gamma function of a half-integer,
$$\mathrm{\Gamma }(n+1/2)=\sqrt{\pi }\frac{(2n)!}{2^{2n}n!},$$
(59)
the series for $`K_1(u)`$ take on the form
$$K_1(u)=\frac{1}{4\pi }\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\frac{2\sqrt{\pi }}{n!}\frac{u^{2n+1}}{(2n+1)}.$$
(60)
The derivative of the kernel is
$$\frac{d}{du}K_1(u)=\frac{1}{4}\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\frac{2}{\sqrt{\pi }}\frac{u^{2n}}{n!}=\frac{1}{4}\frac{2}{\sqrt{\pi }}e^{u^2}.$$
(61)
Integrating the above equation we arrive at
$$K_1(u)=\frac{1}{4}\mathrm{erf}(u).$$
(62)
We adopt a similar approach to determine $`K_2(u)`$,
$$K_2(u)=\frac{1}{2\pi }\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\frac{n!}{(2n+2)!}(2u)^{2n+2}.$$
(63)
We calculate the derivatives
$$f(u)=\frac{d}{du}K_2(u)=\frac{2}{2\pi }\underset{n=0}{\overset{\mathrm{}}{}}(1)^n\frac{n!}{(2n+1)!}(2u)^{2n+1}$$
(64)
and
$`f^{}(u)`$ $`=`$ $`{\displaystyle \frac{d^2}{du^2}}K_2(u)={\displaystyle \frac{4}{2\pi }}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{n!}{(2n)!}}(2u)^{2n}`$ (65)
$`=`$ $`{\displaystyle \frac{2}{2\pi }}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1)^n{\displaystyle \frac{(n1)!}{(2n1)!}}(2u)^{2n}+{\displaystyle \frac{2}{\pi }}.`$ (66)
Thus we have
$$f^{}(u)=2uf(u)+\frac{2}{\pi }.$$
(67)
Let us look for the function $`f(u)`$ in the form
$$f(u)=\frac{2}{\pi }\frac{g(u)}{g^{}(u)}.$$
(68)
Substituting this into the above equation, we finish with
$`g^{\prime \prime }(u)`$ $`=`$ $`2ug^{}(u),`$ (69)
$`g(u)`$ $`=`$ $`{\displaystyle \frac{\sqrt{\pi }}{2}}\mathrm{erfi}(u).`$ (70)
Inserting this into Eq. (68), we conclude that
$$K_2(u)=\frac{1}{\sqrt{\pi }}_0^ue^{y^2}\mathrm{erfi}(y)𝑑y.$$
(71)
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# Detection of interstellar CH₃ 1footnote 11footnote 1Based on observations made with ISO, a project of ESA with participation of ISAS and NASA, and the SWS, a joint project of SRON and MPE with contributions from KU Leuven, Steward Observatory and Phillips Laboratory.
## 1 Introduction
The methyl radical $`\mathrm{CH}_3`$ is an important intermediate product in the basic ion-molecule gas-phase chemistry networks in the interstellar medium driven by cosmic-ray ionization (Herbst & Klemperer 1973). Together with $`\mathrm{CH}`$ and $`\mathrm{CH}_2`$, it is produced by a series of reactions starting with C + H$`{}_{}{}^{+}{}_{3}{}^{}`$ $``$ CH<sup>+</sup> \+ H or the radiative association of C<sup>+</sup> \+ H<sub>2</sub> $``$ CH$`{}_{}{}^{+}{}_{2}{}^{}`$ \+ h$`\nu `$, followed by hydrogen abstraction reactions and dissociative recombination. Alternatively, it can be produced by photodissociation of methane (CH<sub>4</sub>). Subsequent reactions of C<sup>+</sup> with CH<sub>3</sub> form one of the most important steps in the formation of more complex hydrocarbons.
The ion $`\mathrm{H}_3^+`$ which initiates the chemistry in the cold gas has been detected only recently toward the Galactic center by Geballe et al. (1999) at a surprisingly high abundance. This unique line of sight turns out to be an extremely valuable environment to study abundances in the cold low density interstellar medium, since even minor species like $`\mathrm{CH}_3`$ may be detectable due to its long path.
Although $`\mathrm{CH}_3`$ is a simple species, it is difficult to obtain accurate laboratory measurements of its molecular parameters since, as a radical, it recombines very fast with other particles in a gas. Herzberg (1961) and Herzberg & Shoosmith (1956) were the first to determine that the molecule is planar, but definite proof came only from measurements of the out-of-plane bending mode $`\nu _2`$ at 16 $`\mu `$m by Yamada et al. (1981).
Observations of the 16–16.5 $`\mu `$m wavelength range are strongly hampered from the ground due to the Earth’s atmosphere. The first detections of CH<sub>3</sub> in space have become possible only using the Infrared Space Observatory (ISO) (Kessler et al. 1996). Bézard et al. (1998, 1999) have recently detected CH<sub>3</sub> in the atmospheres of Saturn and Neptune respectively, but no previous searches for the molecule in interstellar space have been reported.
## 2 $`\mathrm{CH}_3`$ spectroscopy
Since the CH<sub>3</sub> radical is planar and symmetric, it does not have electric dipole allowed rotational lines which could be detected in the (sub-)millimeter wavelength range. The planar nature also implies that the symmetric stretch $`\nu _1`$ is infrared inactive and the asymmetric stretch $`\nu _3`$ at 3.16 $`\mu `$m relatively weak. Indeed, the transition dipole moment of the $`\nu _3`$ band is found to be a factor of three weaker than the out-of-plane bending mode $`\nu _2`$ (Triggs et al. 1992, Amano et al. 1992).
To calculate the $`\nu _2`$ spectrum, the term energies were taken from Yamada et al. (1981). The nuclear spin of the H-atoms can couple either to a quartet or doublet state, with nuclear spin statistical weights of 4 and 2, respectively. Because CH<sub>3</sub> follows Fermi-Dirac statistics, the $`K=3,6,9,\mathrm{}`$ levels are quartet states, and the other K-values doublet states. The strongest $`Q`$-branch has $`N=K`$ and is located at 16.5 $`\mu `$m; the strongest other feature is the $`{}_{}{}^{Q}R(0)`$ line at 16.0 $`\mu `$m thanks to its favorable Hönl-London factor. The band strength of $`(2.5\pm 0.8)\times 10^{17}`$ cm<sup>-1</sup> (molecule cm<sup>-2</sup>)<sup>-1</sup> was taken from Wormhoudt & McCurdy (1989). The calculation of the spectrum was performed as described in Helmich (1996).
The shape of the spectrum is very sensitive to the excitation temperature (see Figure 7.14 of Helmich 1996). Besides the strong $`{}_{}{}^{Q}Q`$-branch and $`{}_{}{}^{Q}R(0)`$ lines, many more features become visible at excitation temperatures above 25 K, most notable the satellite $`{}_{}{}^{R}Q`$-branch at 16.53 $`\mu `$m and the $`P`$(2) line at 17.60 $`\mu `$m.
## 3 Observations and Data Reduction
Observations were carried out in the SWS grating mode AOT06 (de Graauw et al. 1996) at a spectral resolving power $`R=\lambda /\mathrm{\Delta }\lambda 15002200`$. The spectral range covering the $`\mathrm{CH}_3`$ $`Q`$-branch and the $`P`$(2) line has been measured on 1997 February 21 15:24:27-19:03:01 UT, whereas that covering the $`R`$(0) and $`R`$(1) lines has been obtained on the same day at 19:03:45-20:55:41 UT. The SWS aperture size was $`14^{\prime \prime }\times 27^{\prime \prime }`$ and has been centered on the position of Sgr A RA $`17^h45^m40^s.0`$, Dec $`29^{}00^{}28^{\prime \prime }.6`$ (J2000 coordinates), with the long side of the slit oriented within 1 degree of the north-south direction. Due to the rather large aperture size, the Galactic center sources IRS 1, 2, 3 and 7 also fall inside the beam (see, e.g., Geballe et al. 1989), whereas IRS 5 and 6 are positioned just outside the slit.
Data were processed within the SWS interactive analysis system, based on the standard ISO pipeline OLP V8.7 products. The data reduction adhered to the recommendations of Salama et al. (1997). Raw data were rebinned to $`R=5000`$, a value significantly larger than the actual spectral resolving power of the SWS to avoid losing spectral detail when convolving the observed data samples by the bin.
The absolute calibration of the SWS data has about $`\pm 20\%`$ uncertainty on average longwards of $``$15 $`\mu `$m (Salama et al. 1997). However, since our analysis is entirely based on spectra where the continuum is divided out, the actual uncertainty in the results is determined by the noise in the data rather than the actual calibration uncertainty. The main limitation of our analysis originates from the $`\pm 30\%`$ uncertainty in the $`\mathrm{CH}_3`$ $`\nu _2`$ band strength (Wormhoudt & McCurdy 1989, Yamada & Hirota 1983).
## 4 CH<sub>3</sub> results
As shown in Fig. 1, the $`Q`$-branch at 16.5 $`\mu `$m and the $`R`$(0) line at 16.0 $`\mu `$m are clearly detected. This represents the first unambiguous detection of $`\mathrm{CH}_3`$ in the interstellar medium. The upper limit on the $`P`$(2) line at 17.60 $`\mu `$m provides an important constraint on the temperature. Due to a blend with the \[Ne III\] 15.555 $`\mu `$m atomic fine structure line at the SWS grating resolution, no information from the $`R`$(1) line at 15.54 $`\mu `$m could be obtained.
Both the $`Q`$-branch and the $`R`$(0) lines are shifted with respect to their expected LSR wavelengths by about $`20`$ km s<sup>-1</sup>. Although such a shift is close to the SWS wavelength calibration accuracy (Valentijn et al. 1996), a $`V_{\mathrm{LSR}}=30`$ km s<sup>-1</sup> component of cold molecular gas has been reported previously by several authors from observations at radio and millimeter wavelengths with similar beam sizes (Serabyn et al. 1986 and Sutton et al. 1990, $`\mathrm{CO}`$; Pauls et al. 1996, $`\mathrm{H}_2\mathrm{CO}`$; Serabyn & Güsten 1986, $`\mathrm{NH}_3`$; Marr et al. 1992, $`\mathrm{HCO}^+`$; Güsten et al. 1987, $`\mathrm{HCN}`$; Bolton et al. 1964, $`\mathrm{OH}`$). In all cases several velocity components at $`50,30`$ and 0 km s<sup>-1</sup> have been observed at much higher spectral resolutions. The relative strengths of these three components vary between the observed species with the 0 km s<sup>-1</sup> component often the largest. At the SWS spectral resolution of $`150`$ km s<sup>-1</sup> it is not possible to distinguish between these different velocity components, but our observed shift is consistent with a mix of them. The location of the absorbing gas can therefore not be attributed to one particular feature, but is possibly spread along the line of sight toward Sgr A among spiral arms and molecular clouds.
Fits of synthetic spectra to the data, as described in §2 and matching the SWS resolution, were performed for different excitation temperatures (10 to 50 K), different Doppler parameters ($`b`$ between 1.5 and 30 km s<sup>-1</sup>) and column densities. The best fit to the individual absorption depths (Fig. 1 a,b) as well as their ratio (Fig. 1 c) is given in Table 1. The absorption depths are almost independent of Doppler parameter, and are mainly a function of column density. The ratio of the depths is a strong function of the excitation temperature. The inferred low excitation temperature of $`(17\pm 2)`$ K from the 16.0/16.5 $`\mu `$m ratio of (0.8$`\pm `$0.15) is consistent with the non-detection of the P(2) line. Because $`\mathrm{CH}_3`$ has no dipole moment, the populations of the lowest rotational levels are controlled by collisions, so that the excitation temperature is close to the kinetic temperature.
The H<sub>2</sub> column density along the line of sight has been constrained by several sets of observations. First, the measured extinction of 31 mag (Rieke et al. 1989) implies $`N`$(H<sub>2</sub>)$`2\times 10^{22}`$ cm<sup>-2</sup> using $`N_H/A_V=1.9\times 10^{21}`$ cm<sup>-2</sup> mag<sup>-1</sup> and assuming that at least half of the hydrogen is in molecular form. Second, the detection of at least one optically thin C<sup>18</sup>O line (R(0) at 4.7716 $`\mu `$m) in our ISO-SWS observation together with the measured CO excitation temperature of 8-13 K implies $`N`$(C<sup>18</sup>O)=$`(2\pm 0.5)\times 10^{16}`$ cm<sup>-2</sup>. This is in good agreement with the analysis by Moneti & Cernicharo (2000) on the same data. Using <sup>16</sup>O/<sup>18</sup>O= 300 (Wilson & Rood 1994) and CO/H<sub>2</sub>=$`10^4`$ implies $`N`$(H<sub>2</sub>)$`6\times 10^{22}`$ cm<sup>-2</sup>. We adopt $`N`$(H<sub>2</sub>)=$`(6\pm 3)\times 10^{22}`$ cm<sup>-2</sup>, leading to a CH<sub>3</sub> abundance with respect to H<sub>2</sub> of $`x`$(CH<sub>3</sub>)=$`(1.3\genfrac{}{}{0pt}{}{+2.2}{0.7})\times 10^8`$. Note that all abundances would be increased by a factor of 3 if the lower H<sub>2</sub> column density derived from the extinction is used.
## 5 Related species: CH<sub>4</sub>, C<sub>2</sub>H<sub>2</sub> and CH
The availability of the full SWS scan allows searches for other chemically related molecules. Specifically, the $`\nu _2/\nu _4`$ dyad of CH<sub>4</sub> occurs around 7.7 $`\mu `$m and has been observed with the ISO-SWS toward massive protostars by Boogert et al. (1998). Toward the Galactic center, however, gas-phase CH<sub>4</sub> is not detected. Adopting the same excitation temperature as found for $`\mathrm{CH}_3`$, an upper limit for its abundance of $``$1$`\times 10^{15}`$ cm<sup>-2</sup> is found. Solid $`\mathrm{CH}_4`$ is clearly detected by Chiar et al. (2000) toward Sgr A with a column density of ($`3.0\pm 0.7`$)$`\times 10^{16}`$ cm<sup>-2</sup>. Thus, most of the CH<sub>4</sub> is in solid form, consistent with the low temperature.
Detection of a blend of the pure rotational lines of $`\mathrm{CH}`$ at 149.09 and 149.39 $`\mu `$m towards the Galactic center has been reported by White et al. (1999). We have re-analyzed the LWS observations carried out on 1998 February 20 10:11:34-11:06:44 in the LWS01 grating mode (Clegg et al. 1996) at a resolution of $``$ 1500 km s<sup>-1</sup> (Fig. 2 a). The LWS data reduction has been based on OLP V8.7 products and has been carried out within ISAP (Sturm et al. 1997). Outliers have been removed by iterative sigma clipping and the different scans have been flatfielded to their mean value by applying a second order polynomial offset to each individual scan. The fringing in the LWS data, present in all extended source observations has been removed by the dedicated module within ISAP. The inferred equivalent width for the unresolved doublet is $`(85\pm 5)`$ km s<sup>-1</sup>. Using the formulae from Stacey et al. (1987) and assuming $`T_{\mathrm{ex}}`$=17 K, this leads to $`N`$(CH)$`(1.1\pm 0.1)\times 10^{15}`$ cm<sup>-2</sup>.
Finally, the $`\nu _5`$ band of gas-phase C<sub>2</sub>H<sub>2</sub> at 13.7 $`\mu `$m is clearly detected (Fig. 2 b). Following the analysis of Lahuis & van Dishoeck (2000), we find $`N`$(C<sub>2</sub>H<sub>2</sub>)=($`5.5\pm 0.8)\times 10^{14}`$ cm<sup>-2</sup> assuming $`T_{\mathrm{ex}}=17`$ K. Table 1 summarizes the results obtained from ISO observations. Note that the relative abundances of the molecules have much smaller error bars than the absolute values since the uncertainty in the $`\mathrm{H}_2`$ column density cancels out.
## 6 Chemistry
The absolute and relative abundances of the observed molecules have been compared with a wide variety of models, including time- and depth-dependent models. None of the published pure gas-phase dense cloud models can reproduce the observations of all species (e.g., Millar et al. 1997, Lee et al. 1996, see Table 2). In general, the model CH<sub>3</sub> abundances are too low and the CH<sub>4</sub> abundances too high. Also, the model abundance of C<sub>2</sub>H<sub>2</sub> is significantly smaller than that of CH<sub>4</sub>, in contrast with observations. The only models which come close to matching the absolute and relative abundances of CH<sub>3</sub> and CH<sub>3</sub>/CH and CH<sub>3</sub>/CH<sub>4</sub> are low-density translucent cloud models with $`n`$(H<sub>2</sub>)$`10^3`$ cm<sup>-3</sup> and $`A_V`$few mag, so that photodissociation of CH<sub>4</sub> to CH<sub>3</sub> and CH<sub>2</sub> is still efficient. Table 2 lists recent model calculations by Terzieva & Herbst (priv. communication) and results based on the models by van Dishoeck & Black (1986) and Jansen et al. (1995) using updated branching ratios for the dissociative recombination of the hydrocarbon ions (Andersen et al. 2000). Low metal abundances are favored, to prevent destruction of the hydrocarbons by oxygen and by sulfur atoms and ions. Note, however, that even though the abundances at $`A_V3`$ may match the data within a factor of a few, the CH<sub>3</sub> column density in such models is only $`1\times 10^{13}`$ cm<sup>-2</sup>, nearly two orders of magnitude below observations. At the same time, the CH column density of $`1.4\times 10^{14}`$ cm<sup>-2</sup> is a factor of 10 below observations. Because of the strong depth dependence of the CH, CH<sub>3</sub> and CH<sub>4</sub> abundances, it is not possible to reproduce the column density ratios with these same models. The large observed $`\mathrm{H}_3^+`$ column density suggests that there are several clouds along the line of sight. Some combination of low-density diffuse clouds to produce the CH and denser clouds to account for the solid CH<sub>4</sub> may explain those data, but the mix would have to be tailored very specifically to simultaneously approach the large column densities of CH<sub>3</sub> and C<sub>2</sub>H<sub>2</sub>.
An alternative suggestion is to invoke turbulent chemistry, in which a high $`\mathrm{CH}^+`$ abundance leads to enhancements of other hydrocarbons by 1-2 orders of magnitude (e.g. Hogerheijde et al. 1995, Joulain et al. 1998). However, the relative ratios of CH and CH<sub>3</sub> are unlikely to change in such models.
Given the detection of solid CH<sub>4</sub> and the low inferred temperatures, it is plausible that gas-solid interactions and grain-surface chemistry also play a role in producing the hydrocarbons. In this respect, the situation for CH<sub>3</sub> may be similar to that for NH in diffuse clouds (Mann & Williams 1984, van Dishoeck 1998). Conversion of atomic carbon to small hydrides on grain surfaces may be significant, but no model results exist yet for these conditions. Such models should also explain the C<sub>2</sub>H<sub>2</sub> abundances and lack of complete C<sub>2</sub>H<sub>2</sub> freeze-out. Alternatively, reactions of atomic H with PAHs and solid aliphatic hydrocarbon material, known to be present toward Sgr A from the 3.4 $`\mu `$m absorption feature, may lead to CH<sub>3</sub>. Shock chemistry is not likely to be important for this line of sight because of the low temperatures.
Future high spectral resolution observations of CH<sub>3</sub> toward Sgr A to constrain the velocity structure, as well as observations of CH<sub>3</sub> and other molecules in different types of diffuse clouds are needed to constrain the basic hydrocarbon chemistry.
The authors are grateful to the SWS instrument teams, to E. Herbst, R. Terzieva and D.J. Jansen for updated model results on CH<sub>3</sub> and to W. Duley for inspiring discussions. This work was supported by DARA grants no 50 QI9402 3 and 50 QI 8610 8 and by NWO grant 614.41.003. CMW acknowledges receipt of an ARC Australian Postdoctoral Fellowship.
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# Application of the Exact Muffin-Tin Orbitals Theory: the Spherical Cell Approximation
## I Introduction
During the last decades many attempts have been made to develop accurate and at the same time efficient methods for solving the Kohn-Sham equations in an application of the Density Functional Theory for condensed matter. The accuracy of the methods is crucial e.g. when one searches for the answers given by different density functional approximations. The full-potential techniques have been specially designed to fulfill this requirement. Though, in principle, they give highly accurate results, they have their own limitations. In many cases a compromise has been made between the accuracy and efficiency, and methods based on approximate one electron potentials have been developed. The most commonly used muffin tin approach, albeit its mathematical formulation is very elegant, presents a rather poor representation of the exact potential. Though, the Atomic Sphere Approximation (ASA) brings a real improvement to the potential, most of the conventional methods based on the ASA use similar approximation to the one electron energies and charge density as well . Therefore, using these methods, reasonably accurate results can only be obtained for close packed systems, and they are not suitable to treat systems of low symmetry. In order to maintain or increase the accuracy different corrections should be included and, therefore, the ASA based methods lose their elegance and efficiency.
A few years ago breakthrough was made by developing the Exact Muffin-Tin Orbitals (EMTO) Theory . Within the EMTO Theory the one electron states are calculated exactly for the overlapping muffin-tin potential, while the solution of Poisson’s equation can include certain shape approximations, if required. By separating the two approaches used for the potential and one electron states the accuracy can be sustained at a level comparable with that of the full-potential techniques without losing significantly from the efficiency. The EMTO Theory can be considered as an improved screened KKR method , within that large overlapping potential spheres can be used for accurate representation of the exact one electron potential .
In this work we present a self-consistent implementation of the EMTO Theory within the Spherical Cell Approximation (SCA) for the Poisson’s equation. In the first part we review the EMTO Theory , the definition of the screened spherical waves and the matching equation. Furthermore, we establish the expressions for the number of states and electron density using the Green’s function formalism. In the second part of the paper we discuss the impact of the SCA, used for the shape of the Wigner-Seitz cell, on the total charge density and on the overlapping muffin tin potential. Finally, we establish the accuracy of the SCA-EMTO method by performing test calculations for some systems where reliable full-potential data are available. An approximate solution of the kink cancellation equation in order to reduce the number of iterations needed in a self-consistent calculation is presented in the Appendix.
## II Overview: the EMTO Theory
In the following we review the basic concepts of the EMTO Theory developed by O. K. Andersen and co-workers . Assume that the one-electron Kohn-Sham equations,
$$\left[^2+v(𝐫)\right]\mathrm{\Psi }_j(𝐫)=ϵ_j\mathrm{\Psi }_j(𝐫),$$
(1)
are solved within the muffin tin approximation for the effective potential,
$$v(𝐫)v_0+\underset{R}{}\left[v_R(r_R)v_0\right],$$
(2)
where $`R`$ runs over the lattice sites. Here and in the following we use the notation $`𝐫_R𝐫𝐑`$ and omit the vector notation for the index $`R`$. In (2) $`v_0`$ denotes a parameter that reduces to the muffin tin zero in the case of non-overlapping muffin tins. The spherical potentials $`v_R(r_R)`$ become equal to $`v_0`$ outside the potential spheres of radii $`s_R`$. These radii can be chosen as the linear overlap between the spheres to be as large as $`3040\%`$ . It has turned out that for a good representation of the real full-potential in terms of overlapping muffin tin wells usually a big overlap is preferred between the potential spheres .
In order to solve the Schrödinger equation (1) for the muffin tin potential (2) one chooses different basis functions inside the potential spheres and in the interstitial region. Inside the sphere at R the partial waves are chosen as the basis function which are defined as the products of the regular solutions of the radial Schrödinger equation,
$$\frac{^2\left[r_R\varphi _{Rl}(ϵ,r_R)\right]}{r_{R}^{}{}_{}{}^{2}}=\left[\frac{l(l+1)}{r_R^2}+v_R(r_R)ϵ\right]r_R\varphi _{Rl}(ϵ,r_R),$$
(3)
and the real spherical harmonics, viz.
$$\varphi _{RL}(ϵ,𝐫_R)=\varphi _{Rl}(ϵ,r_R)Y_L(\widehat{r}_R),$$
(4)
where $`L=(l,m)`$. Outside the spheres the so called screened spherical waves, $`\psi _{RL}^a(\kappa ,𝐫_R)`$, are used as basis functions. Therefore, the wave function for the energy $`ϵ_j`$ can be written as
$$\mathrm{\Psi }_j(𝐫)=\underset{RL}{}\varphi _{RL}(ϵ_j,𝐫_R)\mathrm{\Theta }_R(𝐫_R)u_{RL,j}^a+\underset{RL}{}\psi _{RL}^a(\kappa _j,𝐫_R)[1\mathrm{\Theta }_R(𝐫_R)]v_{RL,j}^a.$$
(5)
Here $`\kappa ^2=ϵv_0`$, and in the non-overlapping muffin tins limit it denotes the interstitial one-electron kinetic energy. The $`\mathrm{\Theta }_R(𝐫_R)`$ is one inside the sphere of radius $`s_R`$ centered at R and zero outside. The expansion coefficients $`u_{RL,j}^a`$ and $`v_{RL,j}^a`$ as well as the energies $`ϵ_j`$ are determined from the condition that the wave function $`\mathrm{\Psi }_j(𝐫)`$ and its first derivative should be continuous at the potential spheres. The algebraic formulation of this matching condition in the EMTO formalism is the so the called kink cancellation equation, which is equivalent to the KKR (Korringa-Kohn-Rostoker) equation in an arbitrarily screened representation .
### A The screened spherical waves
The screened spherical waves can be defined in conjunction with hard spheres centered at all sites $`𝐑`$ with radii $`a_{Rl}`$. They are solutions of the wave equation,
$$\left[^2+\kappa ^2\right]\psi _{RL}^a(\kappa ^2,𝐫_R)=\mathrm{\hspace{0.33em}0},$$
(6)
with the boundary condition that on their own $`a`$spheres they behave like a pure real spherical harmonic, while the $`Y_L^{}(\widehat{r}_R^{})`$ projections on all the other $`a`$spheres, $`R^{}R`$, vanish. They form a complete basis set in the "$`a`$" interstitial region and may be expressed in terms of the "value", $`f_{RL}^a`$, and the "slope", $`g_{RL}^a`$, functions , whose radial part satisfy the following boundary conditions
$`f_{Rl}^a(\kappa ,r)|_{a_{Rl}}=\mathrm{\hspace{0.33em}1}\text{and}{\displaystyle \frac{f_{Rl}^a(\kappa ,r)}{r}}|_{a_{Rl}}=\mathrm{\hspace{0.33em}0},`$ (7)
$`g_{Rl}^a(\kappa ,r)|_{a_{Rl}}=\mathrm{\hspace{0.33em}0}\text{and}{\displaystyle \frac{g_{Rl}^a(\kappa ,r)}{r}}|_{a_{Rl}}={\displaystyle \frac{1}{a_{Rl}}}.`$ (8)
These functions may of course be expressed in terms of the usual spherical Bessel and Neumann functions
$$f_{Rl}^a(\kappa ,r)=j_l(\kappa r)𝒲_a\{f,\kappa n\}\kappa n_l(\kappa r)𝒲_a\{f,j\}$$
(9)
and
$$g_{Rl}^a(\kappa ,r)=j_l(\kappa r)𝒲_a\{g,\kappa n\}\kappa n_l(\kappa r)𝒲_a\{g,j\}$$
(10)
since $`𝒲_r\{j,n\}=1/\kappa `$, and therefore satisfy the Wronskian
$$𝒲_r\{f_{Rl}^a,g_{Rl}^a\}r^2\left[f_l^a(\kappa ,r)\frac{g_l^a(\kappa ,r)}{r}\frac{f_l^a(\kappa ,r)}{r}g_l^a(\kappa ,r)\right]=a_{Rl}.$$
(11)
The screened spherical wave $`\psi _{RL}^a(\kappa ^2,𝐫_R)`$ may be expanded in spherical harmonics $`Y_L^{}(\widehat{r}_R^{})`$ about any site $`𝐑^{}`$, as
$$\psi _{RL}^a(\kappa ^2,𝐫_R)=f_{Rl}^a(\kappa ,r_R)Y_L(\widehat{r}_R)\delta _{RR^{}}\delta _{LL^{}}+\underset{L^{}}{}g_{R^{}l^{}}^a(\kappa ,r_R^{})Y_L^{}(\widehat{r}_R^{})S_{R^{}L^{}RL}^a(\kappa ^2),$$
(12)
where the expansion coefficients, $`S_{R^{}L^{}RL}^a(\kappa ^2)`$, are the elements of the so called slope matrix. The slope matrix can be derived from the bare KKR structure constant matrix $`B_{R^{}L^{},RL}(\kappa )`$, by matrix inversion
$$S^a(\kappa ^2)=𝒟\{j(\kappa ,a)\}\frac{1}{aj(\kappa ,a)}\left[B(\kappa )+\kappa cot\alpha (\kappa )\right]^1\frac{1}{j(\kappa ,a)},$$
(13)
where $`𝒟`$ denotes the logarithmic derivative, $`𝒟\{j(r)\}r\left[j(r)/r\right]/j(r)`$, and for simplicity where we have used matrix notation. We note that this equation is equivalent to Eq. (3.26) from Ref. and Eq. (15) from Ref. . The bare KKR structure constants are defined as the expansion coefficients of the $`\kappa n_L(\kappa ,𝐫_R)\kappa n_l(\kappa r_R)Y_L(\widehat{r}_R)`$ functions around site $`R^{}`$ in terms of the $`j_L^{}(\kappa ,𝐫_R^{})j_l^{}(\kappa r_R^{})Y_L^{}(\widehat{r}_R^{})`$ functions, i.e.
$`\kappa n_L(\kappa ,𝐫_R)={\displaystyle \underset{L^{}}{}}j_L^{}(\kappa ,𝐫_R^{})B_{R^{}L^{}RL}(\kappa ),`$with
$$B_{R^{}L^{}RL}(\kappa )\mathrm{\hspace{0.33em}4}\pi \underset{L^{\prime \prime }}{}C_{LL^{}}^{L^{\prime \prime }}i^{l+l^{}l^{\prime \prime }}\kappa n_{L^{\prime \prime }}(\kappa ,𝐑^{}𝐑),$$
(14)
and where $`C_{LL^{}}^{L^{\prime \prime }}`$ are the real Gaunt numbers.
For the partial waves explicitly included in the formalism, the so called low partial waves with $`ll_{low}=23`$, $`\alpha _{Rl}(\kappa )`$ are the hard sphere phase shifts given by
$`cot\alpha _{Rl}(\kappa )=n_l(\kappa a_{Rl})/j_l(\kappa a_{Rl})`$and for the remaining $`Rl`$-chanels, $`\alpha _{Rl}(\kappa )`$ are the proper phase shifts. For high $`l`$’s the latter vanish, and at that point the matrix to be inverted in (13) can be truncated.
When the hard sphere radii, $`a_R`$, are properly chosen and $`\kappa ^2`$ lies below the bottom of the hard sphere continuum, the screened spherical waves have short range. Therefore, the slope matrix can be calculated in real space and the method is suitable to treat impurities, defects, surfaces, etc. It was shown in Ref. that the shortest range of the screened spherical waves can be achieved for non-overlapping spheres with $`a_R0.50.85s_R^i`$, depending on the maximal orbital quantum number $`l`$ of the partial waves explicitly included in the formalism. The $`s_R^i`$ denotes the inscribed or touching sphere radii. In the KKR community, it is customary to determine the $`\alpha _{Rl}(\kappa )`$’s as the phase shifts of repulsive potentials.
Because a screened spherical wave has pure $`(l,m)`$ character only on its own $`a`$-sphere, the matching condition in Eq. (5) should be set up at this sphere. The connection onto the potential sphere (s) is done by introducing a free electron solution $`\phi _{Rl}(ϵ,r_R)Y_L(\widehat{r}_R)`$ from the potential sphere back to the hard sphere, which joins continuously and differentiable to the partial wave, $`\varphi _{RL}(ϵ,𝐫_R)`$, at $`s_R`$ and continuously to the screened spherical wave at $`a_{Rl}`$. The radial part of this backwards extrapolated free-electron solution, after normalizing it to one at its $`a`$-sphere, is given by
$$\phi _{Rl}^a(ϵ,r)\frac{\phi _{Rl}(ϵ,r)}{\phi _{Rl}(ϵ,a_{Rl})}=f_{Rl}^a(\kappa ,r)+g_{Rl}^a(\kappa ,r)D_{Rl}^a(ϵ),$$
(15)
where $`D_{Rl}^a(ϵ)`$ is the logarithmic derivative of $`\phi _{Rl}(ϵ,r)`$ calculated at the hard sphere $`a_{Rl}`$. This can be determined from the matching condition between $`\varphi _{RL}(ϵ,r)`$ and $`\phi _{RL}(ϵ,r)`$ at $`r_R=s_R`$,
$$D_{Rl}^a(ϵ)𝒟\{\phi _{Rl}^a(ϵ,a_{Rl})\}=\frac{f_{Rl}^a(\kappa ,s_R)}{g_{Rl}^a(\kappa ,s_R)}\frac{𝒟\{\varphi _{Rl}(ϵ,s_R)\}𝒟\{f_{Rl}^a(\kappa ,s_R)\}}{𝒟\{\varphi _{Rl}(ϵ,s_R)\}𝒟\{g_{Rl}^a(\kappa ,s_R)\}}.$$
(16)
The relation between the values of the free electron function at $`a`$ and the partial wave at $`s`$ is
$$\frac{\phi _{Rl}(ϵ,a_{Rl})}{\varphi _{Rl}(ϵ,s_R)}=\frac{\phi _{Rl}(ϵ,a_{Rl})}{\phi _{Rl}(ϵ,s_R)}=\frac{1}{f_{Rl}^a(\kappa ,s_R)}\frac{𝒟\{\varphi _{Rl}(ϵ,s_R)\}𝒟\{g_{Rl}^a(\kappa ,s_R)\}}{𝒟\{f_{Rl}^a(\kappa ,s_R)\}𝒟\{g_{Rl}^a(\kappa ,s_R)\}},$$
(17)
In Fig. 1 we have plotted the logarithmic derivative at $`a=0.7w`$, where $`w`$ denotes the average Wigner-Seitz radius, and the normalization function $`\phi _{Rl}(ϵ,a_{Rl})`$ given in (17) in the case of fcc Ga. The logarithmic derivative is a never increasing function of energy and it has a pole above the top of the respective band. Between these poles $`D_{Rl}^a(ϵ)`$ is smooth functions of energy, which varies more slowly than $`𝒟\{\varphi _{Rl}(ϵ,s)\}`$, because $`a<s`$. The poles of $`D_{Rl}^a(ϵ)`$ depend on the representation $`(a)`$ and they are not related directly to the band structure. The $`\phi _{Rl}(ϵ,a_R)`$ from the figure was obtained for partial waves normalized in the $`w`$-sphere. It is always a smooth function of the energy and vanishes at the poles of $`D_{Rl}^a(ϵ)`$.
The slope matrix, Eq. (13), the logarithmic derivative, Eq. (16), and the normalization function, Eq. (17), play a central role in the present implementation of the EMTO Theory.
### B Kink cancellation equation
Using the free electron solutions from (12) and (15) and the partial waves $`\varphi _{Rl}(ϵ,𝐫_R)`$ we can introduce a complete basis set defined in the whole space. These exact muffin tin orbitals or kinked partial waves may be written in the form
$$\overline{\psi }_{RL}^a(ϵ,𝐫_R)=\left(\varphi _{Rl}^a(ϵ,r_R)\phi _{Rl}^a(ϵ,r_R)\right)Y_L(\widehat{r}_R)+\psi _{RL}^a(\kappa ^2,𝐫_R),$$
(18)
where the radial part of the functions $`\varphi _{Rl}^a`$ and $`\phi _{Rl}^a`$ are truncated outside the sphere of radius $`s_R`$ and outside $`s_R`$ and inside $`a_R`$ , respectively. Moreover, the $`ll_{low}`$ projection of the $`\psi _{Rl}^a`$ function is truncated inside the sphere of radius $`a_R`$, while the high-$`l`$ components penetrate into the hard spheres. The $`\overline{\psi }_{RL}^a(ϵ,𝐫_R)`$ functions are continuous and differentiable in the whole space, except at the hard spheres, where they have non zero kinks. In Eq. (18) the partial waves are renormalized according to Eq. (15)
$$\varphi _{RL}^a(ϵ,r_R)\frac{\varphi _{RL}(ϵ,r_R)}{\phi _{RL}(ϵ,a_R)}.$$
(19)
From Eq. (17) and (19) it is immediately seen that the multiplicativ normalization of the partial waves does not enter in the expression of the kinked partial wave. Forming a linear combination of the kinked partial waves,
$$\mathrm{\Psi }_j(𝐫)=\underset{RL}{}\overline{\psi }_{RL}^a(ϵ_j,𝐫_R)v_{RL,j}^a,$$
(20)
and asking for the kinks be canceled we arrive to the kink cancellation or screened KKR equations
$`{\displaystyle \underset{RL}{}}K_{R^{}L^{}RL}^a(ϵ_j)v_{RL,j}^a`$ $``$ (21)
$`{\displaystyle \underset{RL}{}}a_R^{}\left[S_{R^{}L^{}RL}^a(\kappa _j^2)\delta _{R^{}R}\delta _{L^{}L}D_{RL}^a(ϵ_j)\right]v_{RL,j}^a`$ $`=`$ $`\mathrm{\hspace{0.33em}0}\text{for all}R^{}L^{}.`$ (22)
Here we have $`l,l^{}l_{low}`$. The solutions of this equation are the one-electron energies and eigenfunctions, which, using Eq. (5) are given by
$$u_{RL,j}=\frac{v_{RL,j}^a}{\phi _{Rl}(ϵ_j,a_R)}.$$
(23)
It is worth to note that in the final expression of the wavefunction $`\mathrm{\Psi }_j(𝐫)`$, Eq. (5), the backwards extrapolated free electron solution does not enters.
In the case of translation symmetry in Eq. (21) $`R`$ and $`R^{}`$ run over the atoms in the primitive cell only, and the slope matrix, and thus the kink matrix $`K_{R^{}L^{}RL}^a`$ as well, depend on the Bloch vector $`𝐤`$ from the first Brillouin zone. In Fig. 2 we plotted the diagonal elements of the fcc slope matrix (symbols) calculated at the center of the Brillouin zone as a function of the dimensionless energy parameter $`(\kappa w)^2`$. In this calculation the matrix inversion in (13) was performed in real space for $`5`$ coordination shells plus the central site using the $`s,p`$ and $`d`$ orbitals and $`0.7w`$ for the hard sphere radius. The figure demonstrates the weak and smooth energy dependence of the slope matrix up to the bottom of the continuum, $`(\kappa w)^26`$. Therefore in the practical solution of the kink cancellation equation (21) the slope matrix can be estimated using a Taylor expansion around a fixed energy $`\kappa _0^2`$,
$$S_{R^{}L^{}RL}^a(\kappa ^2)=S_{R^{}L^{}RL}^a(\kappa _0^2)+\frac{1}{1!}\dot{S}_{R^{}L^{}RL}^a(\kappa _0^2)(\kappa ^2\kappa _0^2)+\mathrm{},$$
(24)
where the overdot indicates energy derivative. The first and higher order energy derivatives are calculated analytically as described in Ref. . In equation (24) $`\kappa ^2`$ is a complex energy not too far from $`\kappa _0^2`$. In Fig. 2 the solid lines were calculated with a fourth order expansion around $`\kappa _0=0`$. As one can observe, this expansion gives highly accurate energy dependence of the slope matrix over an energy range of approximately $`(1,+1)Ry`$.
### C The electron density
In order to construct the new one-electron potential for a self-consistent calculation first we need to construct the electron density given by
$$n(𝐫)=\underset{j}{\overset{ϵ_jϵ_F}{}}|\mathrm{\Psi }_j(𝐫)|^2,$$
(25)
where the sum runs over the one-electron states below the Fermi level $`ϵ_F`$. In the present implementation of the method instead of calculating explicitly the wave functions (5) and performing the summation in Eq. (25) we introduce the path operator $`g_{R^{}L^{}RL}^a(z,𝐤)`$ defined for a complex energy $`z`$ and Bloch vector $`𝐤`$ by
$$\underset{R^{\prime \prime }L^{\prime \prime }}{}K_{R^{}L^{}R^{\prime \prime }L^{\prime \prime }}^a(z,𝐤)g_{R^{\prime \prime }L^{\prime \prime }RL}^a(z,𝐤)=\delta _{R^{}R}\delta _{L^{}L}.$$
(26)
This function is analytic in the complex plane and it has poles at the one-electron energies along the real axis. Therefore, using the residue theorem, for the total number of electrons we find
$$N(ϵ_F)=\frac{1}{2\pi i}_{ϵ_F}\underset{R^{}L^{}RL}{}_{BZ}g_{R^{}L^{}RL}^a(z,𝐤)\dot{K}_{RLR^{}L^{}}^a(z,𝐤)𝑑𝐤𝑑z,$$
(27)
where the first integration is performed on a complex contour and the second one in the first Brillouin zone. The contour is chosen in a way that it cuts the real axis below the bottom of the valence band and at $`ϵ_F`$. In (27) $`l,l^{}l_{low}`$. The $`z`$ dependent partial waves, $`\varphi _{Rl}(z,r_R)`$, and logarithmic derivatives, $`D_{Rl}^a(z)`$, are obtained by solving Eq. (3) for complex energy. The energy derivative of the kink matrix,
$$\dot{K}_{R^{}L^{}RL}^a(z,𝐤)=a_R^{}\left[\dot{S}_{R^{}L^{}RL}^a(zv_0,𝐤)\delta _{R^{}R}\delta _{L^{}L}\dot{D}_{RL}^a(z)\right],$$
(28)
is calculated by taking the derivatives of Eq. (16) and (24), where the energy derivatives of the basis functions $`\{f^a,g^a\}`$ are calculated analytically. The energy derivative of the logarithmic derivative function is given by
$$\frac{𝒟\{\varphi _{Rl}(z,s_R)\}}{z}=\frac{_0^{s_R}\varphi _{Rl}^2(z,r_R)r_R^2𝑑r_R}{s_R\varphi _{Rl}^2(z,s_R)}.$$
(29)
Because the eigenvectors are normalized as (see Ref. )
$$\mathrm{\Psi }_j^{}(𝐫)\mathrm{\Psi }_j(𝐫)𝑑𝐫=\underset{R^{}L^{}RL}{}v_{R^{}L^{},j}^a\dot{K}_{R^{}L^{}RL}^a(ϵ)v_{RL,j}^a$$
(30)
the expression (27) gives the exact number of states at the Fermi level for the muffin tin potential (2). In (30) the negligible terms due to the overlap between $`s`$-spheres are omitted .
Inside the unit cell at $`R`$ the electron density in terms of the path operator can be expressed as
$$n(𝐫_R)=\frac{1}{2\pi i}_{ϵ_F}\underset{L^{}L}{}Z_{RL^{}}^a(z,𝐫_R)_{BZ}\stackrel{~}{g}_{RL^{}RL}^a(z,𝐤)𝑑𝐤Z_{RL}^a(z,𝐫_R)𝑑z,$$
(31)
where we have introduced the functions
$`Z_{RL}^a(z,𝐫_R)=\{\begin{array}{ccc}\varphi _{Rl}^a(z,r_R)Y_L(\widehat{r}_R)\hfill & \text{if}l_{low}\text{and}r_Rs_R\hfill & \\ \phi _{Rl}^a(z,r_R)Y_L(\widehat{r}_R)\hfill & \text{if}l_{low}\text{and}r_R>s_R\hfill & \\ j_l(z,r_R)Y_L(\widehat{r}_R)\hfill & \text{if}l>l_{low}\text{for all}r_R\hfill & \end{array},`$ (35)
and where the sums over $`l^{}`$ and $`l`$ include the high-$`l`$ terms as well. These functions are equivalent to the scattering solutions of Faulkner and Stocks . In Eq. (31) we have introduced the following matrix
$`\stackrel{~}{g}_{RL^{}RL}^a(z,𝐤)\{\begin{array}{ccccc}g_{RL^{}RL}^a(z,𝐤)\hfill & \text{if}l,l^{}l_{low}\hfill & & & \\ _{R^{\prime \prime }L^{\prime \prime }}g_{RL^{}R^{\prime \prime }L^{\prime \prime }}^a(z,𝐤)S_{R^{\prime \prime }L^{\prime \prime }RL}^a(z,𝐤)\hfill & \text{if}l^{}l_{low}\text{and}l>l_{low}\hfill & & & \\ _{R^{\prime \prime }L^{\prime \prime }}S_{RL^{}R^{\prime \prime }L^{\prime \prime }}^a(z,𝐤)g_{R^{\prime \prime }L^{\prime \prime }RL}^a(z,𝐤)\hfill & \text{if}l^{}>l_{low}\text{and}ll_{low}\hfill & & & \\ _{R^{\prime \prime }L^{\prime \prime }}_{R^{\prime \prime \prime }L^{\prime \prime \prime }}S_{RL^{}R^{\prime \prime }L^{\prime \prime }}^a(z,𝐤)\hfill & & & & \\ \times g_{R^{\prime \prime }L^{\prime \prime }R^{\prime \prime \prime }L^{\prime \prime \prime }}^a(z,𝐤)S_{R^{\prime \prime \prime }L^{\prime \prime \prime }RL}^a(z,𝐤)\hfill & \text{if}l^{},l>l_{low}\hfill & & & \end{array}`$ (41)
where the high-low and the low-high subblocks of the slope matrix are calculated by the usual blowing-up technique .
In principle Eq. (27) and (31) give the exact number of states and electron density. However, in some cases, like for the metals from the $`IIB`$ and $`IIIVA`$ groups, where one of the $`d`$ bands is completely filled, around the top of this band the normalization function (17) goes through zero. This happens, for example, in the case of fcc Ga around the top of the $`3d`$ band, as it can be seen from Fig. 1. For this energy not only the logarithmic derivative but also its energy derivative $`\dot{D}_{Rl}^a`$, appearing in the diagonal of the $`\dot{K}_{R^{}L^{}RL}^a`$ matrix, has poles. In order to cancel these nonphysical poles we rewrite the expression for the number of states as
$`N(ϵ_F)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{ϵ_F}}G(z)𝑑z,`$ (42)
where
$`G(z){\displaystyle \underset{R^{}L^{}RL}{}}{\displaystyle _{BZ}}g_{R^{}L^{}RL}^a(z,𝐤)\dot{K}_{RLR^{}L^{}}^a(z,𝐤)𝑑𝐤{\displaystyle \underset{RL}{}}\left[{\displaystyle \frac{\dot{D}_{Rl}^a(z)}{D_{Rl}^a(z)}}{\displaystyle \underset{ϵ_{Rl}^D}{}}{\displaystyle \frac{1}{zϵ_{Rl}^D}}\right],`$ (43)
with $`l,l^{}l_{low}`$, and that of the electron density as
$`n(𝐫_R)`$ $`=`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{ϵ_F}}{\displaystyle \underset{L^{}L}{}}Z_{RL^{}}^a(z,𝐫_R){\displaystyle _{BZ}}\stackrel{~}{g}_{RL^{}RL}^a(z,𝐤)𝑑𝐤Z_{RL}^a(z,𝐫_R)𝑑z`$ (44)
$`+`$ $`{\displaystyle \frac{1}{2\pi i}}{\displaystyle _{ϵ_F}}{\displaystyle \underset{L}{}}{\displaystyle \frac{Z_{}^{a}{}_{RL}{}^{2}(z,𝐫_R)}{a_RD_{Rl}^a(z)}}dz{\displaystyle \underset{L}{}}{\displaystyle \underset{ϵ_{Rl}^D}{}}{\displaystyle \frac{Z_{}^{a}{}_{RL}{}^{2}(ϵ_{Rl}^D,𝐫_R)}{a_R\dot{D}_{Rl}^a(ϵ_{Rl}^D)}},`$ (45)
where $`ϵ_{Rl}^D`$ are the zeros of the logarithmic derivative function, $`D_{Rl}^a(ϵ)`$. Because the logarithmic derivative is a smooth decreasing function of energy $`ϵ_{Rl}^D`$’s can be easily determined with high accuracy. The second and third terms from the right hand side of (44) are included only for $`ll_{low}`$. Using the fact that the residuum of the $`1/D_{Rl}^a`$ around $`ϵ_{Rl}^D`$ is $`1/\dot{D}_{Rl}^a`$, it is easy to show that the poles of $`\dot{D}_{Rl}^a(z)`$ and those of $`1/\phi _{Rl}(z,a_R)^2`$ are canceled out in Eqs. (43) and (44).
From the electron density (44) we determine the overlapping muffin tin wells and repeat the iterations until self consistency, of the total energy for example, is achieved. In the present implementation of the EMTO Theory the solution of the Poisson’s equation involves the SCA for the shape of the Wigner-Seitz cell, therefore, the construction of the muffin tin potential will be discussed only within this context.
## III The SCA-EMTO method
Equations (21) and (42-44), derived in the previous section, constitute the basis of the present method. In order to perform a self-consistent calculation one constructs the electron density from the solutions of the kink cancellation equation and calculates the new one-electron potential. In this section we describe these steps using the SCA for the shape of the Wigner-Seitz cell.
In the SCA, for solving the Poisson’s equation, we substitute the Wigner-Seitz cells by spherical cells with volumes equal to the volumes of the real cells. If $`\mathrm{\Omega }_R`$ denotes the volume of the Wigner-Seitz cell (Voronoi polyhedron) centered at $`𝐑`$ we have $`\mathrm{\Omega }_R=\mathrm{\Omega }_{w_R}\frac{4\pi }{3}w_R^3`$, where $`w_R`$ is the atomic sphere radius. Thus within the SCA, like in the conventional ASA, the whole space is "covered" by the $`\mathrm{\Omega }_{w_R}`$ spheres.
### A The SCA charge density
During the self-consistent calculation the Fermi level of a $`N`$ electron system is determined by solving the $`N(ϵ_F)=N`$ equation, where $`N(ϵ_F)`$ is given in (42). For this $`ϵ_F`$ the electron density is constructed from Eq. (44). Due to the normalization (30) the so constructed density is exactly normalized within the unit cell cell but not within the SCA spheres of volumes $`\mathrm{\Omega }_{w_R}`$. Therefore, in order to solve the Poisson’s equation within the SCA we have to renormalize the total density inside the spheres. In the present implementation of the method this is realized by
$$n^{SCA}(𝐫_R)=n(𝐫_R)+aY_{00}(\widehat{𝐫}_R),$$
(46)
where the site independent $`a`$ constant is determined from the condition of the charge neutrality within the whole unit cell
$$\underset{R}{}_{\mathrm{\Omega }_{w_R}}n^{SCA}(𝐫_R)𝑑𝐫_R=\underset{R}{}Z_R.$$
(47)
Here $`Z_R`$ denotes the nuclear charge at $`R`$. The sum runs over the atoms from the unit cell, and the integrals are performed inside the SCA spheres. Throughout this section the charge density is normalized within the SCA spheres according to (46) and (47), however, for the sake of simplicity we neglect the $`SCA`$ index for the $`n^{SCA}(𝐫)`$.
### B The SCA muffin tin potential
The spherical symmetric potentials, $`v_R(r_R)`$, that enter in Eq. (3) have to be chosen in a way that, together with the parameter $`v_0`$, to give the best approximation to the full potential $`v(𝐫)`$. The original idea in Ref. is to minimize the mean of the squared deviation between the left and the right hand side of Eq. (2). This leads to a set of integral or differential equations for $`v_R(r_R)`$ and $`v_0`$. In the non-overlapping muffin tins case the equation for $`v_R(r_R)`$ reduce to the well known expression
$$v_R(r_R)=\frac{1}{4\pi }v(𝐫)𝑑\widehat{𝐫}_R,$$
(48)
and $`v_0`$ reduces to the muffin tin zero, i.e. to the average of the full potential calculated in the interstitial region,
$`\mathrm{\Omega }^I\mathrm{\Omega }{\displaystyle \underset{R}{}}V_R\mathrm{\Omega }{\displaystyle \underset{R}{}}{\displaystyle \frac{4\pi }{3}}s_R^3,`$where $`\mathrm{\Omega }`$ is the volume of the region where the approximation (2) is valid (unit cell), and $`V_R`$ denotes the volume of the potential sphere.
In the overlapping muffin tins case the equation for the $`v_0`$ can be written in the following simple form
$$\underset{R}{}\frac{4\pi }{\mathrm{\Omega }}_0^{s_R}\left[v_R(r_R)v_0\right]r_R^2𝑑r+v_0=\frac{1}{\mathrm{\Omega }}_\mathrm{\Omega }v(𝐫)𝑑𝐫,$$
(49)
while the equation for $`v_R(r_R)`$ involves terms coming from the overlapping region, and which give rise to kinks of $`v_R(r_R)`$ when $`r_R`$ touches other muffin tin spheres. In the present implementation of the method, instead of solving the $`v_R(r_R)`$ equations, we all the time, for non-overlapping and for overlapping muffin tin wells as well, fix the $`v_R(r_R)`$ functions to the spherical average of the full potential given by (48). In this case from Eq. (49) we get the expression for the $`v_0`$ as
$$v_0=\frac{1}{\mathrm{\Omega }_RV_R}\left[_\mathrm{\Omega }v(𝐫)𝑑𝐫\underset{R}{}_{V_R}v(𝐫)𝑑𝐫\right],$$
(50)
or
$$v_0=\frac{_R\left[_{\mathrm{\Omega }_R^I}v(𝐫)𝑑𝐫_{\mathrm{\Omega }_R^{ov}}v(𝐫)𝑑𝐫\right]}{_R\left[\mathrm{\Omega }_R^I\mathrm{\Omega }_R^{ov}\right]},$$
(51)
where $`\mathrm{\Omega }_R^I`$ is the real interstitial within a Wigner-Seitz cell centered at $`R`$ with volume $`\mathrm{\Omega }_R`$, and $`\mathrm{\Omega }_R^{ov}`$ is that part of the potential sphere that is outside of the cell $`\mathrm{\Omega }_R`$, i.e.
$$\mathrm{\Omega }\underset{R}{}\mathrm{\Omega }_R\text{and}\mathrm{\Omega }_R^I\mathrm{\Omega }_R^{ov}\mathrm{\Omega }_RV_R.$$
(52)
Eq. (51) assumes the knowledge of the full potential $`v(𝐫)`$ in $`\mathrm{\Omega }_R^I`$ and $`\mathrm{\Omega }_R^{ov}`$ regions. However, the time consuming calculation of the full potential can be avoided by using the SCA for the unit cell. In the non-overlapping SCA case, i.e. $`s_R<w_R`$, we have
$`\mathrm{\Omega }_R^I\mathrm{\Omega }_R^{ov}=4\pi {\displaystyle _{s_R}^{w_R}}r_R^2𝑑r_R,`$and
$$_{\mathrm{\Omega }_R^I}v(𝐫)𝑑𝐫_{\mathrm{\Omega }_R^{ov}}v(𝐫)𝑑𝐫=_{s_R}^{w_R}\left[v(𝐫)𝑑\widehat{𝐫}_R\right]r_R^2𝑑r_R,$$
(53)
while for the overlapping SCA case, i.e. $`s_R>w_R`$, we have
$`\mathrm{\Omega }_R^I\mathrm{\Omega }_R^{ov}=4\pi {\displaystyle _{w_R}^{s_R}}r_R^2𝑑r_R,`$and
$$_{\mathrm{\Omega }_R^I}v(𝐫)𝑑𝐫_{\mathrm{\Omega }_R^{ov}}v(𝐫)𝑑𝐫=_{w_R}^{s_R}\left[v(𝐫)𝑑\widehat{𝐫}_R\right]r_R^2𝑑r_R.$$
(54)
From these equations we get the expression for the parameter $`v_0`$ valid within the SCA
$$v_0=\underset{R}{}_{s_R}^{w_R}r_R^2\left[v(𝐫)𝑑\widehat{𝐫}_R\right]𝑑r_R/\underset{R}{}W_R$$
(55)
where $`W_R4\pi (w_R^3s_R^3)/3`$. Therefore both of the $`v_R(r_R)`$ function and the $`v_0`$ parameter are given in terms of the spherical symmetric part of the full potential.
The many-body part, $`\mu _{xc}[n(𝐫)]`$, of the one-electron effective potential,
$$v(𝐫)=v^C(𝐫)+\mu _{xc}[n(𝐫)],$$
(56)
is calculated within the local density or generalized gradient approximation, while the electrostatic part is derived solving the Poisson’s equation,
$$^2v^C(𝐫)=8\pi \left[n(𝐫)\underset{R}{}Z_R\delta (r_R)\right],$$
(57)
for the electronic and nuclear charge densities. The electrostatic potential can be divided into intercell and intercell component. The spherical symmetric part of the intercell or Madelung potential is given by
$$v_R^M(r_R)=\frac{1}{w}\underset{R^{}L^{}}{}M_{RLR^{}L^{}}Q_{R^{}L^{}}\text{with}L=(0,0),$$
(58)
where $`M_{RLR^{}L^{}}`$ is the Madelung matrix, which can be evaluated by the usual Ewald technique, and
$$Q_{RL}=\frac{\sqrt{4\pi }}{2l+1}_{\mathrm{\Omega }_{w_R}}\left(\frac{r_R}{w}\right)^l\left[n_R(𝐫_R)Z_R\delta (r_R)\right]Y_L(\widehat{r}_R)𝑑𝐫_R.$$
(59)
The Hartree part of the intracell Coulomb potential can be obtained as the solution of the Poisson’s equation using the proper boundary condition at the atomic sphere radius. Alternatively, this term is given by
$`v_R^I(r_R)=\{\begin{array}{cc}8\pi \left[\frac{1}{r_R}_0^{r_R}r_{R}^{}{}_{}{}^{2}n_R(r_R^{})𝑑r_R^{}+_{r_R}^{w_R}r_R^{}n_R(r_R^{})𝑑r_R^{}\right]\text{for}r_Rw_R\hfill & \\ 8\pi \frac{1}{r_R}_0^{w_R}r_{R}^{}{}_{}{}^{2}n_R(r_R^{})𝑑r_R^{}\text{for}r_R>w_R\hfill & \end{array},`$ (62)
that is valid inside the potential sphere $`s_R`$, for $`s_Rw_R`$ as well as for $`s_R<w_R`$. The total potential within the potential sphere is obtained as the sum of Eq. (58), (62), the Coulomb potential of the nucleus and the spherical symmetric exchange-correlation potential, namely
$$v_R(r_R)=v_R^M+v_R^I(r_R)\frac{2Z_R}{r_R}+\mu _{xcR}(r_R).$$
(63)
If the spherical symmetric part of the exchange-correlation potential is approximated by $`\mu _{xcR}[n_R(r_R)]`$ besides the higher order multipole moments from (58), which in many cases can be neglected, all of the potential components from (63) depend only on the spherical symmetric density $`n_R(r_R)`$.
Within the SCA-EMTO method the atomic and potential spheres can be and usually they are chosen differently. The sizes of the atomic spheres, $`w_R`$, are fixed by the volume, and the ratio between them should be chosen in a way that minimizes the errors coming from approximate solution of the Poisson’s equation. We have found that the best representation of the potential can be achieved by choosing the potential sphere radii, $`s_R`$, larger or equal with the atomic sphere radii. For an optimal choice of the potential spheres the potentials at $`s_R`$ should be the same, i.e. $`v_R(s_R)constant`$ for each $`R`$, and this $`constant`$ should have the maximum possible value for linear overlaps bellow $`3040\%`$.
## IV Applications: test calculations
In this section we present a few applications of the SCA-EMTO method. We chose particular systems where the conventional ASA based methods failed and the inclusion of the correction terms or of the exact potential seemed to be unavoidable. First we describe the most important numerical details and after we analyze the present results comparing to the available full-potential calculations.
### A Numerical details
The hard sphere radii are chosen at $`a_R=0.7w`$. In the matrix inversion from Eq. (13) we includ $`79`$ sites in the case of $`fcc`$-based structures ($`fcc,L1_2`$ and $`L1_0`$), and $`89`$ sites in the case of $`bcc`$-based structures ($`bcc,B2`$ and $`B32`$). The Taylor expansion of the slope matrix is carried out for $`\kappa _0=0`$ and includes terms up to the fourth to sixth order energy derivative.
The path operator is calculated for $`1632`$ complex energy points, depending on the band structure, distributed exponentially on a semi-circular contour. The $`k`$-point sampling is performed on a uniform grid in the $`3D`$ Brillouin zones. All the calculations are scalar-relativistic and employ the frozen-core approximation. The basis set include $`s`$, $`p`$ and $`d`$ orbitals and the valence electrons are treated self-consistently within the local density approximation to density functional theory using the Ceperley and Alder exchange-correlation functional and parametrized by Perdew and Wang . The atomic sphere radii, $`w_R`$-s, are chosen in a way that the atomic spheres should have the same volumes as the corresponding Voronoi polyhedra. The electrostatic and exchange-correlation contribution to the total energy is calculated within the SCA as described, for example, in Refs. . The kinetic energy is given by
$$T^{SCA}=\frac{1}{2\pi i}_{ϵ_F}zG(z)𝑑z\underset{R}{}_0^{w_R}v_R(r_R)n_R(r_R)r_R^2𝑑r_R,$$
(64)
where the first term from the right hand side is the sum of the one electron energies and $`G(z)`$ is given in (43).
### B Results
Before starting on the evaluation of the results we address the question of the accuracy of the Taylor expansion for the slope matrix, Eq. (24). In Fig. 3 the total energy of fcc Cu is shown for different number of terms included in the Taylor expansion. The inclusion of the fourth order energy derivative term changes the total energy by $`1.13`$ mRy. The effect of the fifth order term is already less than $`0.2`$ mRy and that of the sixth order term is about $`0.04`$ mRy. Therefore, we conclude that for a reasonable accuracy it is sufficient to include five terms in the expansion of the slope matrix, viz. up to the fourth order energy derivatives. In the case of open structures, wide energy bands or semicore states, however, more terms should be included .
The variations of the SCA-EMTO total energy with the potential sphere radius for fcc and bcc Cu are shown on Fig. 4. For this test calculation the total potential from (55) was weighted by the valence part of the total density according to
$$v_0=\frac{_R_{s_R}^{w_R}r_R^2\left[n(𝐫)v(𝐫)𝑑\widehat{𝐫}_R\right]𝑑r_R}{_R_{s_R}^{w_R}r_R^2\left[n(𝐫)𝑑\widehat{𝐫}_R\right]𝑑r_R},$$
(65)
and the cut-off in (44) for the spherical part of the total density was $`l_{max}=l_{max}^{}=8`$. The inscribed sphere radii are $`s_{fcc}^i=0.91w_{fcc}`$ and $`s_{bcc}^i=0.88w_{bcc}`$, where the theoretical atomic sphere radii, $`w_{fcc}`$ and $`w_{bcc}`$ are shown on the figure. The total energy in both of the fcc and bcc structures beginning from $`s0.80w`$ becomes almost flat with a negligible slope up to $`s1.20w`$, which means about $`3236\%`$ linear overlap for the fcc and bcc structures, respectively. For $`s>0.80w`$ the further increase of the potential sphere radius has little effect on the energy, that means the potential in the corners of the Wigner-Seitz cell is, with a very good approximation, constant. However, for big overlaps, $`s>1.20w`$, the errors coming from the overlap region, and neglected in the kink-cancellation equation and in the charge density as well, become important .
There is a comprehensive study of the structural stability of the transition metals done either by full-potential or by muffin-tin or ASA based methods. In the latter case correction terms are needed for calculation of the accurate total energies. The conventional ASA without correction terms gives, for example, with about $`2`$ mRy lower total energy for the Cu in the bcc phase than in the fcc phase. This underestimation of the bcc total energy is due to the incorrect kinetic energy term, and the inclusion of the exact Hartree energy would lower even more the bcc energy. From a more sophisticated full-potential method a structural energy difference of $`E_{bcc}E_{fcc}0.5`$ mRy was obtained. This number should be compared with our results of $`0.4`$ mRy from Fig. 4. One should note that this difference is almost constant for a wide range of linear overlap.
The second example is the binary Li<sub>x</sub>Al<sub>1-x</sub> ordered compound in different phases. There are two reasons of this choice: i) this system is very well studied through accurate full-potential calculations , and ii) most of the experimentally observed interesting trends, the contraction of the volume of the Al-based alloys, the asymmetric heats of formation with respect to the equiatomic concentration etc., can not be reproduced by an ordinary spherically symmetric calculation. Besides the three different compositions, $`x=0.25,0.50,0.75`$, we consider the pure Al and Li limits in fcc and bcc phases as well. For $`x=0.25`$ and $`0.75`$ the calculations were performed only for the $`L1_2`$ structure, while for $`x=0.50`$ we considered three different structures: $`L1_0,B2`$ and $`B32`$.
In Fig. 5 the charge density contour plots are shown for pure fcc Al, and Al<sub>3</sub>Li in $`L1_0`$ structure, as calculated from Eq. (44) using $`s,p,d`$ basis set, and the maximum orbital quantum number $`l^{}`$ included in (12) was $`10`$. The agreement between these plots and those from Ref. is very good.
The calculated equilibrium Wigner-Seitz radii and the bulk moduli are tabulated in Table I and plotted in Figs. 6 and 7. In the case of pure Al and Li the structure energy differences and in the case of compounds the heats of formation are included in the Table and plotted in Fig. 8. The heat of formation is defined as
$$\mathrm{\Delta }HE_{Li_xAl_{1x}}xE_{Li}(1x)E_{Al},$$
(66)
where all the energies are obtained for the proper equilibrium volume and they are expressed per atom. In Figs. 5-7 and Table I the full-potential values from Ref. are also included.
The mean deviations between the present and the full-potential results for the equilibrium radii, bulk moduli and heats of formations are $`9.8\%,7.5\%`$ and $`17\%`$, respectively. Taking into account the minor discrepancies between the numerical details used in the calculations, for example the exchange-correlation functional, the way the core electrons were treated etc., and the fact that the full-potential methods have their own error limits as well, we can conclude that the agreement between the two sets of results is very good. One should appreciate how well the trends obtained in the full-potential calculation are reproduced by the present method.
## V Conclusions
We have presented a self-consistent implementation based on the Green’s function technique of the Exact Muffin-Tin Orbitals Theory, developed by O.K. Andersen et al. The accuracy of the present implementation was tested on different systems, where we have found a good agreement between the present results and the results obtained by full-potential techniques. In order to gain some experience about the efficiency of the present method we compare the CPU times of a self-consistent calculation of the tight-binding ASA-LMTO method , based also on the Green’s function technique, and that of the SCA-EMTO method. We found that the present implementation of the SCA-EMTO method needs with about $`3`$ times larger CPU time than the tight-binding ASA-LMTO method.
Finally we remark that if the radii of the potential spheres are chosen to be equal with the radii of the atomic spheres, i.e. $`w_r=s_R`$, the SCA-EMTO method can be considered as an ASA based Green’s function technique that involves the so called combined correction term . It gives exact one electron energies and charge densities for the optimized overlapping muffin tin wells. The natural extension of the present SCA-EMTO method to compute the total energies from the output total charge density via the Full Charge Density technique is in progress.
## ACKNOWLEDGMENTS
L.V. acknowledges the interesting and helpful discussions with Prof. O. K. Andersen, the assistance from Drs. C. Arcangeli and R. W. Tank, and the hospitality of the Max-Planck Institute from Stuttgart where the first part of this work was performed. Thanks are also due to Dr. A. V. Ruban for his valuable observations. The Swedish Natural Science Research Council, the Swedish Foundation for Strategic Research and Royal Swedish Academy of Sciences are acknowledged for financial support. Center for Atomic-scale Materials Physics is sponsored by the Danish National Research Foundation. Part of this work was supported by the research project OTKA 023390 of the Hungarian Scientific Research Fund.
## VI Appendix
During the self-consistent procedure the Eq. (26), (42), (44) and (63) are solved iteratively. In order to construct the electron density (44), as the input for the next iteration, we have to invert the kink matrix for each complex energy $`z`$ along the contour and for each Bloch vector $`𝐤`$ from the Brillouin zone. For a reasonably high accuracy we need at least a few hundreds of Bloch vectors in the irreducible part of the Brillouin zone, therefore this solution means the most time-consuming step of the self-consistent procedure. Here we apply a similar "two-step" scheme introduced in Ref. in order to reduce the number of time-consuming iterations. Within this scheme after each iteration an approximate charge self-consistency is achieved by solving self-consistently the following equation written for the $`𝐤`$-integrated path operator
$$g^a(z)=\left[1+g^{a\mathrm{\hspace{0.33em}0}}(z)\left(D^{a\mathrm{\hspace{0.33em}0}}(z)D^a(z)\right)\right]^1g^{a\mathrm{\hspace{0.33em}0}}(z),$$
(67)
where where the index $`0`$ denotes quantities obtained from the previous iteration, and which are kept fixed during the solution of Eq. (67).
In the expression for the number of state (42), we need the $`𝐤`$-integrated trace of the product between the path operator and the energy derivative of the kink matrix, therefore, a similar equation to (67) has to be established for the quantity
$$G_{R^{}L^{}RL}^a(z)_{BZ}g_{R^{}L^{}RL}^a(z,𝐤)\dot{K}_{RLR^{}L^{}}^a(z,𝐤)𝑑𝐤.$$
(68)
Using the definition of the kink matrix after some manipulations we arrive to the equation
$$G_{R^{}L^{}RL}^a(z)=g_{R^{}L^{}RL}^a(z)\left[\frac{G_{R^{}L^{}RL}^{a\mathrm{\hspace{0.33em}0}}(z)}{g_{R^{}L^{}RL}^{a\mathrm{\hspace{0.33em}0}}(z)}+\delta _{R^{}R}\delta _{L^{}L}a_R\left(\dot{D}_{Rl}^{a\mathrm{\hspace{0.33em}0}}(z)\dot{D}_{Rl}^a(z)\right)\right].$$
(69)
It is worth to note that in this expression we do not have matrix multiplication. Finally we mention that as soon as the self-consistency is achieved, i.e.
$$D_{Rl}^a(z)D_{Rl}^{a\mathrm{\hspace{0.33em}0}}(z)\text{and}\dot{D}_{Rl}^a(z)\dot{D}_{Rl}^{a\mathrm{\hspace{0.33em}0}}(z),$$
(70)
both equations, (67) and (69), become exact.
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# Nonlinear acoustic and microwave absorption in disordered semiconductors
## I Introduction
The subject of this article is nonlinear microwave and acoustic properties of amorphous semiconductors and lightly doped crystalline semiconductors in the regime of hopping conductance. We are interested in absorption to electron transitions between localized states associated with defects or impurity atoms. We consider the case where the absorption is due to electron hopping within pairs of neighboring defects containing one electron per pair. The distance between the centers within the pair must be small enough to allow tunneling, while the distance to other impurities should be large enough to prevent tunneling to impurities outside the pair. In a weakly doped semiconductor we can expect these pairs to be relatively rare, and triplets of the same kind will be even less likely. Thus a natural basis to treat the problem is the so-called two level approximation according to which only the lowest energy level of each of two neighboring impurities are taken in consideration. This approach, as well as its range of applicability, was first discussed in detail by Pollack and Geballe.
For brevity, in the following we shall discuss the case of acoustic attenuation and then specify what changes in the formulae should be introduced to allow for electromagnetic absorption.
An external AC electric or acoustic field causes transitions between the electron states. Direct inter-level transitions leading to absorption of quanta give rise to the so-called *resonant* absorption. For low intensities, the resonant contribution to the absorption coefficient of an acoustic wave can be expressed as
$$\mathrm{\Gamma }^{(\text{res})}=\alpha _1\omega /s\mathrm{tanh}\left(\mathrm{}\omega /2k_BT\right)$$
(1)
where $`\alpha _1`$ is a dimensionless coupling parameter, weakly dependent on temperature and frequency, $`s`$ is the sound velocity, $`\omega `$ is the sound frequency, and $`T`$ is the temperature. We can always assume the relation $`\mathrm{}\omega /2k_BT1`$. $`\alpha _1`$ will be specified later, see Eq. (20). Defining $`n_0`$ as the Fermi function $`n_0(E)=[\mathrm{exp}(E/k_BT)+1]^1`$, we can write the factor $`\mathrm{tanh}\left(\mathrm{}\omega /2k_BT\right)`$ as $`n_0(\mathrm{}\omega /2)n_0(\mathrm{}\omega /2)`$. This can be recognized as the difference in the equilibrium population of the two levels of the pair with an energy splitting of $`\mathrm{}\omega /2(\mathrm{}\omega /2)=\mathrm{}\omega `$, that is, the pairs which can directly absorb a phonon.
The *relaxation* absorption is due to a modulation of the electron inter-level spacing $`2ϵ`$ by the AC field. Such a modulation leads to a periodic change of the occupation numbers of the two levels which lags in phase the variation of $`ϵ`$. This lag leads to the energy dissipation. In the linear regime the coefficient of relaxation absorption has been calculated as
$$\mathrm{\Gamma }^{\text{(rel)}}\frac{\alpha _2}{s}\{\begin{array}{cc}\tau ^1,\hfill & \omega \tau _01\hfill \\ \omega ,\hfill & \omega \tau _01\hfill \end{array}$$
(2)
where $`\alpha _2\alpha _1`$ and $`\tau _0`$ represents a minimal relaxation time for $`ϵk_BT`$. The physical meaning of $`\tau _0`$, as well as estimates of this quantity, will be will be discussed later.
Comparing Eq. (2) with Eq. (1) we conclude that the relaxation absorption always predominates at $`\omega \tau _01`$. If $`\omega \tau _01`$ the ratio $`\mathrm{\Gamma }^{\text{(res)}}/\mathrm{\Gamma }^{\text{(rel)}}\omega \tau _0\mathrm{tanh}(\mathrm{}\omega /2k_BT)`$ can be either greater or less than one under experimentally accessible conditions.
For higher intensities this comparison is no longer valid, as both the resonant and the relaxation absorption show strongly nonlinear behavior. For the resonant absorption, the nonlinearity is due to an equalization of the population numbers of the two electron states. From Eq. (1) we see that this leads to a strong reduction of the resonant absorption. Usually this suppression takes place at very low intensities.
The nonlinearity of the relaxation absorption is due to the following. In the limit of high intensities, when the perturbing potential is amplitude $`k_BT`$, there will be times when $`2ϵ(t)k_BT`$ and both states are almost equally occupied. During such part of the wave period no transitions and thereby no absorption will occur. Consequently, the absorption coefficient decreases with the sound amplitude.
The linear relaxation absorption in semiconductors has previously been studied both theoretically and experimentally (for a review, see Ref. ), while for the strongly nonlinear regime only theoretical results were obtained. The reason why this type of nonlinearity has not been observed experimentally is probably some masking by other mechanisms leading to a nonlinear behavior. One of such mechanisms could be wave-induced ionization of impurity atoms into the conduction band observed in InSb. Consequently, to observe the mechanism of nonlinear absorption one should carefully choose the material. That does not seem to be an easy task, and we are not aware of any experiments of this type.
In this article we will address the case of *low* intensities when nonlinear effect manifest themselves as small corrections to the linear absorption. The main message is that already the lowest order corrections offer interesting information and anomalous effects worth studying. Hopefully, these effects will be pronounced within an experimentally achievable parameter space where other mechanisms of nonlinear behavior are still not important.
The two-site approximation for semiconductors allows us to describe our system in terms of the so-called Two Level System (TLS) model. This model was first proposed independently by Anderson et al. and Philips to explain the low temperature specific heat in amorphous materials, and has been successful in describing several other phenomena. The model has given a theoretical explanation for a surprising universality in the behavior of very different disordered materials at low temperatures.
The authors have previously used the TLS model to analyze nonlinear corrections to the absorption in both dielectric and metallic glasses. However, there are several important differences between the situations in glasses and in disordered or doped semiconductors. The main difference is due to the electric charge of the particle involved in the transitions. This introduces two modifications to the earlier results.
First, the electric charge gives the pair a dipole moment proportional to the distance between the pairs, the coupling being proportional to the dipole moment’s component along the direction of the electric field. This leads to a specific orientational dependence of the absorption of one pair, and after integration over all pairs it significantly influences the absorption. Second, we must expect our system to behave differently with the application of a magnetic field. More specifically, as the magnetic field leads to a stronger localization of the electrons, we must expect the absorption to decrease with an applied field. The effect of a magnetic field will therefore also be analyzed, in the special cases of a weak or a strong field. The connection between the orientation of the dipole moment and the magnetic field direction also leads to a dependence of the absorption coefficient on the relative directions of the radiation wave vector and the magnetic field. We will solve the problem for a magnetic field parallel or perpendicular to the wave vector of the radiation.
The magnetic field dependence gives us the possibility to separate the relaxational absorption from other contributions. It will also be shown that the effect of a magnetic field depends on the specific relaxation mechanism, thus providing us with a tool to further understand the relaxation processes for localized carriers.
The paper is organized as follows. First, we give a short introduction to the theory based on the TLS model and how it can be used to solve the problem of relaxation absorption in glasses. Then we will fit the two-site approximation for a semiconductor to the TLS model and show what modifications are needed for this. Finally, we will show how the effects of a magnetic field can readily be included in the analysis, and what effects to expect from this.
## II Two Level Systems and relaxation losses in glasses
The TLS model deals with a particle (in our case, an electron) moving in a slightly asymmetric double well potential. It is assumed that only the two ground levels are accessible. The ground levels in the isolated wells are assumed to have a slight separation in energy, $`2\mathrm{\Delta }`$, and to be coupled via a tunneling energy overlap integral $`\mathrm{\Lambda }`$. The Hamiltonian of such a system is traditionally written as
$$_0=\mathrm{\Delta }\sigma _z\mathrm{\Lambda }\sigma _x$$
(3)
where $`\sigma _i`$ are the Pauli matrices.
Let us now apply an external periodic perturbing potential and study the power absorbed by a single pair, $`p(\mathrm{\Delta },\mathrm{\Lambda })`$. As is will be shown, this power is in general a non-monotonous function of its parameters, and there exists an “optimal” region which dominates absorption. What is important is that those regions are *different* for the linear absorption and for the nonlinear correction. Let us define $`\delta `$ as the typical value of the inter-level splitting $`\mathrm{\Delta }`$ which is important for the onset of nonlinear behavior. Estimates for the quantity $`\delta `$ will be given in the discussion section.
At $`\mathrm{}\omega \delta `$ one can employ the adiabatic approximation and neglect time derivatives of the external field while solving the Schrödinger equation for the TLS. In this approximation we write the interaction Hamiltonian as
$$_{}=\sigma _zd\mathrm{cos}\omega t$$
(4)
ignoring possible off-diagonal items. The quantity $`d`$ is just the coupling constant between the field and the TLS. In the physics of low-temperature properties of glasses $`d`$ is assumed to be a random quantity, uncorrelated with $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }`$. This assumption is generally not valid for the case of semiconductors.
In the case of a sound wave $`d=\gamma _{ik}u_{ik}^{(0)}`$, where $`\gamma _{ik}`$ is the deformational potential of the TLS and $`u_{ik}^{(0)}`$ is the amplitude value of the deformation tensor.
For semiconductors the value of $`\mathrm{\Lambda }`$ depends on the spatial separation of the wells. In the case of electromagnetic waves $`d=\eta _0`$, where $`\eta `$ is the dipole moment of the TLS while $`_0`$ is the amplitude of the electric field. For charged particles the dipole moment is proportional to the distance between the wells, and is thereby strongly correlated with the value of $`\mathrm{\Lambda }`$.
The total Hamiltonian of the TLS may now be written as
$$=(\mathrm{\Delta }+d\mathrm{cos}\omega t)\sigma _z\mathrm{\Lambda }\sigma _x.$$
(5)
The difference between the eigenvalues of this new Hamiltonian, $`2ϵ(t)`$, where
$$ϵ(t)=\sqrt{(\mathrm{\Delta }+d\mathrm{cos}\omega t)^2+\mathrm{\Lambda }^2}.$$
(6)
is just the energy splitting of the TLS. To characterize a TLS we need not only the energy spacing, but also the occupation numbers of the upper ($`n`$) and lower ($`1n`$) levels. The non-equilibrium occupation numbers can be found from the balance equation
$$\frac{dn}{dt}=\frac{nn_0(t)}{\tau (t)}$$
(7)
where $`n_0(t)`$ is the adiabatic equilibrium occupation number. It depends on the energy spacing $`ϵ(t)`$, temperature $`T`$ and time $`t`$ as
$$n_0(t)=\left[e^{2ϵ(t)/k_BT}+1\right]^1.$$
(8)
The relaxation time $`\tau (t)`$ is a function of the energy splitting, the tunneling barrier and temperature, and also depends on the exact relaxation mechanism.
The power absorbed by a single TLS can be determined by the expression
$$p(\mathrm{\Delta },\mathrm{\Lambda })=\frac{2}{\mathrm{\Theta }}_0^\mathrm{\Theta }𝑑tn(t)\frac{dϵ}{dt},\mathrm{\Theta }\frac{2\pi }{\omega }.$$
(9)
The contributions of individual TLS must be added and such a summation can be performed in a conventional way using the distribution function $`N(\mathrm{\Delta },\mathrm{\Lambda })`$ of the random parameters $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }`$ and replacing the deformational potential $`\gamma _{ik}`$ by its average value.
To analyze the nonlinear absorption we use the exact periodic in time solution of Eq. (7) to obtain the following result for the total absorbed power,
$`P`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{\Theta }}}{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}{\displaystyle \frac{N(\mathrm{\Delta },\mathrm{\Lambda })d\mathrm{\Delta }d\mathrm{\Lambda }}{k_BT}}\left(1e^{_0^\mathrm{\Theta }𝑑t_1/\tau (t_1)}\right)^1`$ (10)
$`\times `$ $`{\displaystyle _0^\mathrm{\Theta }}{\displaystyle _0^\mathrm{\Theta }}{\displaystyle \frac{dtdt^{}\dot{ϵ}(t)\dot{ϵ}(tt^{})}{\mathrm{cosh}^2[ϵ(tt^{})/k_BT]}}e^{_0^t^{}𝑑t_1/\tau (tt_1)}.`$ (11)
Unfortunately, this integral cannot be calculated analytically in the general case. Earlier discussions of the strongly nonlinear regime have considered situations where only very limited time or energy ranges contribute to the absorption. This corresponds to the asymptotic high intensity limits of the integral. Our approach has rather been expanding the integral in powers of the amplitude of the modulation of the energy splitting, $`d`$, to get insights into the first onset of nonlinear effects. So we will concentrate on the regime of weak nonlinearity.
## III Semiconductors in the TLS formalism
We will now show how hopping in amorphous or doped crystalline semiconductors can be described by the TLS model.
A TLS is naturally formed by a pair of nearest impurity centers having one electron per pair. To make the calculations for an individual pair we have to specify the coupling constant $`d`$ and the relaxation time $`\tau `$. The latter is essentially dependent on the dominant mechanism of pair population relaxation. To sum over all the pairs we have to specify the proper distribution function for the parameters of the pairs containing one electron.
We will concentrate on the case when the interaction has a dipole character and can be expressed in the form
$$d(𝐫,t)=d_0(𝐫)\mathrm{cos}\omega t,d_0(𝐫)e_0r(𝝂𝐧)$$
(12)
where $`_0`$ is the amplitude of the effective electric field acting upon the electrons. It is just the amplitude of the local electric field created either by the external electric field, our due to piezoelectric interaction with an acoustic wave. $`𝝂=𝓔/`$ is the field polarization vector, $`r`$ is the distance between the components of the pair, $`𝐧=𝐫/r`$ is the pair direction vector.
The relaxation time is strongly dependent on the particular mechanism of interaction between localized electrons and phonons, see for example review Ref. . Here we will only quote the most important results.
### A Relaxation time
In general, the relaxation time can be expressed as
$`{\displaystyle \frac{1}{\tau (ϵ,𝐫)}}`$ $`=`$ $`{\displaystyle \frac{1}{\tau _n(T)}}\left({\displaystyle \frac{ϵ}{k_BT}}\right)^n\left({\displaystyle \frac{\mathrm{\Lambda }(𝐫)}{ϵ}}\right)^2`$ (14)
$`\times \mathrm{\Phi }_n\left({\displaystyle \frac{2ϵ}{E_r}}\right){\displaystyle \frac{\mathrm{coth}(ϵ/k_BT)}{[1+(2ϵ/E_a)^2]^4}},`$
where the exponent $`n`$, the factor $`\tau _n^1`$ and the function $`\mathrm{\Phi }_n(x)`$ are dependent on the particular mechanism of interaction between localized electrons and phonons. The meaning of the energies $`E_r`$ and $`E_a`$ will be made clear in a moment.
The energy dependence of $`\tau ^1`$ can be easily understood. The phonon emitted or absorbed during a transition between the electron levels must obviously have an energy equal to the energy splitting of the pair, $`2ϵ`$. The power by which it occurs in the formula (14) is determined by the product of the phonon density of states with the frequency dependence of the squared interaction matrix element. The factor $`\mathrm{coth}(ϵ/k_BT)`$ is equal to $`2N_\omega +1`$ where $`N_\omega `$ is the Planck function for $`\mathrm{}\omega =2ϵ`$. Apart from a proportionality factor, this is the probability of phonon emission, $`N_\omega `$, plus the probability of absorption, $`(N_\omega +1)`$. The factor $`[\mathrm{\Lambda }(𝐫)/ϵ]^2`$ is a dimensionless measure of the tunneling coupling between the bare states in the two wells. It can be seen that $`\tau `$ has a minimum with respect to $`\mathrm{\Lambda }`$ when $`\mathrm{\Lambda }=ϵ`$, which is equivalent to $`\mathrm{\Delta }=0`$. This condition defines the minimal relaxation time $`\tau _0`$ referred to in Eq. (2). The minimal $`\tau `$ corresponds to the symmetrical configuration when the bare energy levels at both sites are equal. To allow for time dependence of $`\tau `$ one should substitute $`ϵ=ϵ(t)`$ from Eq. (6).
The expression (14) contains two specific energy scales, $`E_a2\mathrm{}s/a`$, and $`E_r\mathrm{}s/r`$. The first scale is the energy of a phonon having a wavelength of the order of the single-site localization length, $`a`$. Phonons with larger energies produce rapidly oscillating fields which average out at the distance occupied by a localized electron. Consequently, $`\tau ^1`$ strongly decays at $`ϵE_a`$. The second scale corresponds to phonons with a wavelength of the order of the distance $`r`$ between the components of the pair. If the deformation potentials of both components of the pair are the same, or if the main mechanism of the electron-phonon interaction is piezoelectric, then at $`2ϵE_r`$ both levels move synchronously, and no interaction occurs. The net interaction is hence proportional to some power of the ratio $`2ϵ/E_r`$, see Ref. for a review.
In most realistic cases one can assume $`ϵE_a`$. However, for $`T1`$ K the ratio $`x=2ϵ/E_r`$ can be either less or greater than 1, giving different types of behavior for $`\mathrm{\Phi }_n(x)`$. We will concentrate on the case of $`x1`$ since this limiting case seems to be more easily accessible for experiments. In this regime, $`\mathrm{\Phi }_n(x)`$ can be considered as constant. The validity of this approximation will be considered in the Discussion. The quantities $`\tau _n(T)`$ are listed in Ref. .
The most important feature of the relaxation for our problem is the *energy dependence* of the relaxation rate, namely, the power $`n`$. Under the above-mentioned conditions, $`n=3`$ in the case of deformational interaction and $`n=1`$ for piezoelectric interactions. As will be clear later, only the energies $`ϵk_BT`$ are important for the anomalous nonlinear behavior, so one can approximate $`\mathrm{coth}(ϵ/k_BT)k_BT/ϵ`$. In this way we arrive at the following energy dependences of the relaxation rate:
$$\frac{1}{\tau }\{\begin{array}{cc}\mathrm{\Lambda }^2ϵ^0\hfill & \text{for the deformational interaction}\hfill \\ \mathrm{\Lambda }^2ϵ^2\hfill & \text{for the piezoelectric interaction}\hfill \end{array}$$
(15)
Note that in the first case the relaxation time is independent on $`ϵ`$ and thereby on time. This is the same as is the case in dielectric glasses. Apart from a constant and the dependency on the magnetic field, we can thus expect the same type of behavior from these two very different systems. In the second case the $`ϵ`$-dependence is the same as in metallic glasses, which has also been analyzed by the authors and has proven to be the source of a pronounced anomalous effect. In particular, in metallic glasses the lowest nonlinear contribution is proportional to the intensity to $`3/2`$ rather than to the intensity squared as for dielectric glasses.
In the absence of the magnetic field the energy overlap integral $`\mathrm{\Lambda }`$ is related to the distance $`r`$ between the sites of a pair simply as $`\mathrm{\Lambda }=\mathrm{\Lambda }_0\mathrm{exp}^{r/a}`$, where $`\mathrm{\Lambda }_0=(15)\times me^4/\mathrm{}^2\kappa ^2`$ is of the order of the effective Bohr energy. Here $`m`$ is the electron effective mass while $`\kappa `$ is the dielectric constant. A magnetic field will squeeze the electron wave function, and this effect will be strongest for the direction perpendicular to the field. This introduces an angular dependency to the localization length and thereby also to $`\mathrm{\Lambda }`$. Following Ref. , we will analyze the limiting cases of weak (w) and strong (s) magnetic field where the influence of magnetic field is weak or strong, respectively. The asymptotic expressions for the $`\zeta (𝐫)\mathrm{ln}[\mathrm{\Lambda }(𝐫)/\mathrm{\Lambda }_0]`$ are the following:
$`\zeta _w`$ $`=`$ $`{\displaystyle \frac{r}{a}}+{\displaystyle \frac{r^3a\mathrm{sin}^2\theta }{24\lambda ^4}},`$ (16)
$`\zeta _s`$ $`=`$ $`{\displaystyle \frac{r^2\mathrm{sin}^2\theta }{4\lambda ^2}}+{\displaystyle \frac{|r\mathrm{cos}\theta |}{a_H}}.`$ (17)
Here $`\theta `$ is the angle between $`𝐫`$ and the direction of the magnetic field $`𝐇`$, $`\lambda =\sqrt{\mathrm{}c/eH}`$ is the magnetic length, while $`a_H=\mathrm{}/\sqrt{2mE_H}`$ is the characteristic localization length in the longitudinal direction, where $`E_H`$ is the ionization energy of the ground state of the localized electron in the magnetic field.
### B Pair distribution function
The total absorption is given by a sum of the contributions of the individual pairs. Hence we have to sum over the $`𝐫`$, as well as over the individual energies of the electron levels. The latter summation must take into account the correlation between the level occupation numbers due to Coulomb interaction. As shown in Ref. , the summation over the energies can be split into integration over the pair center-of-gravity and over the bare inter-level spacing $`\mathrm{\Delta }`$. The first integration gives $`2\mathrm{\Delta }+e^2/\kappa r`$ since only the pairs with the center-of-gravity energy between the chemical potential $`\mu +\mathrm{\Delta }`$ and $`\mu \mathrm{\Delta }e^2/\kappa r`$ have one electron per pair. As a result, the pair distribution function can be expressed through the single-electron density of states $`g`$ as,
$$N(\mathrm{\Delta },𝐫)=\frac{g^2V}{4\pi }\left(2\mathrm{\Delta }+\frac{e^2}{\kappa r}\right)$$
(18)
where $`V`$ is the volume contributing to the absorption. This expression is valid if for a typical hopping distance $`e^2/\kappa r\mathrm{\Delta }_C`$ where $`\mathrm{\Delta }_C`$ in the width of the *Coulomb gap* in the single-electron density of states. Inside the Coulomb gap, we would rather have to use the distribution
$$N(\mathrm{\Delta },𝐫)=\frac{3}{40\pi ^3}\left(\frac{\kappa }{e^2}\right)^6\left(2\mathrm{\Delta }+\frac{e^2}{\kappa r}\right)^5.$$
(19)
Calculations for both cases are similar. Note that the distribution function is isotropic. Under the conditions of interest to us the typical hopping distance $`r`$ is small enough to let us neglect $`2\mathrm{\Delta }`$ in comparison with $`e^2/\kappa r`$ in Eqs. (18) and (19). Thus the distribution becomes $`\mathrm{\Delta }`$-independent, and we denote it as $`N(𝐫)`$. Using the notation presented above we can now write the coupling constants $`\alpha _1,alpha_2\alpha `$ from Eqs. 1, and 2 as
$$\alpha =\frac{4\pi ^3}{3}𝒦^2\frac{e^4g^2ar_\omega ^3}{\kappa ^2},r_\omega =a\mathrm{ln}\frac{\mathrm{\Lambda }_0}{\mathrm{}\omega },$$
(20)
where $`𝒦`$ is the coupling constant of the piezoelectric interaction. The power of $`r_\omega `$ may vary for different type of interactions, depending on whether the interaction includes the dipole moment of the pair.
## IV Calculation of absorption
As a result of the previous considerations, the absorbed power can be expressed as
$$P=𝑑𝐧_0^{\mathrm{}}𝑑\mathrm{\Delta }_0^{\mathrm{}}r^2𝑑rN(r)\left\{_0^\mathrm{\Theta }_0^\mathrm{\Theta }\frac{dtdt^{}}{\mathrm{\Theta }k_BT}\frac{\dot{ϵ}(𝐫,t)\dot{ϵ}(𝐫,tt^{})}{\mathrm{cosh}^2[ϵ(𝐫,tt^{})/k_BT]}\frac{\mathrm{exp}\left(_0^t^{}𝑑t_1/\tau (tt_1)\right)}{1\mathrm{exp}[_0^\mathrm{\Theta }𝑑t_1/\tau (t_1)]}\right\}.$$
(21)
Here the expression in the braces is just the power absorbed by an individual pair, $`p(\mathrm{\Delta },𝐫)`$. The energy splitting $`2ϵ`$ depends on $`𝐫`$ through the interaction potential $`d_0(𝐫)`$ given by Eq. (12) and through the tunneling splitting $`\mathrm{\Lambda }=\mathrm{\Lambda }_0e^{\zeta (𝐫)}`$. Technically it is convenient to transform the integral from the set of variables $`r,𝐧`$, to the variables $`\mathrm{\Lambda },\theta ,\varphi `$. Such a transform introduces the factor
$$\left|\frac{\mathrm{\Lambda }(r,𝐧)}{r}\right|^1=\frac{a}{\mathrm{\Lambda }}f(\mathrm{\Lambda },𝐧).$$
Here the dimensionless function $`f(\mathrm{\Lambda },𝐧)`$ is given by the equation
$$f(\mathrm{\Lambda },𝐧)\frac{1}{a}\frac{r_\mathrm{\Lambda }(𝐧)}{\mathrm{ln}\mathrm{\Lambda }}=\frac{1}{a}\frac{r_\mathrm{\Lambda }(𝐧)}{}.$$
(22)
where $`\mathrm{ln}(\mathrm{\Lambda }_0/\mathrm{\Lambda })`$, while $`r_\mathrm{\Lambda }(𝐧)`$ is the solution of the equation
$$\zeta (r_\lambda ,𝐧)=.$$
(23)
In the simplest case of zero magnetic field $`f=1`$, and we have a distribution very similar to those of glasses. The variable transform strongly simplifies the calculations since the relaxation time is a function of the parameters $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }`$ and the quantity $`r_\mathrm{\Lambda }`$ is a weak (logarithmic) function of $`\mathrm{\Lambda }`$. Consequently it can be extracted out of the integral over $`\mathrm{\Lambda }`$, while replacing $`\mathrm{\Lambda }`$ in the expression for $`r_\mathrm{\Lambda }`$ by its characteristic value. As a result, the integrations over $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }`$ will remain the same as previously calculated for glasses, and only the angular integration and the dependence on the characteristic value of $`\mathrm{\Lambda }`$ are different.
The following calculation procedure is similar to that of Ref. . The expression (21) will be expanded in powers of the effective electric field, and the lowest correction will be compared with the linear result. Expanding the individual contributions in powers of $`d_0`$ as
$$p(\mathrm{\Delta },𝐫)=\underset{k=2}{\overset{4}{}}p^{(k)}d_0^k(𝐫)$$
(24)
we notice the the coefficients $`p^{(k)}`$ depend only on the quantities $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }`$. Transforming the variables from $`r,𝐧`$ to $`\mathrm{\Lambda },𝐧`$ we can use the fact that $`r`$ is a weak function of $`\mathrm{\Lambda }`$ and extract of the quantities proportional to the powers of $`r`$ from the integral over $`\mathrm{\Lambda }`$ replacing
$$rr_c(𝐧)=r_\mathrm{\Lambda }(𝐧)|_{\mathrm{\Lambda }=\mathrm{\Lambda }_c}$$
where $`\mathrm{\Lambda }_c`$ is the characteristic value determined by the integrand over $`\mathrm{\Lambda }`$. In a similar way, we replace $`f(\mathrm{\Lambda },𝐧)f_c(𝐧)=f(\mathrm{\Lambda }_c,𝐧)`$. Finally we arrive at the expression $`P=_{k=2}^4P^{(k)}`$ with $`P^{(k)}=(ae_0)^kI_kJ_k`$ where
$`I_k`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}𝑑\mathrm{\Delta }{\displaystyle _0^{\mathrm{}}}𝑑\mathrm{\Lambda }\mathrm{\Lambda }^1p^{(k)}(\mathrm{\Delta },\mathrm{\Lambda }),`$ (25)
$`J_k`$ $`=`$ $`{\displaystyle 𝑑𝐧(𝝂𝐧)^kf_cN(r_c)(r_c/a)^k}`$ (26)
The quantities $`I_k`$ are the same that enter the expressions for nonlinear absorption in glasses, and we quote them from Ref. .
Parallel with $`I_k`$ we can also estimate $`\mathrm{\Lambda }_c`$, and thereby $``$. It can be shown that this value is only weakly dependent on the mechanism of absorption and on $`k`$. For all the mechanisms analyzed in this paper we have
$$=\mathrm{ln}\frac{c_{}\mathrm{\Lambda }_0}{k_BT\sqrt{\omega \tau _n(T)}}.$$
(27)
The value for $`c_{}`$ is of the order one, and varies only with a factor less than two for the situations discussed in this paper. We will therefore treat $``$ as independent of $`k`$ for this discussion.
In this paper we will present the results for the case of low frequencies, $`\omega \tau _n1`$, for which the anomalous nonlinear behavior is most pronounced. In this case the linear in the intensity contribution, $`I_0`$, is independent of $`\tau `$, and moreover, it is the same for any dependence $`\tau (ϵ)`$,
$$I_0=(\pi ^2/16)\omega 0.62\omega .$$
(28)
The case of deformational interaction, for which the parameter $`n`$ in Eq. (14) is equal to 3, corresponds to the situation in dielectric glasses. In this case
$$I_3=0,I_40.054\frac{\omega }{\sqrt{\omega \tau _3(T)}(k_BT)^2}.$$
(29)
The piezoelectric interaction ($`n=1`$) is similar to the case of metallic glasses. It can be shown that the energy dependence of the relaxation time leads to a divergence in the integration over $`\mathrm{\Lambda }`$ and $`\mathrm{\Delta }`$. To get a proper estimate one should cut off the integration at $`ϵd_0`$. As a result, the leading nonlinear contribution appears proportional to $`|d_0|^3`$ and equal to
$$I_30.1\omega /k_BT.$$
(30)
We again expect $``$ to be given by Eq. 27. To calculate the angular integrals $`J_k`$ we solve Eq. (23) in a recursive way. For the case of weak magnetic fields, we obtain
$`r_c`$ $`=`$ $`a\left(1^2a^4\mathrm{sin}^2\theta /24\lambda ^4\right);`$ (31)
$`f_c`$ $`=`$ $`1+a^4^2\mathrm{cos}^2\theta /8\lambda ^4.`$ (32)
For the strong field limit, solving a quadratic equation, we obtain
$`r_c`$ $`=`$ $`{\displaystyle \frac{2\lambda ^2\mathrm{cos}\theta }{a_H\mathrm{sin}^2\theta }}\left(\sqrt{1+\mathrm{tan}^2\theta (a_H/\lambda )^2}1\right);`$ (33)
$`f_c`$ $`=`$ $`(a_H/a)\mathrm{cos}\theta \left(1+(a_H/\lambda )^2\mathrm{tan}^2\theta \right)^{1/2}.`$ (34)
Here $`\theta `$ is the angle between $`𝐧`$ and the direction of magnetic field $`𝐇`$. Hence, $`d𝐧=d(\mathrm{cos}\theta )d\varphi `$ where $`\varphi `$ is the azimuthal angle between the projections of $`𝐧`$ and $`𝓔_0`$ on the plane, perpendicular to $`𝐇`$. After substitution of Eqs. (31)–(34) into Eq. (26) the angular integrals are calculated directly.
## V Results
To set the scale of nonlinear corrections let us start with the expression for the linear absorption in the absence of magnetic field, $`P_0P_0(0)`$. Both linear and nonlinear contributions are dependent on the relaxation mechanism of the relevant pairs.
### A Deformational interaction between localized pairs and phonons
We will first consider deformational interaction between the localized pairs and thermal phonons. For the deformational mechanism,
$$P_0=(\pi ^2/48)(Va^4g^2e^4\omega ^3_0^2/\kappa ).$$
(35)
In the absence of magnetic field we obtain
$$P_4(0)=\frac{P(0)P_0(0)}{P_0(0)}=0.26\frac{F^2_c^2}{\sqrt{\omega \tau _0(T)}}$$
(36)
were we have introduced the dimensionless “field amplitude”
$$Fe_0a/k_BT.$$
(37)
In a weak magnetic field, the quadratic in magnetic field corrections arise both to the linear absorption and to the lowest nonlinear contribution. They can be expressed in a unified way as, cf. with Ref. ,
$$P_{0/4}^{(w)}(H)=P_{0/4}(0)\left[1c_w(a/\lambda )^4^2\right],$$
(38)
so the magnetic field produces corrections which are $`H^2`$. The numerical factor $`c_w`$ depends on the direction of the electric field $`𝓔_0`$ with respect to the magnetic field $`𝐇`$. Its values are also different for the linear and nonlinear, contributions, $`c_w^{(0)}`$ and $`c_w^{(4)}`$, respectively. The values of $`c_w`$ are shown in Table I.
| Direction | $`c_w^{(0)}`$ | $`c_w^{(4)}`$ | $`c_w^{(3)}`$ |
| --- | --- | --- | --- |
| $`𝐇𝓔_0`$ | 0.2 | $`0.29`$ | $`0.25`$ |
| $`𝐇𝓔_0`$ | 0.1 | $`0.095`$ | $`0.097`$ |
Table I. Numerical coefficients entering the nonlinear contributions to the absorption.
The decrease of attenuation in the magnetic field has the following physical reason. The presence of a magnetic field squeezes the electron wave functions, and the overlap integrals between the components of the pair decrease. Furthermore, the wave functions are squeezed mostly in the direction perpendicular to $`𝐇`$. On the other hand, the coupling between the wave and the pair is maximal if the pair dipole moment is parallel to $`𝓔_0`$. Thus the reduction of absorption is more pronounced for $`𝐇𝐧`$. For the nonlinear contribution the difference should be even stronger, as it includes higher orders of the dipole moment.
In the limit of strong magnetic fields the results are even more interesting, as the functional dependency on the magnetic field also varies with the different absorption types. Yet we still see the same relative considerations as for the weak field limit. For a magnetic field parallel to the radiation polarization vector we get
$`{\displaystyle \frac{P_0^{(s)}(H)}{P_0}}`$ $`=`$ $`3{\displaystyle \frac{\lambda ^2a_H^2}{a^4}}H^{4/3},`$ (39)
$`{\displaystyle \frac{P_4^{(s)}(H)}{P_4(0)}}`$ $`=`$ $`{\displaystyle \frac{5}{2}}{\displaystyle \frac{\lambda ^2a_H^4}{a^6}}H^{5/3}.`$ (40)
For a perpendicular field the results show both a stronger dependency on $`H`$ and on the order of the expansion in intensity.
$`{\displaystyle \frac{P_0^{(s)}(H)}{P_0(0)}}`$ $`=`$ $`6{\displaystyle \frac{\lambda ^4}{^2a^4}}\mathrm{ln}{\displaystyle \frac{a_H^2}{\lambda ^2}}H^2\mathrm{ln}H,`$ (41)
$`{\displaystyle \frac{P_4^{(s)}(H)}{P_4(0)}}`$ $`=`$ $`30{\displaystyle \frac{\lambda ^6}{^3a^6}}\mathrm{ln}{\displaystyle \frac{a_H^2}{\lambda ^2}}H^3\mathrm{ln}H`$ (42)
### B Piezoelectric interaction between localized pairs and phonons
We will now turn our attention to the piezoelectric interaction. This also interacts via a dipole moment, so the linear results are basically the same, apart from a coupling constant. However, there is a striking difference for the nonlinear contribution – similarly to the case of metallic glasses, the integration over $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }`$ in Eq. (25) results in a term proportional to $`|d_0|^3`$ rather than $`d_0^4`$. This is reflected in a change of the absorption dependencies both on the wave intensity and on the magnetic field. For this case, we restrict ourselves by order-of-magnitude estimates for numerical factors. Calculation of exact numbers would require a great amount of numerical work which would be inadequate to the accuracy of the initial model for the electron density of states. The results read as,
$$P_3(0)=[PP(0)]/P_0=c_m|F|,$$
(43)
where the estimate for $`c_m`$ is 0.1. Note that in this case the expansion of the absorption in powers of intensity appears *non-analytical* which implies relatively strong nonlinearity.
Since now the expansion in $`d`$ of the nonlinear contribution starts with $`|d_0|^3`$ rather than from $`d_0^4`$ one can expect weaker magnetic field effects. This is indeed the case. The values of the numerical coefficients $`c_w^{(3)}`$ are shown in the Table I.
For strong fields we get for parallel and perpendicular fields respectively
$`\left({\displaystyle \frac{P_3^{(s)}(H)}{P_3(0)}}\right)_{}`$ $`=`$ $`{\displaystyle \frac{8}{3}}{\displaystyle \frac{\lambda ^2a_H^3}{a^5}}H^{3/2}`$ (44)
$`\left({\displaystyle \frac{P_3^{(s)}(H)}{P_3(0)}}\right)_{}`$ $`=`$ $`{\displaystyle \frac{64}{3}}{\displaystyle \frac{\lambda ^5}{^{5/2}a^5}}\mathrm{ln}{\displaystyle \frac{a_H^2}{\lambda ^2}}H^{5/2}\mathrm{ln}H.`$ (45)
Let us recall that the above expressions are valid for the low frequency limit only, where $`\omega \tau _n(T)<<1`$. Similar calculations are possible for the high frequency limit, as well as for different pair distribution functions. In particular, for the case of pronounced Coulomb gap the essential differences occur only in powers of $`r`$ and thereby of $``$. Thus the influence of the magnetic field will be different. The relation between linear and nonlinear results remain similar, apart from numerical factors.
## VI Discussion
Let us discuss the relevance of the obtained results for realistic materials and situations. In this connection, several parameters are to be considered.
Regarding the material properties, we consider only *weakly doped* or amorphous semiconductors in the regime of *nearest-neighbor hopping conductance*. Hence, we calculate absorption by close pairs independent of each other. To keep the model adequate we have to require that the typical inter-center distance within the pair, $`r_c`$, should be much smaller that the typical distance between defect centers, $`\overline{r}=(4\pi n_d/3)^{1/3}`$. Here $`n_d`$ is the defect concentration. The *hopping distance* $`r_c`$ is discussed in Sec. IV.
Another requirement is that the impurities are not too shallow, so that the electrons cannot be excited from the localized states to the conduction band by the AC perturbing potential. There are experimental examples of this, where such a excitation serves as a source of nonlinear behavior.
According to the present calculation, the most interesting effects occur at “low” frequencies when $`\omega \tau _n(T)1`$. This requirement also ensures that the relaxation absorption dominates the resonant one. Certainly, the minimal relaxation time $`\tau _n(T)`$ is a material property. Usually the above inequality is met at low temperatures for frequencies in the range $`1001000`$ MHz.
The TLS model in glasses is restricted to very low temperatures where higher energy levels are not excited. The situation is a bit different in semiconductor materials where the inter-level splittings are of the order of the Bohr energy. Consequently, the nearest-neighbor hopping conductance can be effective in the temperature range up to a few K.
The main objectives of this paper is to show that nonlinear effects are anomalously large. Indeed, in the case of deformational absorption an additional parameter $`(\omega \tau _3)^{1/2}1`$ is present in the nonlinear expansion (36), while the the case of piezoelectric interaction the nonlinear expansion starts from the *fist power* of dimensionless amplitude $`|F|`$, Eq, (43). Furthermore, the nonlinear contributions have pronounced magnetic field dependences different from the linear ones. We hope that those features will allow experimentalists to detect the nonlinear behavior and to discriminate between different relaxation mechanisms for localized states.
In course of the present calculations we have assumed the inequality $`ϵ/E_r1`$ to be met, see Eq.(14). It is important that for $`ϵ`$ one has to substitute the value which gives the dominating contribution in the final integration over the pair distribution function. This value is actually different for the linear and the nonlinear contributions, and it is also dependent on the relaxation mechanism. It turns out that for the linear absorption this typical $`ϵk_BT`$, while for the nonlinear contributions it is reduced by a factor $`\sqrt{\omega \tau _3}`$ for the deformational interaction, by $`|F|`$ for the piezoelectric one, both calculated under the condition $`ϵE_r`$. A similar estimate is necessary to choose a proper value for $`r`$ in the expression for $`E_r`$. This value depends on the quantities $`\sqrt{\omega \tau _n}`$, $`\mathrm{\Lambda }_0`$, as well as on the magnetic field. Thus the experimental variables intensity, frequency and magnetic field, in addition to the system parameters $`\tau _0`$ and $`\mathrm{\Lambda }_0`$, influence the behavior of the nonlinear absorption. This rich parameter space allows for a large range of experiments.
## VII Conclusions
In this paper, we have analyzed nonlinear contributions to the acoustic and electromagnetic absorption by localized electron states in semiconductors in the regime of hopping conductance. The most important conclusions are the following.
* The total behavior of absorption is determined by the product $`\omega \tau `$ where $`\tau (T)`$ is the minimal relaxation time for a pair with energy splitting of the order $`k_BT`$.
* The anomalous nonlinear behavior occurs at $`\omega \tau 1`$. In the case of deformational relaxation mechanism for the localized electrons a large additional factor $`(\omega \tau )^{1/2}`$ appears in front of the item $`F^2`$ in the expansion of nonlinear absorption. In the case of piezoelectric relaxation mechanism the expansion starts with the item $`|F|`$ rather than $`F^2`$. Here $`F`$ is the dimensional AC field amplitude.
* The anomalous nonlinear absorption is strongly influenced by an external magnetic field, the influence being dependent both on the electron pair distribution function, on the dominating relaxation mechanism for the localized electrons, and on the direction of the field polarization vector with respect to the magnetic field. The influence of magnetic field on the linear absorption and nonlinear corrections is substantially different.
As a result, the physical picture of weakly nonlinear absorption appears rich and informative. Our estimates show that the effects under consideration are accessible for the modern experiment, and many important characteristics – the dominating relaxation mechanism, the importance of the Coulomb gap, typical hopping distances, etc. – can be extracted by comparison to the present theory provided the experiment will be done.
It should be emphasized that there is a close similarity between the present and the results of our previous calculations for glassy materials. However, the localized states in disordered semiconductors, being charged, can be influenced by magnetic field which makes them easier to investigate. We therefore also hope that the studies of semiconductor systems will also provide a new information regarding nonlinear response of TLSs in glasses.
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# The relativistic eikonal approximation in high-energy 𝐴(𝑒,𝑒'𝑝) reactions
## I Introduction
Exclusive A$`(e,e^{}p)`$B reactions from nuclei constitute an invaluable tool to probe a wide variety of nuclear phenomena. At low values of the virtual photon’s four-momentum transfer $`Q^2=\stackrel{}{q}^2\omega ^2`$ and, accordingly, larger distance scales, the quasi-elastic A$`(e,e^{}p)`$ reaction probes the mean-field structure of nuclei. From systematic investigations for a large number of target nuclei a richness of precise information about the independent-particle wave functions and spectroscopic strengths was assembled . At high $`Q^2`$ and decreasing distance scales, the scope of exclusive $`(e,e^{}p)`$ measurements shifts towards studies of (possible) medium dependencies of the nucleonic properties, and, effects like color transparency and the short-range structure of nuclei. Within the context of exclusive $`(e,e^{}p)`$ reactions, “color transparency” stands for the suggestion that at sufficiently high values of $`Q^2`$ the struck proton may interact in an anomalously weak manner with the “spectator” nucleons in the target nucleus .
The extraction of physical information from measured A$`(e,e^{}p)`$B cross sections usually involves some theoretical modelling of which the major ingredients are the initial (bound) and final (scattering) proton wave functions and the electromagnetic electron-nucleus coupling. At lower values of $`Q^2`$, most theoretical work on $`(e,e^{}p)`$ reactions was performed in the so-called distorted-wave impulse approximation (DWIA). The idea behind the DWIA approach is that the inital (bound) and final (scattering) state of the struck nucleon can be computed in a potential model, whereas for the electron-nucleus coupling an “off-shell corrected” electron-proton form can be used. The wealth of high-quality $`(e,e^{}p)`$ data that electron-scattering experiments have provided over the last 20 years, made sure that the DWIA models are well tested against experimental data. For higher values of the energy and momentum transfer ($`Q^21`$(GeV/c)<sup>2</sup>), most theoretical $`(e,e^{}p)`$ work starts from the non-relativistic Glauber theory . This theory is highly successful in describing small angle proton-nucleus scattering at higher energies and is conceived as a baseline for calculating the effect of final-state interactions in high-energy $`(e,e^{}p)`$ reactions. Glauber theory is a multiple-scattering extension of the standard eikonal approximation that relates through a profile function the ejectile’s distorted wave function to the elastic proton scattering wave function . The Glauber method has frequently been shown to be reliable in describing A$`(p,p^{})`$ processes. Several non-relativistic studies have formally investigated the applicability of the Glauber model for describing A$`(e,e^{}p)`$ reactions at higher energies and momentum transfers. Recently, the first high-quality data for exclusive <sup>16</sup>O$`(e,e^{}p)`$ cross sections at higher four-momentum transfer ($`Q^21`$(GeV/c)<sup>2</sup> became available . Below, we will compare results of relativistic eikonal calculations with these data. We believe that this comparison between model calculations and data provides a stringent test of the applicability of the eikonal approximation in describing $`(e,e^{}p)`$ reactions..
Since relativistic effects are expected to become critical in the GeV energy domain, we explore the possibility of developing a fully relativistic model for A$`(e,e^{}p)`$ processes, thereby using the eikonal limit to solve the equations for the final-state wave functions. We employ a relativistic mean-field approximation to the Walecka model to determine the bound state wave functions and binding energies, as well as nucleon and meson potentials. The same mean-field potentials are then also used to compute the scattering wave function in the Dirac eikonal limit. The work presented here is a small initial step towards the formulation of a fully microscopic relativistic model for the description of $`(e,e^{}p)`$ reactions that could possibly bridge the gap between the low and intermediate-energy regime. The model developed in this work can be formally applied in a wide $`Q^2`$ range. As a matter of fact, we employ the relativistic eikonal method to estimate the sensitivity of $`(e,e^{}p)`$ observables in the few GeV regime to a number of physical effects, including off-shell ambiguities and relativity. We adopt different prescriptions for the electron-nucleus coupling in our calculations. By doing this, we estimate the sensitivity of the observables to the theoretical uncertainties that surround the choice of the off-shell electron-proton vertex. It is often claimed that off-shell ambiguities decrease in importance as the four-momentum transfer increases. Here, we make an attempt to quantify the relative importance of the off-shell effects for the $`(e,e^{}p)`$ structure functions by comparing results obtained with different off-shell electron-proton couplings. Hereby we are primarily concerned with the question how big the uncertainties remain when higher and higher four-momentum transfers are probed.
In Sect. II we introduce a relativistic eikonal formalism for calculating A$`(e,e^{}p)`$ observables. This includes a discussion of the method employed to determine the bound (Sect. II B) and scattering (Sect. II C) states. Various forms for the photon-nucleus interaction vertex are introduced in Sect. II D, where special attention is paid to the issue of current conservation. In Sect. III we present the results of our <sup>12</sup>C$`(e,e^{}p)`$ and <sup>16</sup>O$`(e,e^{}p)`$ numerical calculations. In Sect. III B we focus on the issue of the $`Q^2`$ evolution of the off-shell ambiguities. In Sect. III C we compare the results of a fully relativistic calculation with a calculation in which the explicit coupling between the lower components in the inital and final state are neglected. Finally, our concluding remarks are summarized in Sect. IV.
## II Formalism
### A Reaction observables and kinematics.
In this work we follow the conventions for the $`(\stackrel{}{e},e^{}\stackrel{}{p})`$ kinematics and observables introduced by Donnelly and Raskin in Ref. . The four-momenta of the incident and scattered electrons are labeled as $`K^\mu (ϵ,\stackrel{}{k})`$ and $`K^{}_{}{}^{}\mu (ϵ^{},\stackrel{}{k^{}})`$. The electron momenta $`\stackrel{}{k}`$ and $`\stackrel{}{k^{}}`$ define the scattering plane. The four-momentum transfer is given by $`q^\mu =K^\mu K^{}_{}{}^{}\mu =P_{A1}^\mu +P_f^\mu P_A^\mu `$, where $`P_A^\mu `$ and $`P_{A1}^\mu `$ are the four-momenta of the target and residual nucleus, while $`P_f^\mu `$ is the four-momentum of the ejected nucleon. Also, $`q^\mu =(\omega ,\stackrel{}{q})`$, where the three-momentum transfer $`\stackrel{}{q}=\stackrel{}{k}\stackrel{}{k^{}}=\stackrel{}{k}_{A1}+\stackrel{}{k}_f\stackrel{}{k}_A`$ and the energy transfer $`\omega =ϵϵ^{}=E_{A1}+E_fE_A`$ are defined in the standard manner. The $`xyz`$ coordinate system is chosen such that the $`z`$-axis lies along the momentum transfer $`\stackrel{}{q}`$, the $`y`$-axis lies along $`\stackrel{}{k}\times \stackrel{}{k^{}}`$ and the $`x`$-axis lies in the scattering plane; the reaction plane is then defined by $`\stackrel{}{k}_f`$ and $`\stackrel{}{q}`$. The Bjorken-Drell convention for the gamma matrices and Dirac spinors is followed, so that the normalization condition for Dirac plane waves, characterized by a four-momentum $`K^\mu `$ and spin-state $`S^\mu `$, is $`\overline{u}(K^\mu ,S^\mu )u(K^\mu ,S^\mu )=1`$.
In the one-photon-exchange approximation, the process in which a longitudinally polarized electron with helicity h, impinges on a nucleus and induces the knockout of a single nucleon, leaving the residual nucleus in a certain discrete state, can be written in the following form :
$`{\displaystyle \frac{d^5\sigma }{dϵ^{}d\mathrm{\Omega }_e^{}d\mathrm{\Omega }_x}}`$ $`=`$ $`{\displaystyle \frac{MM_{A1}k_f}{8\pi ^3M_A}}f_{rec}^1\sigma _M[(v_L_L+v_T_T+v_{TT}_{TT}+v_{TL}_{TL})`$ (1)
$`+`$ $`h(v_T^{}_T^{}+v_{TL^{}}_{TL^{}})],`$ (2)
where $`f_{rec}`$ is the hadronic recoil factor
$`f_{rec}={\displaystyle \frac{E_{A1}}{E_A}}\left|1+{\displaystyle \frac{E_f}{E_{A1}}}\left(1{\displaystyle \frac{\stackrel{}{q}\stackrel{}{k}_f}{k_f^2}}\right)\right|=\left|1+{\displaystyle \frac{\omega k_fE_fq\mathrm{cos}\theta _f}{M_Ak_f}}\right|,`$ (3)
with $`\theta _f`$ the angle between $`\stackrel{}{k_f}`$ and $`\stackrel{}{q}`$, and $`\sigma _M`$ the Mott cross section
$`\sigma _M=\left({\displaystyle \frac{\alpha \mathrm{cos}\theta _e/2}{2ϵ\mathrm{sin}^2\theta _e/2}}\right)^2,`$ (4)
with $`\theta _e`$ the angle between the incident and the scattered electron. The electron kinematics is contained in the kinematical factors
$`v_L`$ $`=`$ $`\left({\displaystyle \frac{Q^2}{q^2}}\right)^2,`$ (5)
$`v_T`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{Q^2}{q^2}}\right)+\mathrm{tan}^2{\displaystyle \frac{\theta _e}{2}},`$ (6)
$`v_{TT}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{Q^2}{q^2}}\right),`$ (7)
$`v_{TL}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left({\displaystyle \frac{Q^2}{q^2}}\right)\sqrt{\left({\displaystyle \frac{Q^2}{q^2}}\right)+\mathrm{tan}^2{\displaystyle \frac{\theta _e}{2}}},`$ (8)
$`v_T^{}`$ $`=`$ $`\mathrm{tan}{\displaystyle \frac{\theta _e}{2}}\sqrt{\left({\displaystyle \frac{Q^2}{q^2}}\right)+\mathrm{tan}^2{\displaystyle \frac{\theta _e}{2}}},`$ (9)
$`v_{TL^{}}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\left({\displaystyle \frac{Q^2}{q^2}}\right)\mathrm{tan}{\displaystyle \frac{\theta _e}{2}},`$ (10)
whereas the structure functions are defined in a standard fashion as
$`_L`$ $`=`$ $`|\rho (\stackrel{}{q})_{fi}|^2,`$ (11)
$`_T`$ $`=`$ $`|J(\stackrel{}{q};+1)_{fi}|^2+|J(\stackrel{}{q};1)_{fi}|^2,`$ (12)
$`_{TT}`$ $`=`$ $`2Re\{J^{}(\stackrel{}{q};+1)_{fi}J(\stackrel{}{q};1)_{fi}\},`$ (13)
$`_{TL}`$ $`=`$ $`2Re\{\rho ^{}(\stackrel{}{q})_{fi}(J(\stackrel{}{q};+1)_{fi}J(\stackrel{}{q};1))_{fi}\},`$ (14)
$`_T^{}`$ $`=`$ $`|J(\stackrel{}{q};+1)_{fi}|^2|J(\stackrel{}{q};1)_{fi}|^2,`$ (15)
$`_{TL^{}}`$ $`=`$ $`2Re\{\rho ^{}(\stackrel{}{q})_{fi}(J(\stackrel{}{q};+1)_{fi}+J(\stackrel{}{q};1))_{fi}\},`$ (16)
where $`\rho (\stackrel{}{q})_{fi}`$ is the transition charge density, while $`J(\stackrel{}{q};m=\pm 1)_{fi}`$ is the transition three-current expanded in terms of the standard spherical components.
### B Bound state wave functions.
A relativistic quantum field theory for nucleons ($`\psi `$) interacting with scalar mesons ($`\varphi `$) through a Yukawa coupling $`\overline{\psi }\psi \varphi `$ and with neutral vector mesons ($`V_\mu `$) that couple to the conserved baryon current $`\overline{\psi }\gamma _\mu \psi `$, can be described through a lagrangian density of the type
$`_0`$ $`=`$ $`\overline{\psi }(ı\partial ̸M)\psi +{\displaystyle \frac{1}{2}}(_\mu \varphi ^\mu \varphi m_s^2\varphi ^2){\displaystyle \frac{1}{4}}G_{\mu \nu }G^{\mu \nu }`$ (18)
$`+{\displaystyle \frac{1}{2}}m_v^2V_\mu V^\mu g_v\overline{\psi }\gamma _\mu \psi V^\mu +g_s\overline{\psi }\psi \varphi ,`$
with $`M`$, $`m_s`$ and $`m_v`$ the nucleon, scalar meson and vector meson masses, respectively, and $`G^{\mu \nu }^\mu V^\nu ^\nu V^\mu `$ the vector meson field strength. The scalar and vector fields may be associated with the $`\sigma `$ and $`\omega `$ mesons. The model can be extended to include also $`\pi `$ and $`\rho `$ mesons, as well as the coupling to the photon field. The corresponding lagrangian has the form
$``$ $`=`$ $`_0+{\displaystyle \frac{1}{2}}(_\mu \stackrel{}{\pi }^\mu \stackrel{}{\pi }m_\pi ^2\stackrel{}{\pi }\stackrel{}{\pi })ıg_\pi \overline{\psi }\gamma _5\stackrel{}{\tau }\stackrel{}{\pi }\psi {\displaystyle \frac{1}{4}}\stackrel{}{B}_{\mu \nu }\stackrel{}{B}^{\mu \nu }`$ (21)
$`+{\displaystyle \frac{1}{2}}m_\rho ^2\stackrel{}{b}_\mu \stackrel{}{b}^\mu {\displaystyle \frac{1}{2}}g_\rho \overline{\psi }\gamma _\mu \stackrel{}{\tau }\stackrel{}{b}^\mu \psi {\displaystyle \frac{1}{4}}F_{\mu \nu }F^{\mu \nu }`$
$`eA_\mu [\overline{\psi }\gamma ^\mu {\displaystyle \frac{1}{2}}(1+\tau _3)\psi +(\stackrel{}{b}_\nu \times \stackrel{}{B}^{\nu \mu })_3+(\stackrel{}{\pi }\times (^\mu \stackrel{}{\pi }+g_\rho (\stackrel{}{\pi }\times \stackrel{}{b}^\mu )))_3].`$
Here $`\stackrel{}{\pi }`$, $`\stackrel{}{b}_\mu `$, $`A_\mu `$, $`F_{\mu \nu }`$ are the pion, rho, Maxwell and electromagnetic fields. Further, $`\stackrel{}{B}^{\mu \nu }^\mu \stackrel{}{b}^\nu ^\nu \stackrel{}{b}^\mu g_\rho (\stackrel{}{b}^\mu \times \stackrel{}{b}^\nu )`$ is the $`\rho `$-meson field.
At sufficiently high densities, the meson field operators can be approximated by their expectation values. Within the context of the relativistic Hartree approximation, it can be shown that when starting from the langrangian (21) the following Dirac equation for the baryon field $`\mathrm{\Psi }`$ results :
$`\left[ı\gamma ^\mu _\mu M\mathrm{\Sigma }_H\right]\mathrm{\Psi }=0,`$ (22)
where the self-energy $`\mathrm{\Sigma }_H`$ is defined as
$`\mathrm{\Sigma }_H=g_s\varphi +g_v\gamma _\mu V^\mu +g_\pi \gamma _5\tau _\alpha \pi ^\alpha +{\displaystyle \frac{1}{2}}g_\rho \gamma _\mu \tau _\alpha b^{\mu \alpha }+{\displaystyle \frac{1}{2}}\gamma _\mu (1+\tau _3)A^\mu .`$ (23)
Assuming that the nuclear ground state is spherically symmetric and a parity eigenstate, it can be shown that the pion field does not enter in the Hartree approximation. Furthermore, the meson fields only depend on the radius, and only the time component of the vector fields contribute. The time-independent Dirac equation can then be written as :
$`\widehat{H}\mathrm{\Psi }(\stackrel{}{x})`$ $``$ $`[ı\stackrel{}{\alpha }\stackrel{}{}+g_vV^0(r)+{\displaystyle \frac{1}{2}}g_\rho \tau _\alpha b^{0\alpha }(r)`$ (25)
$`+{\displaystyle \frac{1}{2}}e(1+\tau _3)A^0(r)+\gamma ^0(Mg_s\varphi ^0(r))]=E\mathrm{\Psi }(\stackrel{}{x}).`$
The general solutions to a Dirac equation with spherically symmetric potentials have the form
$$\psi _\alpha (\stackrel{}{x})\psi _{n\kappa mt}(\stackrel{}{x})=\left[\begin{array}{c}ıG_{n\kappa t}(r)/r𝒴_{\kappa m}\eta _t\\ F_{n\kappa t}(r)/r𝒴_{\kappa m}\eta _t\end{array}\right],$$
(26)
where $`n`$ denotes the principal, $`\kappa `$ and $`m`$ the generalized angular momentum and $`t`$ the isospin quantum numbers. The $`𝒴_{\pm \kappa m}`$ are the well-known spin spherical harmonics and determine the angular and spin parts of the wavefunction,
$`𝒴_{\kappa m}={\displaystyle \underset{m_lm_s}{}}<lm_l{\displaystyle \frac{1}{2}}m_s|l{\displaystyle \frac{1}{2}}jm>Y_{l,m_l}\chi _{\frac{1}{2}m_s},`$ (27)
$`j=|\kappa |{\displaystyle \frac{1}{2}},l=\{\begin{array}{cc}\kappa ,\hfill & \kappa >0\hfill \\ (\kappa +1),\hfill & \kappa <0.\hfill \end{array}`$ (30)
The Hartree approximation leads to the following set of coupled equations for the different fields :
$`{\displaystyle \frac{d^2}{dr^2}}\varphi _0(r)+{\displaystyle \frac{2}{r}}{\displaystyle \frac{d}{dr}}\varphi _0(r)m_s^2\varphi _0(r)`$ $`=`$ $`g_s\rho _s(r)`$ (31)
$``$ $`g_s{\displaystyle \underset{\alpha _{occ}}{}}\left({\displaystyle \frac{2ȷ_\alpha +1}{4\pi r^2}}\right)(|G_\alpha (r)|^2|F_\alpha (r)|^2),`$ (32)
$`{\displaystyle \frac{d^2}{dr^2}}V_0(r)+{\displaystyle \frac{2}{r}}{\displaystyle \frac{d}{dr}}V_0(r)m_v^2V_0(r)`$ $`=`$ $`g_v\rho _B(r)`$ (33)
$``$ $`g_v{\displaystyle \underset{\alpha _{occ}}{}}\left({\displaystyle \frac{2ȷ_\alpha +1}{4\pi r^2}}\right)(|G_\alpha (r)|^2+|F_\alpha (r)|^2),`$ (34)
$`{\displaystyle \frac{d^2}{dr^2}}b_0(r)+{\displaystyle \frac{2}{r}}{\displaystyle \frac{d}{dr}}b_0(r)m_\rho ^2\varphi _0(r)`$ $`=`$ $`{\displaystyle \frac{1}{2}}g_\rho \rho _3(r)`$ (35)
$``$ $`{\displaystyle \frac{1}{2}}g_\rho {\displaystyle \underset{\alpha _{occ}}{}}\left({\displaystyle \frac{2ȷ_\alpha +1}{4\pi r^2}}\right)(|G_\alpha (r)|^2+|F_\alpha (r)|^2)(1)^{t_\alpha 1/2},`$ (36)
$`{\displaystyle \frac{d^2}{dr^2}}A_0(r)+{\displaystyle \frac{2}{r}}{\displaystyle \frac{d}{dr}}A_0(r)`$ $`=`$ $`e\rho _P(r)`$ (37)
$``$ $`e{\displaystyle \underset{\alpha _{occ}}{}}\left({\displaystyle \frac{2ȷ_\alpha +1}{4\pi r^2}}\right)(|G_\alpha (r)|^2+|F_\alpha (r)|^2)(t_\alpha +{\displaystyle \frac{1}{2}}),`$ (38)
$`{\displaystyle \frac{d}{dr}}G_\alpha (r)+{\displaystyle \frac{\kappa }{r}}G_\alpha (r)`$ $``$ $`[ϵ_\alpha g_vV_0(r)t_\alpha g_\rho b_0(r)`$ (39)
$``$ $`(t_\alpha +{\displaystyle \frac{1}{2}})eA_0(r)+Mg_s\varphi _0(r)]F_\alpha (r)=0,`$ (40)
$`{\displaystyle \frac{d}{dr}}F_\alpha (r){\displaystyle \frac{\kappa }{r}}F_\alpha (r)`$ $`+`$ $`[ϵ_\alpha g_vV_0(r)t_\alpha g_\rho b_0(r)`$ (41)
$``$ $`(t_\alpha +{\displaystyle \frac{1}{2}})eA_0(r)M+g_s\varphi _0(r)]G_\alpha (r)=0,`$ (42)
$`{\displaystyle _0^{\mathrm{}}}𝑑r(|G_\alpha |^2+|F_\alpha |^2)`$ $`=`$ $`1.`$ (43)
The above equations constitute the basis of the relativistic mean-field approach to the lagrangian of Eq. (21).
A new computer program to solve the above set of coupled non-linear differential equations was developed. Starting from an initial guess of the Woods-Saxon form for the scalar and vector potential, the Dirac equations can be solved iteratively using a shooting point method. Analytic solutions to the equations in the regions of large and small r allow to impose the proper boundary conditions. Once the nucleon wave functions are obtained, the densities and meson fields can be re-evaluated. This procedure is repeated a number of times until convergence for the energy eigenvalues is reached. We adopt the values for the $`\sigma `$, $`\omega `$ and $`\rho `$ masses and coupling constants as they were introduced by Horowitz and Serot .
For the <sup>12</sup>C and <sup>16</sup>O nuclei, the newly developed C-code SOR performed all integrations for a radial extension of the nucleus of 20 fm and a stepsize of 0.01 fm. The coupled Dirac equations were solved for a shooting point lying at 2 fm using a fourth order Runge-Kutta algorithm. As a convergence criterium we imposed a tolerance level as small as 0.001 MeV on all single-particle energy levels. The computed densities for the nuclei <sup>12</sup>C and <sup>16</sup>O, are depicted in Fig. 1. We have verified that these results are comparable to those produced by the TIMORA code , which is widely used to solve the set of Eqs. (31).
### C The eikonal final state.
To construct the scattering states for the ejected nucleons, we consider the hamiltonian (25) that was already used to calculate the bound state wave functions
$`\widehat{H}ı\stackrel{}{\alpha }\stackrel{}{}+\gamma ^0M+\gamma ^0\mathrm{\Sigma }_H(r),`$ (44)
where the self-energy $`\mathrm{\Sigma }_H(r)`$ is given by
$`\mathrm{\Sigma }_H(r)=g_s\varphi _0(r)+g_v\gamma _0V^0(r)+{\displaystyle \frac{1}{2}}g_\rho \gamma _0\tau _\alpha b^{0\alpha }(r)+{\displaystyle \frac{1}{2}}e\gamma _0(1+\tau _3)A^0(r).`$ (45)
With the formal substitutions
$`V_s(r)`$ $``$ $`g_s\varphi _0,`$ (46)
$`V_v(r)`$ $``$ $`g_vV_0(r)+{\displaystyle \frac{1}{2}}g_\rho b_0(r)(1)^{t_\alpha 1/2}+eA_0(r)(t_\alpha +{\displaystyle \frac{1}{2}}),`$ (47)
the time independent Dirac equation for a projectile with relativistic energy $`E=\sqrt{k^2+M^2}`$ and spin state $`s`$, can be cast in the form
$`\widehat{H}\varphi _{\stackrel{}{k},s}^{(+)}=[\stackrel{}{\alpha }\stackrel{}{p}+\beta M+\beta V_s(r)+V_v(r)]\varphi _{\stackrel{}{k},s}^{(+)},`$ (48)
where we have introduced the notation $`\varphi _{\stackrel{}{k},s}^{(+)}`$ for the unbound Dirac states. The computed scalar and vector potentials for the <sup>12</sup>C and <sup>16</sup>O nuclei are displayed in Fig. 2
After some straightforward manipulations, a Schrödinger-like equation for the upper component can be obtained
$`\left[{\displaystyle \frac{p^2}{2M}}+V_c+V_{so}(\stackrel{}{\sigma }\stackrel{}{L}ı\stackrel{}{r}\stackrel{}{p})\right]u_{\stackrel{}{k},s}^{(+)}={\displaystyle \frac{k^2}{2M}}u_{\stackrel{}{k},s}^{(+)},`$ (49)
where the central and spin orbit potentials $`V_c`$ and $`V_{so}`$ are defined as
$`V_c(r)`$ $`=`$ $`V_s(r)+{\displaystyle \frac{E}{M}}V_v(r)+{\displaystyle \frac{V_s(r)^2V_v(r)^2}{2M}},`$ (50)
$`V_{so}(r)`$ $`=`$ $`{\displaystyle \frac{1}{2M[E+M+V_s(r)V_v(r)]}}{\displaystyle \frac{1}{r}}{\displaystyle \frac{d}{dr}}[V_s(r)V_v(r)].`$ (51)
In computing the scattering wave functions, we use the scalar and vector potentials as obtained from the iterative bound state calculations. As a result the initial and final state wave functions are orthogonalized and no spurious contributions can be expected to enter the calculated cross sections.
Since the lower component is related to the upper one through
$`w_{\stackrel{}{k},s}^{(+)}={\displaystyle \frac{1}{E+M+V_sV_v}}\stackrel{}{\sigma }\stackrel{}{p}u_{\stackrel{}{k},s}^{(+)},`$ (52)
the solutions to the equation (49) determine the complete relativistic eigenvalue problem. So far no approximations have been made. Various groups have solved the Dirac equation (49) for the final scattering state using Dirac optical potentials derived from global fits to elastic proton scattering data . Not only are global parametrizations of Dirac optical potentials usually restricted to proton kinetic energies $`T_p1`$ GeV, calculations based on exact solutions of the Dirac equation frequently become impractical at higher energies. This is particularly the case for approaches that rely on partial-wave expansions in determining the transition matrix elements. To overcome these complications, we solve the Dirac equation (49) in the eikonal limit . In intermediate-energy proton scattering ($`T_p`$ 500 MeV) the eikonal approximation was shown to reproduce fairly well the exact Dirac partial wave results . Following the discussion of Ref. , we define the average momentum $`\stackrel{}{K}`$ and the momentum transfer $`\stackrel{}{q}`$ in terms of the projected initial ($`\stackrel{}{k}_i`$) and final momentum ($`\stackrel{}{k}_f`$) of the ejectile
$`\stackrel{}{q}`$ $`=`$ $`\stackrel{}{k}_f\stackrel{}{k}_i,`$ (53)
$`\stackrel{}{K}`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\stackrel{}{k}_f+\stackrel{}{q}).`$ (54)
In the eikonal, or, equivalently, the small-angle approximation ($`qk_i`$) the following operatorial substitution is made in computing the scattering wave function
$`p^2=[(\stackrel{}{p}\stackrel{}{K})+\stackrel{}{K}]^22\stackrel{}{K}\stackrel{}{p}K^2.`$ (55)
After introducing this approximate relation, the Dirac equation for the upper component (49) becomes
$`[ı\stackrel{}{K}\stackrel{}{}K^2+M(V_c+V_{so}[\stackrel{}{\sigma }(\stackrel{}{r}\times \stackrel{}{K})ı\stackrel{}{r}\stackrel{}{K}])]u_{\stackrel{}{k},s}^{(+)}=0,`$ (56)
where the momentum operators in the spin orbit and Darwin terms are substituted by $`\stackrel{}{K}`$. Remark that the above equation is now linear in the momentum operator. In the eikonal limit, the scattering wave functions take on the form
$`u_{\stackrel{}{k},s}^{(+)}=e^{ı\stackrel{}{k}\stackrel{}{r}}e^{ıS(\stackrel{}{r})}\chi _{\frac{1}{2}m_s}.`$ (57)
Inserting this into Eq. (56), yields an expression for the eikonal phase . Defining the $`z`$-axis along the direction of the average momentum $`\stackrel{}{K}`$, this phase can be written in an integral form as :
$`ıS(\stackrel{}{b},z)=ı{\displaystyle \frac{M}{K}}{\displaystyle _{\mathrm{}}^z}𝑑z^{}[V_c(\stackrel{}{b},z^{})+V_{so}(\stackrel{}{b},z^{})[\stackrel{}{\sigma }(\stackrel{}{b}\times \stackrel{}{K})ıKz^{}]],`$ (58)
where we have introduced the notation $`\stackrel{}{r}(\stackrel{}{b},z)`$. The scattering wave function, which is proportional to
$`\varphi _{\stackrel{}{k},s}^{(+)}\left[\begin{array}{c}1\\ \frac{1}{E+M+V_sV_v}\stackrel{}{\sigma }\stackrel{}{p}\end{array}\right]e^{ı\stackrel{}{k}\stackrel{}{r}}e^{ıS(\stackrel{}{r})}\chi _{\frac{1}{2}m_s},`$ (61)
is normalized such that
$`\overline{\varphi _{\stackrel{}{k},s}^{(+)}}\varphi _{\stackrel{}{k},s}^{(+)}=1.`$ (62)
This wave function differs from the plane-wave solution in two respects. First, the lower component exhibits the dynamical enhancement due to the combination of the scalar and vector potentials. Second, the eikonal phase $`e^{ıS(\stackrel{}{r})}`$ accounts for the interactions that the struck nucleon undergoes in its way out of the nucleus. The calculation of the eikonal phase (58) involves a transformation to a reference frame other than the usual laboratory or center-of-mass frame, namely the frame where the average momentum is pointing along the $`z`$-axis. As the eikonal phase has to be re-evaluated for every $`(\stackrel{}{b},z)`$ point in space, the Dirac eikonal $`(e,e^{}p)`$ calculations are very demanding as far as computing power is concerned. In evaluating the matrix elements, the radial integrations were performed on a 0.1 fm mesh. It is worth remarking that the standard Glauber approach followed in many studies involves an extra approximation apart from the ones discussed above. Indeed, in evaluating the eikonal phase from Eq. (58) one frequently approximates the z-dependence of the potentials by a delta function.
### D Off-shell electron-proton coupling
We express the matrix elements of the nucleon current in the usual form
$`<P_fS_f|J^\mu |P_iS_i>=\overline{u}_f\mathrm{\Gamma }^\mu (P_f,P_i)u_i,`$ (63)
where $`\mathrm{\Gamma }^\mu `$ is the electromagnetic vertex function for the nucleon and $`u_i`$ ($`u_f`$) the nucleon spinors. As discussed in many works , some arbitrariness, often referred to as the “off-shell ambiguity”, surrounds the choice for the functional form of the vertex function $`\mathrm{\Gamma }^\mu `$. For a free nucleon, $`\mathrm{\Gamma }^\mu `$ can be expressed in several fully equivalent forms
$`\mathrm{\Gamma }_{cc1}^\mu `$ $`=`$ $`G_M(Q^2)\gamma ^\mu {\displaystyle \frac{\kappa }{2M}}F_2(Q^2)(P_i^\mu +P_f^\mu ),`$ (64)
$`\mathrm{\Gamma }_{cc2}^\mu `$ $`=`$ $`F_1(Q^2)\gamma ^\mu +ı{\displaystyle \frac{\kappa }{2M}}F_2(Q^2)\sigma ^{\mu \nu }q_\nu ,`$ (65)
$`\mathrm{\Gamma }_{cc3}^\mu `$ $`=`$ $`{\displaystyle \frac{1}{2M}}F_1(Q^2)(P_i^\mu +P_f^\mu )+ı{\displaystyle \frac{1}{2M}}G_M(Q^2)\sigma ^{\mu \nu }q_\nu ,`$ (66)
where $`F_1`$ is the Dirac, $`F_2`$ the Pauli form factor and $`\kappa `$ is the anomalous magnetic moment. The relation with the Sachs electric and magnetic form factors is established through $`G_E=F_1\tau \kappa F_2`$ and $`G_M=F_1+\kappa F_2`$, with $`\tau Q^2/4m^2`$.
When considering bound (or, “off-shell”) nucleons, however, the above vertex functions can no longer be guaranteed to produce the same results. As a matter of fact, explicit current conservation is rather an exception than a rule in most calculations that deal with $`(e,e^{}p)`$ reactions from finite nuclei. In nuclear physics, the most widely used procedure to “effectively” restore current conservation is based on modifying the longitudinal component of the nuclear vector current using the substitution
$`J_z{\displaystyle \frac{\omega }{q}}J_0.`$ (67)
This procedure is partly inspired on the observation that meson-exchange and isobar terms enter the charge current operator in a higher relativistic order than they used to do for the vector current. There exist several other prescriptions which are meant to restore current conservation. Along similar lines, the charge operator can be replaced by
$`J_0{\displaystyle \frac{q}{\omega }}J_z.`$ (68)
One can also construct a vertex function that garantuees current conservation for any initial and final nucleon state. This can be achieved for example by adding an extra term to the vertex
$`\mathrm{\Gamma }_{DON}^\mu `$ $`=`$ $`F_1(Q^2)\gamma ^\mu +ı{\displaystyle \frac{\kappa }{2M}}F_2(Q^2)\sigma ^{\mu \nu }q_\nu +F_1(Q^2){\displaystyle \frac{\mathit{}q^\mu }{Q^2}},`$ (69)
which is also equivalent to the Eqs. (64-66) in the free nucleon case. An operator derived from the generalized Ward-Takahashi identity reads
$`\mathrm{\Gamma }_{WT}^\mu `$ $`=`$ $`\gamma ^\mu ı{\displaystyle \frac{\kappa }{2M}}F_2(Q^2)\sigma ^{\mu \nu }q_\nu +[F_1(Q^2)1]{\displaystyle \frac{\mathit{}q^\mu +Q^2\gamma ^\mu }{Q^2}}.`$ (70)
## III Results
### A Final state interactions and the eikonal approximation
We start our $`(e,e^{}p)`$ investigations within the relativistic eikonal approximation for the kinematics of an <sup>16</sup>O$`(e,e^{}p)`$ experiment that was recently performed at Jefferson Lab . In this experiment, the separated <sup>16</sup>O$`(e,e^{}p)`$ structure functions are measured at $`Q^2`$ = 0.8 (GeV/c)<sup>2</sup> and $`\omega `$ = 0.439 GeV for missing (or, initial) proton momenta $`p_m=\stackrel{}{k_f}\stackrel{}{q}`$ below 355 MeV/c. The variation in missing momentum was achieved by varying the detection angle of the ejected proton with respect to the direction of the momentum transfer (“quasi-perpendicular kinematics”). The measured cross sections for knockout from the $`1p_{1/2}`$ and $`1p_{3/2}`$ levels are depicted in Fig. 3 along with the predictions of our relativistic eikonal calculations. A spectroscopic factor of 0.6 was adopted for all bound levels, and the standard dipole form was used for the electromagnetic form factors. At low missing momenta, the eikonal results shown in Fig. 3 produce a fair description of the data. As a comparison, the results of a relativistic plane wave calculation in the impulse approximation (RPWIA) are also displayed. Through comparing the plane-wave and the eikonal calculations, thereby keeping all other ingredients of the calculations identical, one can evaluate how the eikonal method deals with final state interactions (FSI). In the eikonal calculations, the dips of the RPWIA calculations are filled in, and, at low missing momenta the RPWIA cross sections are reduced. These two features reflect nothing but the usual impact of the final-state interactions on the A$`(e,e^{}p)`$ angular cross sections. The limitations of the eikonal approximation ($`qk_i`$) are immediately visible at higher missing momenta ($`p_m`$ 250 MeV/c). Here, the eikonal cross sections largely overshoot both the RPWIA results and the data and should by no means be considered as realistic. It is worth remarking that the data closely follow the trend set by the RPWIA curves. As a matter of fact, whereas the eikonal calculations predict huge effects from final-state interactions at large transverse missing momenta, the data seem to suggest rather the opposite effect. We consider this observation as one of the major findings of this work.
One may wonder whether the observed behaviour of the eikonal results at higher missing momenta in Fig. 3 is a mere consequence of the small-angle approximation contained in Eq. (55), or whether the adopted model assumptions for computing the scattering states is also (partly) at the origin of this pathological behavior. To address this question, we have performed calculations for various fixed recoil angles $`\theta `$ defined as
$$\mathrm{cos}\theta =\frac{\stackrel{}{p}_m\stackrel{}{q}}{\left|\stackrel{}{p}_m\right|\left|\stackrel{}{q}\right|}.$$
(71)
The results are displayed in terms of the reduced cross section $`\rho `$ which is defined in the standard fashion as the differential cross section, divided by a kinematical factor times the $`\mathrm{`}\mathrm{`}CC1^{\prime \prime }`$ off-shell electron-nucleon cross section of Ref. . For the results of Figure 4 we considered in-plane kinematics at a fixed value of the outgoing proton momentum ($`k_f`$=1 (GeV)/c) and an initial electron energy of 2.4 GeV. The variation in missing momentum is achieved by changing the $`q`$. For recoil angles $`\theta `$ = 0<sup>o</sup> (“parallel kinematics”) the eikonal calculations do not exhibit an unrealistic behavior up to $`p_m`$=0.5 GeV/c, which is the highest missing momentum considered here. With increasing recoil angles, and consequently, growing “transverse” components in the missing momenta the “unrealistic” behaviour of the eikonal results becomes manifest. Accordingly, the accuracy of the eikonal method based on the small-angle approximation of Eq. (55) can only be guaranteed for proton knockout in a small cone about the momentum transfer. A similar quantitative behaviour as a function of the recoil angle to what is observed in Fig. 4 was reported in Ref. for d$`(e,e^{}p)`$n cross sections determined in a Glauber framework. We conclude this section with remarking that the eikonal method does not exclude situations with high initial (or, missing) momenta, it only requires that the perpendicular component of ejectiles’s momentum $`\stackrel{}{k}_f`$ is sufficiently small. It speaks for itself that such conditions are best fulfilled as one approaches parallel kinematics. This observation puts serious constraints on the applicability of the Glauber method, that is based on the eikonal approximation, for modelling the final-state interactions in high-energy $`(e,e^{}p)`$ reactions from nuclei. However, it should be noted that our framework does use purely real scalar and vector potentials. More realistic scattering potentials demand an imaginary part that accounts for the inelastic channels that are open during the reaction process. The Glauber approach effectively includes these inelastic channels and on these grounds one may expect that its range of applicability is somewhat wider than what is observed here. With the eye on defining the region of validity for the eikonal approximation more clearly, we have studied differential cross sections for various $`Q^2`$. In Fig. 5, we display the computed differential cross sections for the <sup>12</sup>C$`(e,e^{}p)^{11}`$B$`(1p_{3/2}^1)`$ process against the missing momentum for $`Q^2`$ varying between 1 and 20 (GeV/c)<sup>2</sup>. Hereby, quasi-elastic conditions were imposed. The arrow indicates the missing momentum where the slope of the eikonal differential cross section starts deviating from the trend set by the RPWIA cross section. In the light of the conclusions drawn from the comparison between data and the eikonal curves in Fig. 3, the eikonal results should be regarded with care beyond this missing momentum. Furthermore, it is clear that the change in the slope of the angular cross section becomes more and more pronounced as $`Q^2`$ increases. It is apparent from Fig. 5 that the eikonal differential cross section changes slope at about $`p_m`$=250 MeV/c for all values of $`Q^2`$ considered. We remark that we imposed quasi-elastic conditions for all cases contained in Fig. 5. As a consequence, the momentum of the ejected nucleon varies quite dramatically as one moves up in $`Q^2`$. The uniform behaviour of all curves contained in Fig. 5 allows one to write down a relation between the transferred momentum $`\stackrel{}{q}`$ and the polar scattering angle $`\theta `$ : $`|\stackrel{}{q}|\theta 250`$ MeV rad. This simple relation could serve as a conservative guideline to determine the opening angle of the cone in which the outgoing proton momentum has to reside to ascertain that the eikonal approximation produces “realistic” results. This limitation of the eikonal method can also be inferred from the results contained in Refs. . Indeed, in Figs. 3 and 4 of Ref. one can confirm that the above relation between $`|\stackrel{}{q}|`$ and $`\theta `$ defines the missing momentum at which a sudden change in the $`p_m`$ dependence of the calculated cross sections is observed. The above relation can be understood as follows. In quasi-perpendicular and quasi-elastic kinematics, the missing momentum roughly equals the transverse momentum of the ejected nucleon. With increasing momentum transfer, the longitudinal momentum of the escaping nucleon increases correspondingly while its transverse momentum has to stay smaller than the suggested value of 250 MeV/c. Hence, the sine of the angle between the transferred momentum and the ejectile’s momentum has to decrease. Since we are dealing with small angles, sin($`\theta `$) can be approximated by $`\theta `$. The opening angle of the cone in which the eikonal approximation is valid, can be inferred to be independent of $`Q^2`$ in the Lorentz frame where the ejected nucleon is at rest. When transforming back to the labframe, lateral dimensions become dilated, and, thus, angles contracted.
### B The $`Q^2`$ evolution of off-shell effects
A major point of concern in any A$`(e,e^{}p)`$B calculation are the ambiguities regarding the off-shell electron-proton coupling. Most calculations do not obey current conservation and a variety of prescriptions have been proposed to partially cure this deficiency. Here we adopt a heuristic view and estimate the sensitivity of the calculated observables by comparing the results obtained with different viable prescriptions for the electron-proton coupling. Amongst the infinite number of possible prescriptions for the off-shell electron-proton coupling we have selected four that are frequently used in literature. Figure 6 shows the separated structure functions for $`1p_{1/2}`$ knockout in the kinematics of Fig. 3. Current conservation was imposed by either modifying the longitudinal component of the vector current operator (hereafter denoted as the “J0 method”), or by modifying the charge operator (hereafter denoted as the “J3 method”), along the lines of Eqs. (67) and (68). Note that for the operator of Eq. (69), both methods yield the same results, since, by construction, this operator is current conserving, regardless of the method adopted to compute the wave function for the initial and final state.
Turning to the results shown in Fig. 6, the predicted strengths in the longitudinal structure functions $`R_L`$ and $`R_{TL}`$ depend heavily on the choice made for the electron-proton coupling. For the $`CC1`$ prescription, for example, the values obtained with the J3 method are several times bigger than those obtained within the J0 method. The predicted differences among the various current operators within one scheme (“J0” or “J3”) are also sizeable. The ambiguities are, however, much smaller for the calculations performed with the $`J_z\frac{\omega }{q}J_0`$ substitution. This clearly speaks in favor of this recipe which is mostly used in A$`(e,e^{}p)`$ calculations. The $`R_{TT}`$ and $`R_T`$ structure functions are, obviously, insensitive to whether the “J0” or “J3” method is adopted. All adopted electron-proton couplings but the $`CC1`$ one produce the same results in the $`R_T`$ and $`R_{TT}`$ responses.
With increasing $`Q^2`$ and the corresponding decreasing distance scale, the role of off-shell ambiguities in the photon-nucleus coupling is expected to decline and the impulse approximation is believed to become increasingly accurate. In order to investigate the degree and rate to which this virtue may be realized, we have performed calculations for kinematics in the range of $`0.8Q^220`$(GeV/c)<sup>2</sup>. We use two techniques to estimate the relative importance of the off-shell effects as a function of $`Q^2`$. First, results computed with the “J0” and “J3” method can be compared. Second, predictions with various choices for the electron-proton coupling are confronted with one another. The validity of the IA is then established whenever the final result happens to become independent of the adopted choice. In order to assess the degree to which this independence is realized, we have considered ratios of structure functions for some fixed kinematics but calculated with different choices for the electron-proton coupling. As a benchmark calculation, we have computed <sup>12</sup>C$`(e,e^{}p)^{11}`$B$`(1p_{3/2}^1)`$ observables in quasi-elastic kinematics for several values of the four-momentum transfer. The results are shown in Figs. 7 and 8. Fig. 7 shows for several observables the ratio of the values obtained with the “J3” scheme to the corresponding prediction using the “J0” scheme. Fig. 8 shows the ratio of the strengths obtained with the $`CC1`$ vertex function compared to the corresponding predictions with the $`CC2`$ form. Remark that in the limit of vanishing off-shell effects, these ratios should equal one. It is indeed found that the calculations that are based on the substitution $`J_z\frac{\omega }{q}J_0`$, tend to converge to those based on $`J_0\frac{q}{\omega }J_z`$ with increasing energy transfer. This is particularly the case at low missing momenta, where the decrease in the longitudinal response is almost exponential. The overall behaviour is identical for the higher missing momentum case ($`p_m`$ = 150 MeV/c), but the rate of decrease is somewhat slower. This can be attributed to the fact that at higher momenta, hence, greater angles, the transverse components of the vertex functions play a more important role. Looking at Fig. 8 one can essentially draw the same conclusions. The predictions with the different prescriptions also converge to each other as the energy is increased. Again this convergence is more pronounced for the low missing momentum case. This feature is most apparent in the purely transverse channel, which dominates the cross section at sufficiently high energies. It appears thus as if off-shell ambiguities, speaking in terms of strenghts and absolute cross sections, are of far less concern at higher $`Q^2`$ than they used to be in the $`Q^2`$ 1 (GeV/c)<sup>2</sup> region, where most of the data have been accumulated up to now. The interference structure functions $`R_{TT}`$ and $`R_{TL}`$ are subject to off-shell ambiguities that are apparently extending to the highest four-momentum transfers considered here. This feature was already established in Ref. and explained by referring to the large weight of the negative energy solutions in the interference structure functions $`R_{TL}`$ and $`R_{TT}`$.
### C Relativistic Effects
Recently, there have been several claims for strong indications for genuine (or, “dynamic”) relativistic effects in A$`(\stackrel{}{e},e^{}\stackrel{}{p})`$ observables . In an attempt to implement some of these effects in calculations based on a Schrödinger picture, several techniques to obtain a “relativized version” of the electron-nucleus vertex have been developed. In leading order in a $`p/M`$ expansion these “relativized” electron-nucleus vertices typically miss the coupling between the lower components in the bound and scattering states. For that reason, we interpret the effect of the coupling between the lower components in the bound and scattering states as a measure for the importance of relativistic effects. In Fig. 9 we display results of fully relativistic <sup>12</sup>C$`(e,e^{}p)^{11}`$B$`(1p_{3/2}^1)`$ calculations and calculations in which the specific coupling between the lower components in the bound and scattering states have been left out. We consider quasi-elastic conditions and study the $`Q^2`$ evolution of the structure functions for two values of the missing momentum ($`p_m`$ = 0 and 150 MeV/c) both corresponding with small recoil angles. Hence, the results of Fig. 9 refer to kinematic conditions for which the eikonal approximation is justified. A rather complex and oscillatory $`Q^2`$ dependence of the relativistic effects emerges from our numerical calculations. Looking first at the $`p_m`$ 0 MeV/c case, which nearly corresponds with parallel kinematics, we observe that for both the longitudinal and transverse structure functions, the impact of the coupling amongst the lower components first increases, and then tends to become fairly constant for higher values of $`\omega `$. The genuine relativistic effect stemming from the coupling between the lower components in the initial and final states is larger in the longitudinal than in the transverse channel. It is noteworthy that in the cross section the impact of the “relativistic dynamical effects” never exceeds the 10 % level. If we turn our attention to the interference structure functions $`R_{TL}`$ and $`R_{TT}`$, the relativistic effects grow in importance. Especially for the $`R_{TL}`$ structure function the effects are large and extend to the smallest values of $`Q^2`$ considered here. This enhanced sensitivity of the $`R_{TL}`$ response to relativistic effects, even when relatively low values of $`Q^2`$ are probed, complies with the conclusions drawn in other studies . Also the tendency of the relativistic effects to increase the cross section when higher values of $`p_m`$ are probed complies with the findings of earlier studies . A quantity that is relatively easy to access experimentally and depends heavily upon the $`R_{TL}`$ term, is the so-called left-right asymmetry $`A_{LT}`$
$`A_{LT}={\displaystyle \frac{\sigma (\varphi =0^{})\sigma (\varphi =180^{})}{\sigma (\varphi =0^{})+\sigma (\varphi =180^{})}}={\displaystyle \frac{v_{TL}R_{TL}}{v_LR_L+v_TR_T+v_{TT}R_{TT}}}.`$ (72)
In Fig. 10 we have plotted the left-right asymmetry for both $`1p_{1/2}`$ and $`1p_{3/2}`$ knockout from <sup>16</sup>O in the kinematics of Fig. 3. It is indeed verified that the asymmetry is very sensitive to relativistic effects. As has been reported, relativistic effects enhance the asymmetry further, and this enhancement is more pronounced for the $`1p_{1/2}`$ knockout reaction. The role played by the lower components in this dynamical enhancement of the left-right asymmetry can be further clearified by looking at the results of Fig. 11. In this figure, we plot the left-right asymmetry for $`1p_{3/2}`$ knockout from <sup>12</sup>C, for different $`Q^2`$ and quasi-elastic conditions. Looking at the fully relativistic curves, we observe a gradual decrease of the asymmetry with increasing $`Q^2`$. At the same time, the relative contribution of the “non-relativistic” contribution to $`A_{LT}`$ diminishes. This feauture indicates that the asymmetry $`A_{LT}`$ is nearly exclusively generated by the coupling between the lower components as $`Q^2`$ increases.
## IV Summary
We have outlined a fully relativistic eikonal framework for calculating cross sections for $`(e,e^{}p)`$ reactions from spherical nuclei at intermediate and high four-momentum transfers and carried out <sup>12</sup>C$`(e,e^{}p)`$ and <sup>16</sup>O$`(e,e^{}p)`$ calculations for a variety of kinematical conditions, thereby covering four-momentum transfers in the range $`0.8Q^220`$ (GeV/c)<sup>2</sup>. Our results illustrate that the validity of the eikonal method is confined to proton emission in a cone with a relatively small opening angle about the direction of the virtual photon’s momentum. This observation puts serious constraints on the applicability of the Glauber method, that is based on the eikonal approximation, for modelling the final-state interactions in high-energy $`(e,e^{}p)`$ reactions from nuclei. Incorporation of the inelastic channels in the eikonal method is however needed to fully appreciate the limits of the Glauber model, and work along these lines is in progress. In line with the expectations, our investigations illustrate that the uncertainties induced by off-shell ambiguities on the calculated observables diminish as $`Q^2`$ increases. Nevertheless, in the relativistic eikonal framework four-momentum transfers of the order 5 (GeV/c)<sup>2</sup> appear necessary to assure that the effect of the off-shell ambiguities can be brought down to the percent level. Our theoretical framework permits to assess the impact of the relativistic effects over a wide energy range. The impact of the lower components on the $`(e,e^{}p)`$ observables is observed to be significant over the whole $`Q^2`$ range studied. Especially the left-right asymmetry lends itself very well to study these effects of genuine relativistic origin.
Acknowledgement
This work was supported by the Fund for Scientific Research of Flanders under contract No. 4.0061.99 and the University Research Council.
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# 1 Introduction
## 1 Introduction
Consider a suitable observable in a fully developed hydrodynamic turbulent flow. For example, this may be a longitudinal or transverse velocity difference, a temperature or a pressure difference. The dynamics is effectively described by some highly nonlinear set of model equations—for example the Navier-Stokes equation. Since nobody is able to solve this equation exactly, it is desirable to find some effective statistical description using methods from generalized statistical mechanics.
The underlying idea is quite similar to what was going on more than a hundred years ago, proceeding from classical mechanics to ordinary thermodynamics. Although nobody was able to ‘solve’ the classical $`N`$-body problem with $`N3`$ exactly, one still was able to develop quite a successful effective theory of a gas of $`10^{23}`$ particles by extremizing the Boltzmann-Gibbs entropy. For fully developed turbulent spatio-temporal chaotic systems, the ordinary Boltzmann-Gibbs statistics is not sufficient to describe the (non-Gaussian) stationary state. But still we can try to develop an effective probabilistic theory using more general information measures.
The idea to start from an extremization principle for turbulent flows is actually not new. For example, Cocke has presented some work on turbulence where he extremizes the Fisher information. Castaing et al. also start from another extremum principle to derive probability densities in fully developed turbulence. Here we will work within a new approach to turbulence based on extremizing the Tsallis entropies
$$S_q=\frac{1}{q1}\left(1\underset{i}{}p_i^q\right).$$
(1)
The $`p_i`$ are the probabilities of the various microstates of the physical system, and $`q`$ is the non-extensitivity parameter. The ordinary Boltzmann-Gibbs entropy is obtained in the limit $`q1`$.
Generally, the Tsallis entropies are known to have several nice properties. They are positive, concave, take on their extremum for the uniform distribution and preserve the Legendre transform structure of thermodynamics. On the other hand, they are non-extensive (non-additive for independent subsystems).
Extremizing $`S_q`$ under suitable norm and energy constraints, one arrives at a generalized version of the canonical distribution given by
$$p_i=\frac{1}{Z_q}(1+(q1)\beta ϵ_i)^{\frac{1}{1q}},$$
(2)
where
$$Z_q=\underset{i}{}(1(1q)\beta ϵ_i)^{\frac{1}{1q}}$$
(3)
is the partition function, $`\beta =1/(kT)`$ is a suitable inverse temperature variable, and the $`ϵ_i`$ are the energies of the microstates $`i`$. Ordinary thermodynamics is recovered for $`q1`$. One can also work with the escort distributions , defined by $`P_i=p_i^q/p_i^q`$. If $`\beta `$ is allowed to depend on $`q`$, the escort distribution is of the same form as eq. (2), with a new $`q^{}`$ defined by $`q/(q1)=:1/(q^{}1)`$.
All we have to decide now is what we should take for the effective energy levels $`ϵ_i`$ in the turbulence application. This depends on the problem considered. For example, turbulence in different dimensions ought to yield different effective energy levels. Moreover, different observables will also lead to different effective energies. In a turbulent 3-dimensional flow for example, temperature differences should be described by different effective energies than velocity differences. This is clear from the fact that the experimentally observed stationary probability distributions (slightly) differ.
At this stage one has to turn to some sort of model. In the following we will consider a simple local model for longitudinal velocity differences that seems to reproduce the statistics of true turbulence experiments quite well .
## 2 Perturbative approach to chaotically driven systems
The model is based on a generalization of the Langevin equation to deterministic chaotic driving forces . We denote the local longitudinal velocity difference of two points in the liquid separated by a distance $`r`$ by $`u`$. Clearly $`u`$ relaxes with a certain damping constant $`\gamma `$ and at the same time is driven by deterministic chaotic force differences $`F_{chaot}(t)`$ in the liquid, which are very complicated. Hence a very simple local model is
$$\dot{u}=\gamma u+F_{chaot}(t).$$
(4)
The force $`F_{chaot}(t)`$ is not Gaussian white noise but a complicated deterministic chaotic forcing. It changes on a typical time scale $`\tau `$, which is smaller than the relaxation time $`\gamma ^1`$, so $`\gamma \tau `$ is a small parameter. We may effectively discretize in time and consider the rescaled kick force
$$F_{chaot}(t)=(\gamma \tau )^{1/2}\underset{n=0}{\overset{\mathrm{}}{}}x_n\delta (tn\tau ),$$
(5)
where the $`x_n`$ are the iterates of an appropriate stroboscopic chaotic map $`T`$. Integrating eq. (4) one obtains
$`x_{n+1}`$ $`=`$ $`T(x_n)`$
$`u_{n+1}`$ $`=`$ $`\lambda u_n+\sqrt{\gamma \tau }x_{n+1},`$ (6)
where $`u_n:=u(n\tau +0)`$ and $`\lambda :=e^{\gamma \tau }`$. For $`\gamma \tau 0`$, $`t=n\tau `$ finite, and so-called $`\phi `$-mixing deterministic maps $`T`$ it has been shown that the $`u`$-dynamics converges to the Ornstein-Uhlenbeck process, regarding the initial values $`x_0`$ as random variables. Hence the invariant density of $`u`$ becomes Gaussian in this limit. For finite $`\gamma \tau `$, on the other hand, the invariant density is non-Gaussian. It can have fractal and singular properties if $`\gamma \tau `$ is large . But for small $`\gamma \tau `$ it approaches the Gaussian distribution provided $`T`$ is $`\phi `$-mixing.
The route to the Gaussian limit behaviour has been investigated in detail in for many different chaotic maps $`T`$. Within a well defined universality class (defined by square root scaling of the first order correction) it was found that for small enough $`\gamma \tau `$ the invariant density of $`u`$ is always given by
$$p(u)=\frac{1}{\sqrt{2\pi }}e^{\frac{1}{2}u^2+c\sqrt{\gamma \tau }(u\frac{1}{3}u^3)}+O(\gamma \tau ),$$
(7)
provided the variable $`u`$ is rescaled such that the variance of the distribution is 1. $`c`$ is a non-universal constant. This means, not only the Gaussian limit distribution is universal (i.e. independent of details of $`T`$), but also the way the Gaussian is approached if the time scale ratio $`\gamma \tau `$ goes to 0. Much more details on this can be found in .
It is now reasonable to assume that also the local chaotic forces acting on longitudinal velocity differences in a turbulent flow lie in this universality class. We can then use eq. (7) to construct effective energy levels $`ϵ_i`$ for the non-extensive theory.
## 3 Constructing effective energy levels
Eq. (7) corresponds to a Boltzmann factor
$$p_i=\frac{1}{Z}e^{\beta ϵ_i}$$
(8)
with $`\beta =1/(kT)=1`$, $`Z=\sqrt{2\pi }`$, and energy $`ϵ_i`$ formally given by
$$ϵ_i=\frac{1}{2}u^2c\sqrt{\gamma \tau }(u\frac{1}{3}u^3)+O(\gamma \tau ).$$
(9)
The main contribution is the kinetic energy $`\frac{1}{2}u^2`$, but in addition there is also a small asymmetric term with a universal $`u`$-dependence.
The complicated hydrodynamic interactions and the cascade of energy dissipating from larger to smaller levels is now expected to be effectively described by a non-extensive theory with the above energy levels (see for a related cascade model). We obtain from eq. (2) and (9) the formula
$$p(u)=\frac{1}{Z_q}\left(1+\beta (q1)\left(\frac{1}{2}u^2c\sqrt{\gamma \tau }(u\frac{1}{3}u^3)+O(\gamma \tau )\right)\right)^{\frac{1}{q1}}.$$
(10)
This equation is in very good agreement with experimentally measured probability densities. For detailed comparisons with various turbulence experiments, see . The parameter $`\beta `$ is determined by the condition that the distribution should have variance 1. For $`\gamma \tau =0`$ this is achieved for $`\beta =2/(53q)`$.
Note that we have identified a small parameter $`\gamma \tau `$ in our approach. It is the ratio of two time scales—that of the local forcing and that of the relaxation to the stationary state. One may conjecture that it is related to the inverse Reynolds number $`R_\lambda ^1`$ . The turbulent statistics is determined by a kind of effective non-extensive field theory with the formal coupling constant $`\sqrt{\gamma \tau }`$. A perturbative approach is possible since $`\sqrt{\gamma \tau }`$ is small. In fact, for the dynamics (6) one can work with analogues of Feynman graphs related to higher-order correlations of the chaotic dynamics . Eq. (10) is just obtained by first-order perturbation theory—the complete theory is the infinite-order theory taking into account all orders of $`\sqrt{\gamma \tau }`$. But first we have to understand the ‘free’ turbulent field theory obtained for $`\gamma \tau =0`$. Here almost everything can be calculated analytically.
## 4 The ‘free’ turbulent field theory
If $`\gamma \tau =0`$ the moments can be easily evaluated. In we obtained
$$|u|^m=_{\mathrm{}}^{\mathrm{}}p(u)|u|^m𝑑u=\left(\frac{2k}{\beta }\right)^{\frac{m}{2}}\frac{B(\frac{m+1}{2},k\frac{m+1}{2})}{B(\frac{1}{2},k\frac{1}{2})}$$
(11)
with $`k`$ defined as $`k:=1/(q1)`$ ($`k`$ needs not to be integer). The beta function is defined as
$$B(x,y)=\frac{\mathrm{\Gamma }(x)\mathrm{\Gamma }(y)}{\mathrm{\Gamma }(x+y)}.$$
(12)
We now show that this formula for the moments can be significantly simplified. First, we obtain from the definition of the beta function
$$|u|^m=\left(\frac{2k}{\beta }\right)^{\frac{m}{2}}\frac{\mathrm{\Gamma }\left(\frac{1}{2}+\frac{m}{2}\right)\mathrm{\Gamma }\left(k\frac{1}{2}\frac{m}{2}\right)}{\mathrm{\Gamma }\left(\frac{1}{2}\right)\mathrm{\Gamma }\left(k\frac{1}{2}\right)}$$
(13)
Generally, one has
$$\mathrm{\Gamma }(x+n)=\mathrm{\Gamma }(x)\underset{j=0}{\overset{n1}{}}(x+j)$$
(14)
for natural numbers $`n`$. Suppose $`\frac{m}{2}=:n𝐍`$, then it follows from eq. (13) and (14)
$$u^m=\left(\frac{2k}{\beta }\right)^{\frac{m}{2}}\underset{j=0}{\overset{\frac{m}{2}1}{}}\frac{2j+1}{2j+2km1}$$
(15)
In particular, we obtain for the first few moments
$`u^2`$ $`=`$ $`\left({\displaystyle \frac{2k}{\beta }}\right){\displaystyle \frac{1}{2k3}}`$ (16)
$`u^4`$ $`=`$ $`\left({\displaystyle \frac{2k}{\beta }}\right)^2{\displaystyle \frac{1}{2k5}}{\displaystyle \frac{3}{2k3}}`$ (17)
$`u^6`$ $`=`$ $`\left({\displaystyle \frac{2k}{\beta }}\right)^3{\displaystyle \frac{1}{2k7}}{\displaystyle \frac{3}{2k5}}{\displaystyle \frac{5}{2k3}}`$ (18)
This can be used to evaluate the complete set of hyperflatness factors $`F_m`$ defined as
$$F_m=\frac{u^{2m}}{u^2^m}$$
(20)
We obtain
$`F_1`$ $`=`$ $`1`$ (21)
$`F_2`$ $`=`$ $`3{\displaystyle \frac{2k3}{2k5}}`$ (22)
$`F_3`$ $`=`$ $`35{\displaystyle \frac{(2k3)^2}{(2k7)(2k5)}}`$ (23)
$`F_4`$ $`=`$ $`357{\displaystyle \frac{(2k3)^3}{(2k9)(2k7)(2k5)}}`$ (24)
and generally
$$F_m=(2m1)!!\frac{(2k3)^{m1}}{_{j=5}^{2m+1}(2kj)}$$
(25)
($`j`$ odd). Note that all hyperflatness factors are independent of $`\beta `$.
## 5 Extracting $`q(r)`$ from experimentally measured structure functions
The great advantage of the hyperflatness factors is that they yield a direct way to estimate the $`r`$-dependent non-extensitivity parameter $`q(r)`$ from experimentally measured structure functions $`u^m(r)`$. Eq. (22) yields
$$F_2=\frac{6k9}{2k5}=\frac{159q}{75q}$$
(26)
or
$$k=\frac{5F_29}{2F_26},$$
(27)
equivalent to
$$q=\frac{7F_215}{5F_29}$$
(28)
(see also ).
Another relation follows from eq. (23).
$$k=\frac{1}{2F_330}\left(6F_345\pm \sqrt{F_3^2+120F_3}\right),$$
(29)
equivalent to
$$q=\frac{8F_375\pm \sqrt{F_3^2+120F_3}}{6F_345\pm \sqrt{F_3^2+120F_3}}$$
(30)
In fact, each hyperflatness factor $`F_m`$ with $`m2`$ yields a relation for $`k`$ (or $`q`$), and all relations are the same in case the ‘free’ turbulence theory is exact.
Given some experimentally measured hyperflatness structure functions $`F_m(r)`$ one can now determine the corresponding curves $`q(r)`$. The less these curves differ for the various $`m`$, the more precise is the zeroth-order (free) turbulence theory. For examples of such experimentally measured curves $`q(r)`$, see . Generally, the hyperflatness factors are complicated functions of both the Reynolds number $`R_\lambda `$ and the separation distance $`r`$, and so is $`q(R_\lambda ,r)`$. But one may conjecture that for $`R_\lambda \mathrm{}`$ one obtains a universal function $`q^{}(r)`$. How it looks like in the entire $`r`$-range is still an open question.
## 6 Some remarks on the scaling exponents $`\zeta _m`$
The precise values of the scaling exponents $`\zeta _m`$, which describe the scaling behaviour of the structure functions $`u^mr^{\zeta _m}`$ in the inertial range, are still a rather controversial topic in turbulence theory. In fact, it is not at all clear whether there is exact scaling at all or just approximate scaling, whether the higher moments exist at all or not, how reliable the experimentally measured higher-order exponents are, and what the effects of finite Reynolds numbers are. On the theoretical side, a variety of models and theories have been suggested (see, e.g., ) but a true breakthrough convincing a significant majority of scientists working in the field seems not to have been reached at the moment.
To ‘derive’ values for the scaling exponents using methods from non-extensive statistical mechanics, one needs additional assumptions— just as in all other models and theories dealing with the scaling exponents. In a logarithmic depence of $`k`$ on $`r`$ was suggested to derive the following formula for the scaling exponents
$$\zeta _m=\frac{m}{3}\left\{1\mathrm{log}_2\left(1\frac{3}{2\overline{k}1}\right)\right\}+\mathrm{log}_2\left(1\frac{m}{2\overline{k}1}\right).$$
(31)
Here $`\overline{k}=1/(q1)`$ is the average value of the non-extensitivity parameter in the inertial range. On the other hand, Kolmogorov’s lognormal model (the K62-theory ) predicts
$$\zeta _m=\frac{m}{3}\left(1+\frac{\mu }{2}\right)\frac{1}{18}\mu m^2.$$
(32)
where $`\mu `$ is the intermittency parameter. Now, asuming that $`\overline{k}`$ is rather large we can expand the logarithms in eq. (31). The linear terms cancel, and the first non-trivial terms are the quadratic ones. Neglecting the higher-order cubic terms, one precisely obtains from eq. (31) the result (32), identifying
$$\mu =\frac{18}{\mathrm{ln}2(2\overline{k}1)^2}.$$
(33)
It is encouraging that the simplest non-extensive model assumptions lead to an old theory that has a long tradition in turbulence theory (though it is known that K62 cannot be correct —it is, however, a good approximation for $`m`$ not too large). The value of the intermittency parameter comes out with the correct order of magnitude ($`\mu 0.2`$) if one chooses $`\overline{k}6`$, and the same $`k`$ fits the experimentally measured probability densities correctly.
On the other hand, starting in the derivation leading to eq. (31) from a different dependence of $`k`$ on $`r`$ rather than the logarithmic one, one also obtains different formulas for the scaling exponents. So the final answer to the question of the scaling exponents is still open.
Actually, in eq. (31) was derived using the ‘free’ turbulence theory, but clearly the correct theory is the infinite-order theory. As one can easily see, the infinite-order theory yields probability densities $`p(u)`$ living on a compact support for arbitrary small but finite $`\gamma \tau `$. Indeed, iterating eq. (6) one obtains $`u_n=\lambda ^nu_0+\sqrt{\gamma \tau }_{j=1}^n\lambda ^{nj}x_j`$. Hence, for any chaotic dynamics $`T`$ bounded on some finite phase space (say \[-1,1\]) one obtains for $`n\mathrm{}`$ from the geometric series the rigorous bound $`|u|\sqrt{\gamma \tau }/(1\lambda )1/\sqrt{\gamma \tau }`$. Thus for the infinite-order theory all moments exist for arbitrary small but finite $`\gamma \tau `$ —in contrast to the 0-th order theory, where only moments $`u^m`$ with $`m<2k1`$ exist, since here the densities live on a non-compact support and there is polynomial decay for large $`|u|`$. This once again shows how delicate the problem of the existence of the moments and of the scaling exponents in general is. For densities living on a compact support, as provided by the infinite-order theory, one expects a linear asymptotics of $`\zeta _m`$ for large $`m`$.
Generally, the way non-extensive statistics can be used to ‘derive’ the scaling exponents $`\zeta _m`$ is not unique. An alternative approach, based on an extension of the multifractal model and $`q`$-values smaller than 1, has been suggested by T. and N. Arimitsu . In their model asymptotically there is a logarithmic dependence og $`\zeta _m`$ on $`m`$.
## Acknowledgement
A large part of this research was performed during the author’s stay at the Institute for Theoretical Physics, University of California at Santa Barbara, supported in part by the National Science Foundation under Grant No. PHY94-07194. The author also gratefully acknowledges support by a Leverhulme Trust Senior Research Fellowship of the Royal Society.
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# Chapter 1 Disordered Quantum Solids
## Chapter 1 Disordered Quantum Solids
T. Giamarchi, E. Orignac
## Abstract
Due to the peculiar non-fermi liquid of one dimensional systems, disorder has particularly strong effects. We show that such systems belong to the more general class of disordered quantum solids. We discuss the physics of such disordered interacting systems and the methods that allows to treat them. In addition to, by now standard renormalization group methods, We explain how a simple variational approach allows to treat these problems even in case when the RG fails. We discuss various physical realizations of such disordered quantum solids both in one and higher dimensions (Wigner crystal, Bose glass). We investigate in details the interesting example of a disordered Mott insulator and argue that intermediate disorder can lead to a novel phase, the Mott glass, intermediate between a Mott and and Anderson insulator.
### 1 Introduction
Disorder effects are omnipresent in condensed matter physics, where one has to struggle very hard to deal with clean systems. Quite remarkably fermionic system exhibit marked differences with disorder in classical systems. Indeed the very existence of a Fermi energy $`E_F`$ “reduces” the effects of disorder since the relevant parameter now become the relative strength of the disorder compared to the Fermi energy $`D/E_F`$ or the mean free path compared to the Fermi length $`k_Fl`$. Of course nature would not remain as simple as that, and quantum effects lead in fact also to reinforcement of the disorder effects and turn in low dimension a free electron system into an insulator, as pointed out by Mott and Twose.
We have now gained a very good understanding of the properties of such disordered free electron systems. To tackle them an arsenal of methods ranging from diagrams, scaling theory, replicas, supersymmetry have been developped. Life become much less simple when interactions among fermions are taken into account. Although it is very intuitive to think that when interactions are small the noninteracting problem is a good starting point, such an intuition turns out to be wrong for a number of reasons: (i) even if in the pure system interactions can be “removed” from the system by resorting to Fermi Liquid theory this is not the case when disorder is present. Because disorder renders electrons slowly diffusive slowly rather than ballistic, they feel the interactions much more strongly with explosive results. Effective interactions increase when looking at low energy properties and Fermi liquid theory breaks down. The non-interacting physics is thus only relevant for very low disorder. (ii) When the dimension is small or the interactions strong to start with (like in systems undergoing Mott transitions) it is of course impossible to start from the non-interacting limit and one has to solve the full problem. What puts us at a disadvantage here is that most of the techniques useful for the noninteracting case also fail as soon as interactions are included: (i) supersymmetric method, which rests by construction on the quadratic nature of the hamiltonian is useless. (ii) it is now difficult now to really determine classes of diagram to sum.
Disorder can still be averaged over by using replicas but then one is left with a complicated (and untractable) theory. Renormalisation group attempts have been made with some success but also with the problem that the coupling constants diverge, so that the low energy fixed point remains elusive. Many other attempts in treating this very complicated problem exist in the literature and it is impossible to list them all. Note that numerical studies are also hampered when studying this problem both because of the difficulty in taking the interactions into account (which more or less imposes either Monte-Carlo or exact diagonalization) and then to perform the complicated disorder averaging with enough statistics or large system sizes to get reliable results.
A very peculiar situation occurs when one considers one dimension. On one hand one expects the difficulty to be maximum here. The interactions lead to very strong effects and destroy any trace of Fermi liquid giving rise to what is known as a Lutinger liquid. The disorder is also extremely strong, giving rise for the noninteracting case to a system so localized that the localisation length is simply the mean free path and a diffusive regime is absent. On the other hand one is in $`d=1`$ in a much better situation to tackle the problem since the interactions can be treated essentially exactly using for example techniques such as bosonization so one only needs good techniques to tackle the disorder. Quite interestingly the physics of interacting and disordered one dimensional systems is the one of disordered quantum solids. Higher dimensional examples of such systems are the Wigner crystal, Charge Density Waves and Bose glass. In these notes, we will explain the techniques allowing to treat such systems, ranging from simple physical arguments to a quite sophisticated variational approach. In order to remain pedagogical and keep the algegra simple we mostly discuss the technicalities on the simplest example of spinless fermions. We briefly discuss the specific physical realizations and give references so that the reader can look in more details the physical properties of these specific systems.
Before we embark with the physics, let us point out that these notes results from the synthesis of various lectures. Some arbitrary choice of material had to be made in order to keep some level of clarity. Even if we have made some effort to cover various interesting topics, these notes cannot pretend to be as exhaustive as a full review. We thus apologize in advance to anybody whose pet problem (or paper) is not covered in these few pages.
### 2 Disordered interacting Fermions
#### 2.1 Model
Let us consider spinless fermions hopping on a lattice with a kinetic energy $`t`$ and an interaction $`V`$
$`H`$ $`=`$ $`t{\displaystyle \underset{i}{}}(c_i^{}c_{i+1}+h.c.)`$
$`+`$ $`{\displaystyle \underset{i>j}{}}V_{ij}(n_i\overline{n})(n_j\overline{n})`$
where $`n_i`$ ($`\overline{n}`$) is the local (average) electron density, and the rest of the notation is standard ($`t=1`$ is the unit of energy). For nearest neighbor interactions this is the well known $`tV`$ model. When using a Jordan-Wigner transformation to express the fermions in term of spins this latter model maps to an XXZ spin chain. In addition to (2.1) we want to submit the fermions to a disorder. We concentrate here on site disorder. Again most results/methods will carry over for randomness in the hopping. The randomness is simply
$$H_{\mathrm{int}}=\underset{i}{}\mu _i(n_i\overline{n})$$
(1.2)
where $`\mu _i`$ is a random variable of zero mean.
#### 2.2 Pure system
The low energy properties of the pure system (2.1) are by now well understood. We will thus review the bosonization of (2.1) only very briefly to fix the notations and refer the reader to the various reviews on the subject . To get the low energy physics, it is enough to focus on the excitations aroung the Fermi points. In the continuum limit $`ix`$ this amounts to express the fermion field in term of slow varying (with respect to the lattice spacing) field of right (with momentum close to $`+k_F`$) and left (with momentum close to $`k_F`$) movers
$$c_i^{}\psi ^{}(x)=e^{ik_Fx}\psi _+^{}(x)+e^{ik_Fx}\psi _{}^{}(x)$$
(1.3)
In term of these fields (2.1) becomes
$$H=i\mathrm{}v_F𝑑x(\psi _+^{}_x\psi _+(x)\psi _{}^{}_x\psi _{}(x))+𝑑x_1𝑑x_2V(x_1x_2)\rho (x_1)\rho (x_2)$$
(1.4)
where the density reads
$$\rho (x)=\psi _+^{}\psi _++\psi _{}^{}\psi _{}+(e^{i2k_Fx}\psi _+^{}\psi _{}+\mathrm{h}.\mathrm{c}.)$$
(1.5)
The remarkable feature in $`d=1`$ is that all the excitations of the system can be reexpressed in term of the fluctuations of density. If one introduces a field $`\varphi (x)`$ describing the long wavelength part of the density (1.5) reads
$$\rho (x)=\rho _0\frac{1}{\pi }\varphi (x)+\frac{1}{(2\pi \alpha )}(e^{i2k_Fx+2\varphi (x)}+\mathrm{h}.\mathrm{c}.)$$
(1.6)
where $`\alpha `$ is a lattice spacing, $`\rho _0=\overline{n}/\alpha `$ and:
$`\varphi (r)`$ $`=`$ $`{\displaystyle \underset{\pm }{}}\psi _\pm ^{}(x)\psi _\pm (x)`$ (1.7)
$`\theta (r)`$ $`=`$ $`{\displaystyle \underset{\pm }{}}\pm \psi _\pm ^{}(x)\psi _\pm (x)`$ (1.8)
$`\theta `$ is a similar field but associated with the long wavelength part of the current. $`\varphi `$ and $`\mathrm{\Pi }=\frac{1}{\pi }\theta `$ are canonically conjugate. Both the Hamiltonian and the fermion operator can be expressed in terms of these two fields
$`H`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle 𝑑x\left[uK(\theta (x))^2+\frac{u}{K}(\varphi (x))^2\right]}`$ (1.9)
$`\psi _\pm (r)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\pi \alpha }}}e^{i(\pm \varphi (r)\theta (r))}`$ (1.10)
Using (1.10) in (1.5) one easily recovers (1.6). For the free Hamiltonian (2.1) (with $`V=0`$) one has in (1.9) $`u=v_F`$ and $`K=1`$. What makes the boson representation so useful is the fact that even in the presence of interactions the bosonized form of (2.1) remains (1.9) but with renormalized (Luttinger Liquid) parameters $`u`$ and $`K`$.
#### 2.3 Disorder
Similarly the disorder can be rewritten in the boson representation. It is natural to separate the Fourier components with wavevectors close to $`q0`$ ($`\eta `$) and $`q\pm 2k_F`$ ($`\xi ,\xi ^{}`$). The disorder becomes thus
$$H_{\mathrm{dis}}=\eta (x)[\psi _+^{}\psi _++\psi _{}^{}\psi _{}]+(\xi (x)\psi _+^{}\psi _{}+\mathrm{h}.\mathrm{c}.)$$
(1.11)
where $`\eta `$ and $`\xi `$ are now two independent random potentials. $`\eta `$ is real but when the system is incommensurate ($`2k_F\pi `$) $`\xi `$ is complex. The potentials $`\eta `$ and $`\xi `$ correspond respectively to the forward scattering and the backward scattering of the fermions on impurities. (1.11) reads in the boson representation
$$H_{\mathrm{dis}}=\eta (x)\frac{1}{\pi }\varphi (x)+[\frac{\xi (x)}{(2\pi \alpha )}e^{i2\varphi (x)}+\mathrm{h}.\mathrm{c}.]$$
(1.12)
For incommensurate systems the forward scattering can easily be absorbed by a gauge transformation on the fermions. In the boson language this amounts to shift $`\varphi \varphi \frac{2K}{u}^x𝑑y\eta (y)`$. Since $`\xi `$ is complex and hence has a random phase this shift does not affect the backscattering term (this will be even more explicit on the replicated Hamiltonian (1.47) below). Thus the forward scattering can be accounted for completely. Its main effect is to lead to an exponential decay of the disorder averaged density correlation functions. Since the current and the superconducting correlation functions depends only on $`\theta `$ (or equivalently in an action representation on $`_t\varphi `$) they are not affected by the shift. Transport properties thus depends on the backscattering alone, as is obvious on physical basis since forward scattering cannot change the current. For commensurate potentials the situation is mode complicated since $`\xi `$ is real ($`e^{+i2k_F}=e^{i2k_F}`$) and the forward scattering now does affect the backscattering term (in other words by combining a forward scattering and a scattering on the lattice one can generate a backscattering term). This has consequences that we will examine in more details in section 7. For the moment let us focus on the incommensurate case and get rid of the forward scattering.
### 3 Tackling the disorder
Disordered one dimensional particles are thus described by
$$S/\mathrm{}=dxd\tau [\frac{1}{2\pi K}[\frac{1}{v}(_\tau \varphi )^2+v(_x\varphi )^2]+\frac{\xi (x)}{2\pi \alpha \mathrm{}}e^{i2\varphi (x)}+\mathrm{h}.\mathrm{c}.]$$
(1.13)
where we have rewritten (1.9) and (1.11) as an action. We have reintroduced $`\mathrm{}`$ and other pesky constants to show explicitly the various physical limits. Note that although we are mainly concerned here with fermions (1.9) describes in fact nearly every one dimensional disordered problem ranging from dirty bosons to disordered spin chains, since all this problems have essentially the same boson representation. We will examine these different systems in more details in section 4. (1.13) emphasizes the physics of the problem. The electron system can be viewed as a charge density wave since the density varies as (1.6)
$$\rho (x,\tau )=\rho _0\mathrm{cos}(2k_Fx2\varphi (x,\tau ))$$
(1.14)
The elastic term in (1.13) wants the phase of the density wave to be as constant as possible, and have a nice sinusoidally modulated density. The disorder term on the contrary wants to pin this charge density on the impurities by distorting the phase, as shown on Figure 1.
The problem of localization of interacting fermions is thus very similar to the one of the pinning of classical charge density waves. The charge density wave is here intrinsic to the one-dimensional interacting electron gas and not due to a coupling to phonons. The main features are nevertheless similar, the main difference being the fact that the effective mass of the “CDW” is much smaller in the absence of the electron-phonon coupling and hence the importance of the quantum fluctuations is much higher. Various quantities of special interest have a quite simple expression in terms of the bosons. The conductivity is simply given by
$$\sigma (\omega )=\frac{i}{\omega }\left[_\tau \varphi _\tau \varphi _{i\omega _n,q=0}\right]_{i\omega _n\omega +i\delta }$$
(1.15)
whereas the compressibility is given by a similar correlation function but with differents limits
$$\kappa =\frac{1}{\pi ^2}_x\varphi _x\varphi _{i\omega _n=0,q}$$
(1.16)
The correlations of the $`2k_F`$ part of the density and of the superconducting order parameter $`\psi _+^{}\psi _{}^{}`$ are respectively given by
$`\chi _\rho (r)`$ $`=`$ $`e^{i2\varphi (r)}e^{i2\varphi (0)}`$ (1.17)
$`\chi _s(r)`$ $`=`$ $`e^{i2\theta (r)}e^{i2\theta (0)}`$ (1.18)
where $`r=(x,\tau )`$ and $`\pi \mathrm{\Pi }(x)=\theta (x)`$.
Given the highly nonlinear nature of the coupling to disorder, (1.13) is quite tough to solve. On the other hand, if $`\xi (x)`$ was just a constant we would have a simple sine-Gordon theory for which a great deal is known. We will thus tackle the problem by increasingly sophisticated methods.
#### 3.1 Chisel and Hammer
In the absence of quantum fluctuations, $`\varphi `$ would be a classical field and we would have a good idea of what happens. This is the way Fukuyama and Lee looked at this problem. Such an approximation is of course very good for “phononic” charge density waves since the quantum term is $`\mathrm{\Pi }^2/(2M)`$ and thus very small. For fermions this corresponds to the “classical” limit $`\mathrm{}0`$, $`K0`$ keeping $`\overline{K}=K/\mathrm{}`$ fixed, and thus to very repulsive interactions. In that case we can ignore all quantum fluctuations, and look for a static solution for $`\varphi `$. It is of course crucial for the existence of such solution that the disorder does not depends on time. This solution $`\varphi _0(x)`$ describes the static distortion of the phase imposed by the random potential. In the absence of kinetic energy $`(\varphi )^2`$ , it would be easy to “determine” $`\varphi _0`$. If we write the random field $`\xi `$ as an amplitude $`|\xi (x)|`$ and a random phase $`2\zeta `$, the disorder term writes
$$𝑑x|\xi (x)|e^{i2(\varphi (x)\zeta (x))}+\mathrm{c}.\mathrm{c}.$$
(1.19)
The optimum is thus for $`\varphi _0(x)`$ to follow the random phase on each point. For point like impurities located on random positions $`R_i`$, $`|\xi |`$ would just be the strength of each impurity potential and $`\zeta =k_FR_i`$. Thus $`\varphi _0(x)=\zeta (x)`$ is the generalisation to any type of disorder (and in particular to the Gaussian disorder so dear to the theorist) of the physics expressed in Figure 1: get the density minimum at each impurity. In presence of kinetic energy following the random phase would cost too much kinetic energy. We do not know exactly now how to determine the optimal $`\varphi _0(x)`$ but we can do some scaling arguments. Let us assume that $`\varphi `$ remains constant for a lengthscale $`L_{\mathrm{loc}}`$. On this lengthscale $`\varphi `$ takes the value that optimizes the disorder term, which now reads
$$E_{\mathrm{dis}}=\left[_0^{L_{\mathrm{loc}}}\xi (x)\right]e^{i2\varphi }+\mathrm{c}.\mathrm{c}.$$
(1.20)
If we take for example a gaussian distribution for $`\xi `$
$$\overline{\xi (x)\xi ^{}(x^{})}=𝒟\delta (xx^{})$$
(1.21)
one gets because of the average of a complex random variable on a box of size $`L_{\mathrm{loc}}`$ that the disorder contributes as
$$E_{\mathrm{dis}}=\sqrt{𝒟L_{\mathrm{loc}}}e^{i(2\varphi _02\mathrm{\Xi })}$$
(1.22)
where $`\mathrm{\Xi }`$ is some phase. It clear that the optimum energy is reached if $`\varphi _0`$ adjusts to this (now unknown) phase. The global energy gain now scales as $`\sqrt{L_{\mathrm{loc}}}`$. Between two segments of size $`L_{\mathrm{loc}}`$ the phase has to distort to reach the next optimal value. The distortion being of the order of $`2\pi `$ the cost in kinetic energy reads
$$E_{\mathrm{kin}}\frac{1}{L_{\mathrm{loc}}}$$
(1.23)
minimizing the total cost shows that the length over which $`\varphi _0`$ remains constant is given by
$$L_{\mathrm{loc}}\left(\frac{1}{𝒟}\right)^{\frac{1}{3}}$$
(1.24)
This tells us that the system does pin on the impurities and that below $`L_{\mathrm{loc}}`$ the system looks very much like an undistorted system. Since at the scale $`L_{\mathrm{loc}}`$, $`\varphi _0`$ varies randomly the $`2k_F`$ density density correlations will decay exponentially with a characteristic size $`L_{\mathrm{loc}}`$. It is thus very tempting to associate $`L_{\mathrm{loc}}`$ with the Anderson localization length. Note that for the free fermion point $`L_{\mathrm{loc}}1/𝒟`$ instead of (1.24), so the above formula is clearly missing a piece of physics when $`K`$ is not zero. Nevertheless from this simple scaling argument we have obtained: (i) the fact that classical CDW or very very repulsive fermions are pinned (localized) by disorder; (ii) the localization length; (iii) the fact that the ground state should contain a static distortion of the phase due to the disorder. Unfortunately we have no other information on $`\varphi _0`$, which is certainly a drawback.
Even with our limited knowledge of the statics we can nevertheless try to extract the dynamics. Let us assume that all deformations of the phase which are not contained in the static distortion are small and thus that we can write
$$\varphi (x,\tau )=\varphi _0(x)+\delta \varphi (x,\tau )$$
(1.25)
with $`\delta \varphi (x,\tau )\varphi _0(x)`$ in a very vague sense since we deal with random variables. One can try to expand the random term in power of $`\delta \varphi `$
$`S_{\mathrm{dis}}`$ $`=`$ $`{\displaystyle 𝑑\tau 𝑑x|\xi (x)|\mathrm{cos}(2(\varphi (x,\tau )\zeta (x)))}`$
$``$ $`2{\displaystyle 𝑑\tau 𝑑x|\xi (x)|\mathrm{cos}(2(\varphi _0(x)\zeta (x)))(\delta \varphi (x,\tau ))^2}`$
One can thus use in principle (3.1) to compute the various physical quantities. Note that the conductivity (1.15) will not depend directly on the statics solution $`\varphi _0`$ since $`_t\varphi _0=0`$, so we can hope to compute it. Of course the dependence of the fluctuations $`\delta \varphi `$ in $`\varphi _0`$ is hidden in (3.1). If $`\varphi _0`$ was following the random phase at every point, then the disorder term would just lead to a mass term for the fluctations and the optical conductivity would show a gap. In fact this is not true at every point so (3.1) leads to a distribution of masses for the fluctuations. Unfortunately the knowledge of $`\varphi _0`$ is too crude to compute the conductivity accurately and depending on what exactly is $`L_{\mathrm{loc}}`$ one can find either a gap, a non analytic behavior or a $`\sigma (\omega )\omega ^2`$ behavior at small frequencies. Based on physical intuition Fukuyama and Lee opted for the later, but the method shows its limitations here and does not allow a reliable calculation of the physical quantities. More precise calculations of $`\varphi _0`$ and the conductivity can be performed in the classical limit $`K0`$ using a transfer matrix formalism .
One thing that can be obtained from the partial knowledge of the fluctuations is the effect of quantum fluctuations. Indeed when quantum fluctuations are present the expansion (3.1) is not valid any more since the cosine should be normal ordered before it can be expanded. This leads to
$$\mathrm{cos}(2(\varphi _0\zeta +\delta \varphi ))2\mathrm{cos}(2(\varphi _0\zeta ))e^{2\delta \varphi ^2}(\delta \varphi )^2$$
(1.27)
where the average $``$ has now to be computed self-consistently using (1.27). This leads to a modified localisation length of the form
$$L_{\mathrm{loc}}\left(\frac{1}{𝒟}\right)^{\frac{1}{32K}}$$
(1.28)
This expression for the localization length suggests that a delocalization transition is induced by the quantum fluctuations and occurs at $`K=3/2`$. In the fermion language this corresponds to extremely attractive interactions.
#### 3.2 Starting from the metal: RG
The previous method starts directly from the localized phase. It provides some limited information about this phase, but suffers from serious limitations. An alternative approach is to start from the pure Luttinger liquid and investigate the effects of disorder perturbatively, and build a renormalization group analysis. The RG provides us with the best possible description of the delocalized phase and the critical properties of the transition. It also gives a very accurate description of the localized phase up to lengthscales of the order of the localization length $`L_{\mathrm{loc}}`$. Here again we describe the method for simplicity on spinless fermions and discuss more complex systems in section 4.
To have a hint of the RG equations let us expand the disorder term (1.13) to second order. This leads to
$$𝑑\tau 𝑑x𝑑\tau ^{}𝑑x^{}\xi (x)\xi ^{}(x^{})e^{2i(\varphi (r)\varphi (r^{}))}$$
(1.29)
It is easy to see that at the tree level, (1.29) scales as $`𝒟L^{32K}`$. This leads to the scaling of the disorder
$$\frac{𝒟}{dl}=(32K)𝒟$$
(1.30)
This traduces in fact the dressing of the scattering on the disorder by the interactions and has been derived using either diagrams or RG . In itself it seems to confirm the result of (1.28) i.e. the existence of a transition. Note that the advantage of the bosonization derivation is to allow to reach the non perturbative point in interactions where such a metal-insulator transition would take place.
In fact (1.30) would not allow in itself to really determine the metal-insulator transition point. This can be seen by using the RG to compute the finite temperature (or finite frequency) conductivity of the system . The idea is simply to renormalize until the cutoff is of the order of the thermal length $`l_Tu/T`$ corresponding to $`e^l^{}l_T/\alpha `$. At this length scale the disorder can be treated in the Born approximation. As the conductivity is a physical quantity it is not changed under renormalization and we have:
$`\sigma (n(0),D(0),0)=\sigma (n(l),D(l),l)=\sigma _0{\displaystyle \frac{n(l)D(0)}{n(0)D(l)}}=\sigma _0{\displaystyle \frac{e^lD(0)}{D(l)}}`$ (1.31)
where $`\sigma (n(l),D(l),l)=\sigma (l)`$ and $`n(l)`$ are respectively the conductivity and the electronic density at the scale $`l`$. $`\sigma _0=e^2v_F^2/2\pi \mathrm{}𝒟`$ is the conductivity in the Born approximation, expressed with the initial parameters. Using (1.30) one gets from (1.31)
$$\sigma (T)\frac{1}{𝒟}T^{22K}$$
(1.32)
This result is the direct consequence of the renormalization of the scattering on impurities by interactions (1.30). One immediately sees that (1.32) alone would lead to a paradox since (1.30) gives a localized-delocalized boundary at $`K=3/2`$ whereas (1.32) gives perfect conductivity above $`K=1`$ (i.e. the noninteracting point). One could also immediately see that if one introduces a new variable such as
$$\stackrel{~}{𝒟}=e^{al}𝒟$$
(1.33)
the dimension of such a variable would be $`(3a2K)`$, leaving the location of the transition point as determined from (1.30) quite arbitrary. Although such a transformation seems arbitrary, if one considers that the disorder stems from impurities with a concentration $`n_i`$ and a strength $`V`$, the limit of Gaussian disorder corresponds simply in taking $`n_i\mathrm{}`$ with $`V0`$ keeping $`𝒟=n_iV^2`$ fixed. Thus the choice $`a=1`$ in (1.33) simply corresponds to $`\stackrel{~}{𝒟}=V^2`$ i.e. writting an RG equation for the impurity strength.
The answer to this simple paradox is of course that (1.30) should be complemented by another RG equation. In addition to renormalizing $`𝒟`$ (1.29) generates as well quadratic terms that renormalize the free part of the Hamiltonian, i.e. the velocity $`v`$ and the Luttinger parameter $`K`$. Details can be found in . The main equation is the renormalization of the Luttinger parameter $`K`$ and reads
$$\frac{K}{dl}=K^2𝒟/2$$
(1.34)
This equation describes the renormalization of the interactions by the disorder. Both RG equations (1.30) and (1.34) have a diagramatic representation shown on figure 2.
Using the flow (1.30) and (1.34) one can easily check that two phases exist as shown on figure 3.
For large $`K`$ one in the delocalized phase where the disorder is irrelevant and the system is a Luttinger liquid with renormalized coefficients $`u^{}`$ and $`K^{}`$. All correlation functions decay as power laws, and because $`K^{}>3/2`$ the system is dominated by superconducting fluctuations. On the transition line the exponent flows to the universal value $`K^{}=3/2`$ Below this line $`𝒟`$ flows to large values, indicating that the disorder is relevant. This phase is the localized phase. This is obvious for physical reasons but can also be guessed from the exact solution known for the noninteracting line $`K=1`$ (and any $`𝒟`$) which belongs to this phase. As can be seen from (1.30) and (1.34) the transition is Berezinskii-Kosterlitz-Thouless (BKT) like in the $`K`$, $`\sqrt{𝒟}`$ variables.
In addition to the phase diagram itself a host of physical properties can be extracted from the RG. The simplest one is the localization length. One can use that for $`𝒟(l)1`$ the localization length is of the order of the (renormalized) lattice spacing $`\alpha e^l`$. The full determination needs an integration of both (1.30) and (1.34). Close to the transition the divergence of the localization length is BKT like (setting $`K=3/2+\eta `$)
$$L_{\mathrm{loc}}e^{2\pi /\sqrt{9𝒟\eta ^2}}$$
(1.35)
Deep in the localized regime, and for weak disorder, a good approximation is to neglect the renormalization of $`K`$ in (1.30). A trivial integration of (1.30) then gives back (1.28). This we see that the SCHA calculation corresponds in fact, both for the phase diagram and for the localization length to the limit of infinitesimal disorder.
Out of the RG one can also extract, using (1.31) the behavior of the temperature or frequency dependence of the conductivity. In the localized phase this can only be used up to the energy scale corresponding to the localization length i.e. $`E_{\mathrm{pin}}=k_BT_{\mathrm{pin}}=\mathrm{}\omega _{\mathrm{pin}},\omega _{\mathrm{pin}}=v/L_{\mathrm{loc}}`$. Below this lengthscale another method than the weak coupling RG should be used. We will come back to that point in section 6. Here again, although the full flow should be taken into account one can get an approximate formula by ignoring the renormalization of $`K`$, which leads to (1.32). For $`K<3/2`$ (including the noninteracting point) any small but finite disorder grows, renormalizing the exponents and ultimately leading to a decrease of the conductivity, even if one started initially from $`K>1`$. A very crude way of taking into account both equations (1.30) and (1.34) would be to say that one can still use (1.32) but with scale dependent exponents (see )
$$\sigma (T)T^{22K(T)}$$
(1.36)
This renormalization of exponents and the faster decay of conductivity is in fact the signature of Anderson localization. The equivalent frequency dependence is
$$\sigma (\omega )\omega ^{2K4}$$
(1.37)
Here again this formula break down below the scale $`\omega _{\mathrm{pin}}v/L_{\mathrm{loc}}`$ which is the pinning frequency. Similarly below $`L_{\mathrm{loc}}`$ correlation functions can again be computed using the RG, but of course the asymptotic behavior cannot be obtained.
### 4 Other systems and RG
Despite its limitations to physics above $`E_{\mathrm{pin}}`$ the RG is an extremely efficient method given its simplicity. It allows in addition a perfect description of the delocalized phase and of the critical behavior, something unatainable through the methods of section 3.1 and allows for interesting extensions.
First let us note that the equations (1.30) and (1.34) also describe the case of a single impurity. Indeed in that case one can go back to the definition $`𝒟=n_iV^2`$ and take the limit $`n_i0`$, i.e. $`𝒟0`$. (1.34) shows that in that case $`K`$ cannot be renormalized since a single impurity cannot change the thermodynamic behavior. Only (1.30) remains, leading directly to temperature dependence of the form (1.32), and a localized-delocalized transition at $`K=1`$. More details on such a relation between the two problems and on the remaining open question can be found in .
Quite remarkably the set (1.30-1.34) seems wrong. Indeed $`K`$ naively depends on the (inelastic) interations. Perturbatively, for the pure systems $`K=1V/(2\pi v_F)`$. If one start for $`K=1`$, i.e. for the non interacting system, it would thus seem from (1.34) that the elastic scattering on the impurities can generate inelastic fermion-fermion interactions. The solution of this paradox is hidden in the precise way the RG procedure is build. In order to have the elastic nature of the scattering on impurities, the time integrations in (1.29) should be dome independently for $`\tau `$ and $`\tau ^{}`$. When one performs the RG one introduces a cutoff and imposes $`|\tau \tau ^{}|>\alpha `$. Thus a part is left out of (1.29) which is
$`𝒟{\displaystyle 𝑑x_{|\tau \tau ^{}|<\alpha }𝑑\tau 𝑑\tau ^{}\rho (x,\tau )\rho (x,\tau ^{})}`$
$`2𝒟\alpha {\displaystyle 𝑑x𝑑\tau \rho (x,\tau )\rho (x,\tau )}`$ (1.38)
which is exactly an inelastic interaction term. Thus in fact $`K`$ contains not only the original inelastic interactions $`V`$ but also a small correction coming from the disorder itself. In order to determine the flow for $`V`$ it is thus necessary to take this small correction into account which gives the flow of Figure 3-b. One thus sees that the elastic case $`V=0`$ indeed remains elastic and also that for spinless fermions, the perturbative flow seems to indicate that the inelastic interactions are reduced by the disorder. This is compatible with the physical image that one would get at strong disorder: fermions localize individually and since the overlap of wavefunctions is exponentially small, so is the effect of interactions. One could thus naively expect that below $`L_{\mathrm{loc}}`$ the effect of interactions are strong but disappear above $`L_{\mathrm{loc}}`$. As we will see in section 6 the variational approach confirms this image. For fermions with spins (to be discussed below) we expect only the interaction in the charge sector to vanish. Interactions in the spin sector do remain and lead to random exchange.
Of course many more physical systems can be studied by this method. This is the case for spin chains, than can directly be mapped onto spinless fermions. A spin chain under a random magnetic field along $`z`$ is directly the problem that we solved in the previous section, with $`K`$ being the anisotropy ($`K=1`$ for an XY chain and $`K=1/2`$ for an Heisenberg one). However although for fermion problems it is unlikely that one is at a commensurate filling this is the natural situation for a spin chain (since the magnetization is zero in the absence of external field the filling of the equivalent spinless fermion system is $`n=1/2`$). A spin chain with random exchange is thus like a commensurate fermionic system with random hopping. We will come back to this peculiar case in section 7. Quite remarkably, in one dimension, bosons lead to physics similar to fermions. This is due to the fact that in one dimension statistics cannot be separated from interactions. Interacting bosons can thus be represented by a bosonization representation quite similar to the one of fermions. The density can be written as
$$\rho (r)=\rho _0\rho _0_x\varphi +\underset{n}{}e^{in(2\pi \rho _0x2\varphi (r))}$$
(1.39)
very similar to the fermionic form. while the single particle operator is now
$$\psi (r)\rho _0^{1/2}e^{i\theta }$$
(1.40)
(note the difference with the fermionic operator). The hamiltonian is still described by the quadractic form (1.13). Now $`K=\mathrm{}`$ for noninteracting bosons and $`K=1`$ for hard core ones. There is thus a transition between a superfluid state (for $`K>3/2`$) to a Bose glass state where the bosons are localized by the random potential at $`K=3/2`$.
In a similar way one can of course treat fermions with spins. Equations are more complicated since they involve the charge and spin sectors, and we will not discuss the full physics here but refer the reader to . Let us just draw attention to one interesting consequence of the renormalization equation of the disorder for the problem with spins, which reads
$$\frac{d𝒟}{dl}=(3K_\rho K_\sigma g_1)𝒟$$
(1.41)
where $`g_1`$ is the backscattering interaction between opposite spins. For a Hubbard type interaction $`g_1=U`$. For spin isotropic systems $`K_\sigma =1+g_1/(2\pi v_F)`$ and $`g_1`$ is marginal with a flow
$$\frac{dg_1}{dl}=g_1^2$$
(1.42)
For more general spin couplings either $`g_10`$ and $`K_\sigma K_\sigma ^{}`$, or $`g_1`$ is relevant and a spin gap opens. The physics of the localization transition depends thus on the sign of the interactions with special physical consequences. The transition point move from $`K_\rho =2`$ (for infinitesimal disorder) for repulsive interactions to $`K_\rho =3`$ for attractive ones. For a Hubbard type interaction for which $`1/2<K_\rho <2`$ the delocalization point can never be reached and the system is localized regardless of the strength of interaction and disorder. Another consequence is that in the presence of spin degrees of freedom the divergence of the localization length at the transition is not BKT like any more. One could naively think that in the repulsive case the $`g_1`$ term could be omitted and that the renormalization of the disorder could be written $`(3K_\rho K_\sigma ^{})`$. For spin isotropic interactions this misses an important part of the physics. Indeed, for a Hubbard interaction $`K_\rho =1U/(2\pi v_F)`$. Substituting in (1.41) leads (for the initial steps of the flow)
$$\frac{d𝒟}{dl}=(1\frac{U}{\pi v_F})𝒟$$
(1.43)
whereas the incorrect substitution at the fixed point would lead to $`(1+\frac{U}{\pi v_F})`$, leading to quite different physics. (1.43) implies that for Hubbard type interactions repulsive interactions make the system less localized than for attractive interactions, i.e. that the $`L_{\mathrm{loc}}^{U>0}>L_{\mathrm{loc}}^{U<0}`$. Similar effects exist for the charge stiffness and the persistent currents, i.e. for a system with spins the persistent currents are in fact enhanced by repulsive interactions. This counter-intuitive statement can be explained physically: interactions have two effects: (i) they tend to reinforce, when attractive the superconducting fluctuations in the system. This screens disorder and makes it less effective. This is the only effect occurring for spinless fermions. (ii) when spin degrees of freedom exists, repulsive interactions also tend to make the density more uniform by spreading the charge. This makes it more difficult to couple to disorder. These two effect compete. This has several consequences, in particular for mesoscopic systems. Of course, for fermions, true delocalization can only be achieved with attractive interactions reaching at least the nearest neighbor.
Many other systems have been treated by such method. there are one-dimensional systems with long-range $`1/r`$ interactions, that lead to a pinned Wigner crystal , doped spin 1 chains , fermionic and bosonic ladders, spin 1 chains in a magnetic disorder , spin ladders . Since we want to focus here on the methods we refer the reader to the above references for a detailed discussion of the physics of such systems.
### 5 A zest of numerics
Although we are mainly concerned about analytical method in these notes, let us mention some numerical results and methods that have been used in connection with the RG predictions. Although numerical studies have become very powerful in one dimension for pure systems the presence of disorder complicates matters. Three main methods have been used.
Exact diagonalizations, have been used to study both the phase diagram and the charge stiffness of both spinless fermions (or equivalently XXZ spin chains) and fermions with spins with short range or long range interactions . Using the finite size scaling of the spin stiffness $`\rho _s=\frac{1}{L}\frac{^2E(\theta )}{\phi ^2}=f\left(\frac{L}{L_{\mathrm{loc}}}\right)`$, where $`E(\theta )`$ is the ground state energy of the disordered XXZ chain with boundary conditions $`S_L^+=e^{i\phi }S_1^+`$, the localization length $`L_{\mathrm{loc}}`$ can be obtained , and is in good agreement with the RG results of section 3.2. The behavior of the correlation length close to the transition point appeared consistent with the predicted BKT-like behavior. The results also suggested that a finite disorder was needed to disorder the ground state for $`K>3/2`$. A similar study with systems sizes of up to $`L=18`$ sites was also made for XXZ spin chains with a random exchange in . Analysis of persistent currents was also in agrement with the RG prediction of section 4. Unfortunately the exact Diagonalization approach of the last section is limited to zero temperature and small system size.
In order to consider bigger system sizes, one can use Quantum Monte Carlo methods. In , such a study was performed for disordered bosons. The superfluid density was obtained as a function of interaction for a given disorder strength. It was shown that for not too repulsive equations, there was a phase with a finite superfluid density. For more repulsive interactions, a phase with finite compressiblity by zero superfluid density was obtained, in agreement with the Bose Glass theory of section 4.
The most promising recent method is the Density Matrix Renormalization Group. It been introduced in the recent years as a method specially designed to calculate the ground state of correlated one-dimensional systems. This method has been also applied to the problem of the XXZ chain in a random magnetic field parallel to the z axis by Schmitteckert et al. . The authors of have been able to consider system size of up to $`L=60`$ sites, and average over several hundred realizations of the disorder. Localization and phase diagram were also in good agreeement with the RG predictions.
Clearly, various numerical checks confirm the predictions of the RG. Unfortunately so far only the phase diagram, stiffness and localization length have been computed. This is clearly related to the complexity of the problem at hand. What would be extremely useful would be informations on quantities deep in the localized phase such as the single particle Green’s function, the ac or dc conductivity. Analysis of such quantities would nicely complement the RG analytical study and allow for comparison with other analytical techniques more suited for the localized phase such as the variational method we analyze in the next section.
### 6 Variational Method
Let us now study this problem using a completely different and at first sight more formal method. As usual it is very convenient to get rid of the disorder from the start. Given the non quadratic nature of (1.13) supersymmetric methods are unapplicable and we have to turn to replicas. The idea of the replica method in itself is quite simple. If we want to compute an observable $`O`$ we have to do both average over disorder and thermodynamic average
$$\overline{O}=𝒟Vp(V)O_V=DVp(V)\frac{𝒟\varphi O[\varphi ]e^{S_V[\varphi ]}}{D\varphi e^{S_V[\varphi ]}}$$
(1.44)
The action is usually linear in disorder and for Gaussian disorder the distribution of random potential is $`p(V)e^{{\scriptscriptstyle 𝑑xV(x)^2}}`$, so the average would be quite trivial without the denominator in (1.44). The idea is thus to introduce $`n`$ fields and to compute
$`{\displaystyle D\varphi _1D\varphi _2\mathrm{}D\varphi _nO[\varphi _1]e^{_{i=1}^nS_V[\varphi _i]}}=`$
$`{\displaystyle D\varphi O[\varphi ]e^{S_V[\varphi ]}\left[D\varphi e^{S_V[\varphi ]}\right]^{n1}}`$ (1.45)
which is exactly the quantity we want to average over disorder in (1.44) if one takes the formal limit $`n0`$. Since (6) has no denominator averaging over disorder is trivial. Of course there is a price to pay: before the averaging the replicas are all independent fields but the averaging introduces an interaction between them. We have thus traded a theory depending on a random variable $`V`$ but a single field for a theory without disorder but with $`n`$ coupled fields. Usually this is still a situation we are better equipped to solve because of the large number of field theoretic method dealing with “normal” (i.e. transitionally invariant actions). For the particular case (1.13) the replicated action is
$`S/\mathrm{}`$ $`=`$ $`{\displaystyle 𝑑x𝑑\tau \frac{1}{2\pi K}\underset{a}{}\left[\frac{1}{v}(_\tau \varphi _a)^2+v(_x\varphi _a)^2\right]}`$ (1.47)
$`{\displaystyle \frac{𝒟}{(2\pi \alpha )^2\mathrm{}}}{\displaystyle \underset{ab}{}}{\displaystyle 𝑑x𝑑\tau 𝑑\tau ^{}\mathrm{cos}(\varphi _a(x,\tau )\varphi _b(x,\tau ^{}))}`$
where $`a=1,\mathrm{},n`$ is the replica index. Disorder averaging has coupled the replicas via the cosine term. Because the disorder is time independent this coupling contains two fields that can be at arbitrary time and is thus highly non local. For fermions one usually prefers to go to frequency space, where this imply conservation of the frequency for each replica index, but this would not simplify things here because of the cosine.
This is up to now a totally formal procedure and nothing has been accomplished. (1.47) is totally equivalent to (1.13) and the difficulty is of course to solve it. Based on the RG equation (1.30) one could think naively that since the localized phase corresponds to $`𝒟\mathrm{}`$ it would be safe to expand the cosine term in (1.47). Unfortunately it is easy to check that fails seriously when $`n0`$ is taken (it of course works perfectly for a finite number of field $`n2`$). In order to circumvent this problem let us try to improve over this simple minded expansion of the cosine. Let us try a variational ansatz. We introduce a trial action $`S_0`$
$$S_0/\mathrm{}=\frac{1}{2\beta L\mathrm{}}\underset{q,\omega _n}{}\underset{ab}{}\varphi _a(q,\omega _n)G_{ab}^1(q,\omega _n)\varphi _b(q,\omega _n)$$
(1.48)
where the propagators $`G^1`$ are our variational “parameters”. As usual $`\frac{1}{L}_q\frac{dq}{(2\pi )}`$. If we introduce
$$Z=𝒟\varphi e^{S/\mathrm{}}$$
(1.49)
We then have the variational theorem for the free energy $`F=\mathrm{}\mathrm{log}(Z)`$
$$FF_{\mathrm{tr}}=F_0+SS_0_{S_0}$$
(1.50)
Since $`S_0`$ is quadratic, (1.50) can be in general computed quite explicitely as a function of the (unknown) propagators $`G`$. The “best” quadratic action $`S_0`$ is thus the one that satisfies the saddle point equations
$$\frac{F_{\mathrm{tr}}}{G_{ab}(q,\omega _n)}=0$$
(1.51)
which gives a set of integral equations allowing to determine the unknown functions $`G`$.
The observables are simply defined by quantities diagonal in replica indices as can be seen from (6). For some quantities such as the compressibility it is necessary to be more careful since one has to substract the average, which is usually zero in a pure system or after averaging over disorder but non zero for a specific realization of the disorder. Let us introduce the various propagators (time ordering in $`\tau `$ is always implied):
$`B_{ab}(x,\tau )`$ $`=`$ $`[\varphi _a(x,\tau )\varphi _b(0,0)]^2=`$ (1.52)
$`(G_{aa}(0,0)+G_{bb}(0,0)2G_{ab}(x,\tau ))`$
$`G_{ab}(q_x,\omega _n)`$ $`=`$ $`\varphi _a(q_x,\omega _n)\varphi _b(q_x,\omega _n)`$ (1.53)
The compressibility is given by
$`\chi (q,\omega _n)`$ $`=`$ $`{\displaystyle \frac{1}{\mathrm{}}}{\displaystyle }dx{\displaystyle _0^\beta \mathrm{}}d\tau e^{ı(qx\omega _n\tau )}\times `$ (1.54)
$`\times \overline{T_\tau (n(x,\tau )n(x,\tau ))(n(0,0)n(0,0))}`$
which leads to the average static compressibility $`\chi _s=lim_{q0}(lim_{\omega 0}\chi (q,\omega ))`$ (see (1.16)). When expressed in terms of the replicated bosonized operators (1.54) gives
$$\chi _s=\underset{q0}{lim}\underset{\omega 0}{lim}q^2G_c(q,\omega )$$
(1.55)
where we introduced an important propagator: the connected one defined as $`G_c^1(q)=_bG_{ab}^1(q)`$.
Without the replicas this method is nothing but the well known Self Consistent Harmonic Approximation (SCHA), which is known to work very well for sine-Gordon type Hamiltonians. Such a method gives in particular correctly the two phases (massless and massive). Extension of this method to disordered systems was done in the context of classical elastic systems such as interfaces . In quantum problems another level of complexity occurs because of the aforementioned non locality of the interaction in time. But before going to these problems, specific to the quantum systems, let us illustrate the aspects of this variational method when applied to disordered systems, on a technically simpler example (for which this method was extremely fruitful ): the case of classical periodic systems.
#### 6.1 A classical example
Let us take the action (1.47) but with only a single time integral for the disorder term. Such action would be the result of the average on a disorder both dependent on space and time. Of course such a disorder would be quite unrealistic for quantum problems. However (6.1) would be a perfectly natural Hamiltonian for a classical problem where $`z=v\tau `$ is now just one of the spatial dimensions . To make the analogy more transparent let us use $`z=v\tau `$, and replace the integral over $`x`$ by an integral in $`d1`$ dimensions. If denote by $`r`$ the $`d`$-dimensional space variable $`r=(x,z)`$ the starting action is
$`S/\mathrm{}`$ $`=`$ $`{\displaystyle d^dr\frac{1}{2\pi K}\underset{a}{}(_r\varphi _a)^2}`$
$``$ $`{\displaystyle \frac{𝒟}{(2\pi \alpha )^2\mathrm{}v}}{\displaystyle \underset{ab}{}}{\displaystyle d^dr\mathrm{cos}(2(\varphi _a(r)\varphi _b(r)))}`$
One can see that (6.1) is exactly the Hamiltonian describing an elastic system such as a vortex lattice or a classical CDW in the presence of point like defects in $`d`$ dimensions. $`\mathrm{}`$ plays the role of the temperature for the classical system, the elastic constant $`c`$ is given by $`c=1/(\pi \overline{K})`$ with $`\overline{K}=K/\mathrm{}`$ and $`\rho _0^2\mathrm{\Delta }/2=𝒟\mathrm{}/(2\pi \alpha )^2`$ would be the correlator of the classical disorder .
If we call $`q=(q_x,\omega )`$ the $`d`$-dimensional momentum, without loss of generality, the matrix $`G_{ab}^1(q)`$ can be chosen of the form $`G_{ab}^1=cq^2\delta _{ab}\sigma _{ab}`$. We obtain by minimization of the variational free energy the saddle point equations
$$G_c^1(q)=cq^2,\sigma _{ab}=\frac{2𝒟}{(\pi \alpha )^2}e^{2B_{ab}(r=0)}$$
(1.57)
Using (1.52-1.53) one obtains
$$B_{ab}(r)=\mathrm{}\frac{d^dq}{(2\pi )^d}(G_{aa}(q)+G_{bb}(q)2\mathrm{cos}(qr)G_{ab}(q))$$
(1.58)
For this particular problem, the connected part is not affected by disorder. This is the consequence of a hidden symmetry (statistical tilt symmetry) of (6.1), whose disorder part is not affected by any local shift of $`\varphi _a(r)`$ such as $`\varphi _a(r)\varphi _a(r)+f(r)`$, where $`f`$ is an arbitrary function. Such a symmetry does not exist for the time correlated disorder natural in a quantum problem, with important physical consequences on which we will come back. The only interesting equation here is thus the equation for the off-diagonal part $`\sigma _{ab}`$. Because of the locality of the interaction term between replicas in (6.1) the self energy $`\sigma _{ab}`$ is simply a matrix of constants.
Given the symmetry of the original action/Hamiltonian (6.1) by permutation of the replica indices, it is very natural to look for a variational ansatz with the same symmetry. This would mean that the $`G^1`$ matrix would have only (for each value of $`q`$) two independent values: the diagonal one $`G_{aa}=\stackrel{~}{G}`$ and the off diagonal one $`G_{ab}`$. Such a matrix can easily be inverted for any $`n`$, and the analytic continuation for $`n0`$ gives
$`G_c`$ $`=`$ $`\stackrel{~}{G}G_{ab}={\displaystyle \frac{1}{G_c^1}}`$ (1.59)
$`G_{ab}`$ $`=`$ $`{\displaystyle \frac{G_{ab}^1}{(G_c^1)^2}}`$ (1.60)
Using these inversion formulas is it easy to solve for (1.57). In $`d>2`$, $`B`$ depends on $`G_c`$ only and thus $`\sigma _{ab}`$ is simply a constant proportional to disorder. Given the gaussian nature of the trial action (1.48) the correlation functions such as the density-density can easily be computed
$$\overline{\rho (r)\rho (0)}=\rho _a(r)\rho _a(0)e^{2B_{aa}(r)}$$
(1.61)
Using (1.59) shows that $`G_{aa}(q)𝒟/q^4`$ leading to a growth
$$B_{aa}(r)𝒟r^{4d}$$
(1.62)
Although this solution is perfectly well behaved, a stability analysis of the replica symmetric saddle point shows that it is unstable. This can be checked from the eigenvalue $`\lambda `$ of the replicon mode .
$$\lambda =1\frac{8\mathrm{}𝒟}{(\pi \alpha )^2}e^{4\mathrm{}{\scriptscriptstyle {\scriptscriptstyle \frac{d^dp}{(2\pi )^d}}G_c(p)}}\frac{d^dq}{(2\pi )^d}G_c^2(q)$$
(1.63)
A negative eigenvalue $`\lambda `$ indicates an instability of the replica symmetric solution. We introduce a small regularizing mass in $`G_c`$: $`G_c(q)^1=cq^2+\mu ^2`$ and take the limit $`\mu 0`$. It is easy to see from (1.63) that for $`d<2`$ the replica symmetric solution is always stable. In that case disorder is in fact irrelevant, due to the strong quantum (or thermal for the associated classical system) fluctuations. For $`d=2`$ the condition becomes $`\mu ^{2(\overline{K}1)}<1`$ for small $`\mu `$. Thus there is a transition at $`\overline{K}=1`$ between a replica symmetric stable high temperature phase where disorder is irrelevant and a low temperature (glassy) phase where the symmetric saddle point is unstable. For the classical system this is the well known Cardy-Ostlund transition , with very interesting physical aspects of its own. This transition is the equivalent for time dependent disorder of the localization transition studied in section 3.2. For time dependent disorder one time integral drops in (1.29) and (1.30) would become $`(22K)𝒟`$ giving the transition at $`K=1`$. More details can be found in . For $`2<d<4`$ the replica symmetric solution is always unstable.
One should thus look for another way of inverting the $`0\times 0`$ matrices than the replica symmetric one. Fortunately such a scheme was invented in the context of spin glasses. Instead of having a single value $`\sigma `$ for the off diagonal term, one introduces a whole set of values. Let us briefly illustrate the procedure here, refering the reader to for details. Let us introduce a set of integers $`m_0=n`$, $`m_1`$, …, $`m_{k+1}=1`$ such that $`m_i/m_{i+1}`$ is an integer. One cuts the matrix in blocks of size $`m_1`$ as illustrated on Figure 4.
Elements outside the blocks have the value $`\sigma _0`$. The procedure is then recursively applied for the inner blocks. At the last step the value on the block of size $`m_k`$ is $`\sigma _k`$ and a diagonal value $`\stackrel{~}{\sigma }`$. Quite generally the matrix is thus now parametrized by a diagonal element and a whole function (in the limit $`n0`$) $`\sigma (u)`$, where $`u[0,1]`$ (notice the range of variation of the $`m_i`$ when $`n0`$). Such matrices can also be inverted in the limit $`n0`$), albeit with more complicated inversion rules than for the RS solution. When one has a continuous function (see Figure 5) the RSB is said to be continuous. Simpler cases are of course a constant function (a single off-diagonal value) which is simply the RS solution or a function continuous by steps. We have represented in Figure 4 (see also Figure 5) the case of a two-step RSB. We denote $`\stackrel{~}{G}(q)=G_{aa}(q)`$, similarly $`\stackrel{~}{B}(x)=B_{aa}(x)`$, and parametrize $`G_{ab}(q)`$ by $`G(q,v)`$ where $`0<v<1`$, and $`B_{ab}(x)`$ by $`B(x,v)`$. Physically, $`v`$ parametrises pairs of low lying states, in the hierarchy of states, $`v=0`$ corresponding to states further apart. The saddle point equations become:
$$\sigma (v)=\frac{2𝒟}{(\pi \alpha )^2}e^{2B(0,v)}$$
(1.64)
where
$$B(0,v)=2\mathrm{}\frac{d^dq}{(2\pi )^d}(\stackrel{~}{G}(q)G(q,v))$$
(1.65)
$`B(0,v)`$ corresponds physically to the mean squared phase fluctuations at the same point in space ($`r=0`$) (for the associated classical system this bould be mean squared relative displacements of the same object) but in two replica states, or more physically in two different low lying metastable states. The large distance behaviour of disorder-averaged correlators is determined by the small $`v`$ behaviour of $`B(0,v)`$. We look for a solution such that $`\sigma (v)`$ is constant for $`v>v_c`$, $`v_c`$ itself being a variational parameter, and has an arbitrary functional form below $`v_c`$. This corresponds to full RSB (see Figure 5). The algebraic rules for inversion of hierarchical matrices give:
$$B(0,v)=B(0,v_c)+_v^{v_c}𝑑w\frac{d^dq}{(2\pi )^d}\frac{2\mathrm{}\sigma ^{}(w)}{(G_c(q)^1+[\sigma ](w))^2}$$
(1.66)
where $`[\sigma ](v)=u\sigma (v)_0^v𝑑w\sigma (w)`$ and
$$B(0,v_c)=\frac{d^dq}{(2\pi )^d}\frac{2\mathrm{}}{G_c(q)^1+[\sigma ](v_c)}$$
(1.67)
This is a simple number. Taking the derivative of (1.64) with respect to $`v`$, using $`[\sigma ]^{}(v)=v\sigma ^{}(v)`$, (1.66), and (1.64) again one finds
$$1=\sigma (v)\frac{d^dq}{(2\pi )^d}\frac{4\mathrm{}}{(cq^2+[\sigma ](v))^2}\sigma (v)\left(\frac{4\mathrm{}c_d}{c^{d/2}}\right)[\sigma (v)]^{(d4)/2}$$
(1.68)
Since the integral is ultraviolet convergent, we have taken the short-distance momentum cutoff to infinity. $`c_d`$ is a simple number
$$c_d=\frac{d^dq}{(2\pi )^d}\left(\frac{1}{q^2+1}\right)^2=\frac{(2d)\pi ^{1d/2}}{2^{d+1}\mathrm{sin}(d\pi /2)\mathrm{\Gamma }(d/2)}$$
(1.69)
with $`c_{d=3}=1/(8\pi )`$, $`c_{d=2}=1/(4\pi )`$. Derivating one more time one gets for the effective self energy:
$$[\sigma ](v)=(u/u_0)^{2/\theta }$$
(1.70)
where $`\theta =(d2)`$ and $`v_0=8\mathrm{}c_dc^{d/2}/(4d)`$. The shape of $`[\sigma ](u)`$ is shown on Figure 5.
The solution (1.70), is a priori valid up to a breakpoint $`u_c`$, above which $`[\sigma ]`$ is constant, since $`\sigma ^{}(u)=0`$ is also a solution of the variational equations. $`u_c`$ can also be extracted from the saddle point equations and we refer the reader to for details. The precise value of $`u_c`$ is unimportant for our purpose but the existence of the two distinct regimes in $`[\sigma ](u)`$ has a simple physical interpretation that we examine in section 6.2 Using (1.70) one can now compute the correlation functions. Larger distances correspond to less massive modes and is dominated by the small $`u`$ behavior of (1.70). One obtains
$`\overline{(\varphi (r)\varphi (0))^2}`$ $`=`$ $`2\mathrm{}{\displaystyle \frac{d^dq}{(2\pi )^d}(1\mathrm{cos}(qr))\stackrel{~}{G}(q)}`$ (1.71)
$`\stackrel{~}{G}(q)`$ $`=`$ $`{\displaystyle \frac{1}{cq^2}}(1+{\displaystyle _0^1}{\displaystyle \frac{dv}{v^2}}{\displaystyle \frac{[\sigma ](v)}{cq^2+[\sigma ](v)}}){\displaystyle \frac{Z_d}{q^d}}`$ (1.72)
with $`Z_d=(4d)/(4\mathrm{}S_d)`$ and $`1/S_d=2^{d1}\pi ^{d/2}\mathrm{\Gamma }[d/2]`$. Thus for $`2<d<4`$ this leads to a logarithmic growth,
$$\overline{(\varphi (x)\varphi (0))^2}=\frac{1}{2}A_d\mathrm{log}|x|$$
(1.73)
with $`A_d=4d`$, instead of the power law growth (1.62) of the replica symmetric solution. Note that the amplitude is independent of disorder.
#### 6.2 If it ain’t broken …
It is thus necessary to break the replica symmetry to get the correct asymptotic physics ($`1/q^d`$ propagator) instead of the $`1/q^4`$ given by the replica symmetric solution. From (1.72) one easily sees that at large enough $`q`$ (i.e. for short distances) one recovers the replica symmetric solution. Thus $`u_c`$ and $`[\sigma (u_c)]`$ (see figure 5) define a lengthscale $`L(c/[\sigma (u_c)])^{1/2}`$ above which the RS solution does not describe the physics. This lengthscale corresponds to the pinning length similar to (1.24) obtained by balancing the elastic energy with the disorder one. Here
$`E_{\mathrm{el}}`$ $``$ $`L^{d2}`$ (1.74)
$`E_{\mathrm{dis}}`$ $``$ $`𝒟^{1/2}L^{d/2}`$ (1.75)
leading to the famous Larkin length $`L_{\mathrm{loc}}(1/𝒟)^{1/(4d)}`$ which corresponds to the distance for which relative displacements are of the order of the lattice spacing (or for which the phase here is of the order $`2\pi `$) . Below this lengthscale the system has a single equilibrium state. It can be described by simply expanding in the displacements in (1.13)
$$H_{\mathrm{dis}}=d^dxf(x)\varphi (x)$$
(1.76)
where $`f`$ is a random force, derivative of the random potential $`V`$. It is easy to see that this model gives the $`1/q^4`$ propagator. This breaks down when the displacements cannot be expanded, i.e. for distances larger than $`L_{\mathrm{loc}}`$. As anticipated by Larkin and Ovchinikov, pinning occurs that tends to keep a domain of size $`L_{\mathrm{loc}}`$ in place despite the thermal fluctuations.
This help us to understand the physics of the RSB solution. To illustrate it let us focus for simplicity to the case in $`d=2`$. As can be seen from (1.70) by letting $`d2`$ the RSB solution in that case is a one-step breaking as shown in Figure 5. There are some additional complications but they are unimportant for our discussion here. A normal self consistent approximation would approximate the energy by a simple Gaussian centered in $`\varphi =0`$. The only way to incorporate the pinning is thus to put a mass term in the propagator
$$q^2q^2+L_{\mathrm{loc}}^2$$
(1.77)
As can be immediately inferred from the discussion of section 3.1, this is a much too crude approximation. The RSB solution is smarter and approximates effectively the distribution of displacements by a hierarchical superposition of Gaussians centered at different randomly located points in space. A pictorial view of these two cases is shown on Figure 6.
The double distribution over environment and thermal (or quantum) fluctuations is approximated as follows . In each environment there are effective “pinning centers” corresponding to the low lying metastable states (preferred configurations). Since all $`q`$ modes are in effect decoupled within this approximation, for each $`q`$ mode a preferred configuration (a state) is $`\varphi _\alpha (q)`$. They are distributed according to:
$`P(\{\varphi _\alpha (q)\}){\displaystyle \underset{\alpha }{}}e^{\frac{c}{2T_g}q^2|\varphi _\alpha (q)|^2}`$ (1.78)
Each is endowed with a free energy $`f_\alpha `$ distributed according to an exponential distribution $`P(f)\mathrm{exp}(u_cf/T)`$ (here $`u_c=T/T_g`$). Once these seed states are constructed, the full thermal distribution of the $`q`$ mode $`\varphi _q`$ is obtained by letting it fluctuate thermally around one of the states:
$$P(\varphi _q)\underset{\alpha }{}W_\alpha e^{\frac{c}{2T}(q^2+L_{\mathrm{loc}}^2)|\varphi _q\varphi _\alpha (q)|^2}$$
(1.79)
where each state is weighted with probability $`W_\alpha =e^{f_\alpha /T}/_\beta e^{f_\beta /T}`$. One thus recovers qualitatively the picture of Larkin Ovchinikov as the solution of the problem with the replica variational method. The Larkin length naturally appears as setting the (internal) size of the elastically correlated domains. The full RSB case corresponds to more levels in this hierarchy of Larkin domains (in some sense there are clusters of domains of size larger than $`L_{\mathrm{loc}}`$) and the way this hierarchy scales with distance reproduces the exponents for displacements and energy fluctuations. We refer the reader to for more details on the use of the variational method in this context and for the physical properties of such systems.
#### 6.3 Quantum problems
We now apply the same method to the quantum problem. Because of the non locality in time of the interaction in (1.47) the solution will have quite different properties. We again use the trial action (1.48), with the parametrization
$$vG_{ab}^1(q,\omega _n)=\frac{((vq)^2+\omega _n^2)}{\pi K}\delta _{ab}\sigma _{ab}(\omega _n)$$
(1.80)
Although, as in section 6.1 $`\sigma _{ab}`$ is still independent of $`q`$ because of the locality in space it is now dependent on $`\omega _n`$. This would render the solution extremely complicated if it were not for a remarquable property of quantum systems. Off diagonal replica terms such as $`\sigma _{ab}`$ only exist for the mode $`\omega _n=0`$ The general argument is that in each realization of the random potential $`V`$, the disorder does not depend on $`\tau `$. Therefore before averaging over disorder:
$$G_{ab,V}=\varphi _a(x,\tau )\varphi _b(0,0)=\varphi _a(x,\tau )\varphi _b(0,0)=\varphi _a(x,0)\varphi _b(0,0)$$
(1.81)
It is important to note that such a property crucially depends on the assumption that the hamiltonian is $`\tau `$-independent and on the fact that equilibrium has being attained. This is the case considered here.
The static mode $`\omega _n=0`$ thus plays a special role. This is quite natural in a time independent disorder. Quite naively one sees already that the properties of the variational solution will thus be very similar to the ones of point like (i.e. totally uncorrelated) disorder in $`d`$ spatial dimensions (here $`d=1`$). This suggest strongly that the variational method will pull out a static solution, reminiscent of the one introduced in section 3.1 and treat the fluctuations around this static solution. We will come back to this point later. The solution can be obtained quite in all dimensions and we refer the reader to for details. We specialize here to the case $`d=1`$. In this case two type of solution exist, a simple RS solution with $`\sigma _{ab}=0`$. This solution is stable for $`K>3/2`$. It corresponds of course to the delocalized regime where the cosine term in (1.47) is irrelevant. The variational method correctly reproduces the (gaussian) delocalized regime, and the correct transition point, but of course misses the renormalization of the Luttinger parameters given by the RG. For $`K<3/2`$ although an RS solution still exists it is unstable and physically obviously incorrect . One should look for an RSB solution. In that case the correct solution is a one-step RSB solution (it would be full RSB for $`d>2`$) of the type shown in Figure 5.
$`G_c^1(q,\omega _n)`$ $`=`$ $`{\displaystyle \frac{\mathrm{}}{\pi K}}(vq^2+{\displaystyle \frac{\omega _n^2}{v}})+{\displaystyle \frac{2𝒟}{\mathrm{}(\pi \alpha )^2}}{\displaystyle _0^\beta \mathrm{}}𝑑\tau (1\mathrm{cos}(\omega _n\tau ))`$ (1.82)
$`\left[\mathrm{exp}(2\mathrm{}\stackrel{~}{B}(x=0,\tau )){\displaystyle _0^1}𝑑u\mathrm{exp}(2\mathrm{}B(u))\right]`$
with
$`\sigma (q,\omega _n,u)={\displaystyle \frac{2𝒟v}{(\pi \alpha )^2}}\beta \mathrm{exp}(\mathrm{}2B(u))\delta _{\omega _n,0}`$ (1.83)
These equations are still formidable to solve. A simple parametrization is
$`vG_c^1(q,\omega _n)={\displaystyle \frac{1}{\pi \overline{K}}}((vq)^2+\omega _n^2)+\mathrm{\Sigma }_1(1\delta _{n,0})+I(\omega _n)`$ (1.84)
$`I(\omega _n)={\displaystyle \frac{2𝒟v}{(\pi \alpha )^2\mathrm{}}}{\displaystyle _0^\beta \mathrm{}}\left[e^{2\mathrm{}\stackrel{~}{B}(\tau )}e^{2\mathrm{}B(u>u_c)}\right](1\mathrm{cos}(\omega _n\tau ))𝑑\tau `$ (1.85)
$`\mathrm{\Sigma }_1=u_c(\sigma (u>u_c)\sigma (u<u_c))=[\sigma ](u>u_c)`$ (1.86)
$`\sigma (u)={\displaystyle \frac{2𝒟v}{(\pi \alpha )^2}}e^{\mathrm{}2B(u)}\beta \delta _{n,0}`$ (1.87)
The parameters $`\mathrm{\Sigma }_1`$, the breakpoint $`u_c`$ and the function $`I(\omega _n)`$ have to be determined self-consistently. Let us examine first the general properties of the solution.
Since $`I(\omega _n=0)=0`$ it is easy to check from (1.16) that the compressibility is unchanged by the disorder, since the “mass” term $`\mathrm{\Sigma }_1`$ goes also away at $`\omega _n=0`$. The variational method thus correctly reproduces that the compressibility of an Anderson insulator is still finite, and practically unchanged (for free electrons where we can compute it) from the value without disorder. Correlation functions are also easy to obtain. We just give here the important physical point without the explicit derivation . Because of the presence of $`\mathrm{\Sigma }_1`$ in the propagator (1.84) they will be massive. This leads to
$$B(x\mathrm{},\tau =0)x/L_{\mathrm{loc}}$$
(1.88)
A full calculation of the correlations shows that $`\stackrel{~}{B}(\tau )`$ grows until $`\tau L_{\mathrm{loc}}`$ when it saturates
$$B(x=0,\tau \mathrm{})\mathrm{Cste}=(\rho _0^1l_{})^2$$
(1.89)
Since the variational action is gaussian one has for the correlation of the $`2k_F`$ part of the density
$$\chi _\rho (x,\tau )=e^{2B(x,\tau )}$$
(1.90)
Two different physical effects are described by the above correlation functions. In the absence of disorder $`\chi _\rho (\tau \mathrm{})`$ would go to zero as a power law, a sign of the wandering of the particles. (1.89) traduces the fact that the time independent potential localizes the particles at a given point in space, instead of letting them fluctuate (unboundedly) due to quantum fluctuations in the absence of disorder. This leads to density correlation in time going to a constant up to a Debye Waller like factor. The behavior (1.89) thus shows that a static solution $`\varphi _0(x)`$ exists. The variational approach thus provides a good justification for the static solution $`\varphi _0`$ on which the method of section 3 is built. The correlation length of this static solution is given by (1.88). Since this length controls through (1.90) the exponential decay of the spatial correlations of density it is thus related to the standard localization length. Note that here both $`L_{\mathrm{loc}}`$ and the full static solution are determined by the variational method. Simple dimensional analysis on (1.84) shows that
$$L_{\mathrm{loc}}1/\sqrt{\mathrm{\Sigma }_1}.$$
(1.91)
The solution of the variational equations leads back to the expression (1.28) for $`L_{\mathrm{loc}}`$. Quite interestingly (1.89) defines a length if one writes the constant in units of the fermion spacing. This length is the “width” of the fluctuations of the particles as shown in Figure 7.
Note the difference between the two lengths $`l_{}`$ and $`L_{\mathrm{loc}}`$. They are related through
$$l_{}^2=\frac{\alpha ^2}{\pi ^2}(\mathrm{}\overline{K})\mathrm{ln}(L_{\mathrm{loc}}/a)$$
(1.92)
At the transition one expects that $`L_{\mathrm{loc}}`$ diverges as $`L_{\mathrm{loc}}\mathrm{exp}(b/(K_cK)^\alpha )`$ with $`\alpha =1/2`$ from the RG (see equation (1.35)). Thus relation (1.92) predicts that:
$$l_{}\frac{1}{(K_cK)^{\alpha /2}}$$
(1.93)
diverges as a power law, which could be measured in numerical simulations.
Let us now look at the conductivity, for which we need $`I(\omega _n)`$. If we call $`I^{}(\omega )`$ and $`I^{\prime \prime }(\omega )`$ the real and imaginary parts of the analytic continuation of $`I(\omega _n)`$, then the conductivity is given by
$$\sigma (\omega )=\frac{\omega I^{\prime \prime }+\omega (\omega ^2+I^{}+\mathrm{\Sigma }_1)}{(\omega ^2+\mathrm{\Sigma }_1+I^{})^2+(I^{\prime \prime })^2}$$
(1.94)
It is easy to see that $`\mathrm{}\sigma (\omega )`$ goes to zero at zero frequency and thus the phase is indeed localized. But because of the analytic continuation, the existence or not of a gap in the optical conductivity is not linked to the existence of a “mass” $`\mathrm{\Sigma }_1`$ but to whether $`I^{\prime \prime }`$ is nonzero at small frequencies or not. The equation for $`I(\omega _n)`$ takes a particularly simple form in the limit $`K,\mathrm{}0`$ while keeping $`\overline{K}`$ fixed. In this limit, we can write
$$I(\omega _n)=\mathrm{\Sigma }_1f(\frac{\omega _n}{\sqrt{\pi \overline{K}\mathrm{\Sigma }_1}}),$$
(1.95)
where the scaling function $`f`$ satisfies:
$$f(x)=2\left[1\frac{1}{\sqrt{1+x^2+f(x)}}\right].$$
(1.96)
It can be shown that for $`\omega 0`$
$$\mathrm{}\sigma (\omega )\omega ^2$$
(1.97)
In a similar way (1.94) gives
$$\mathrm{}\sigma (\omega )\omega /\mathrm{\Sigma }_1$$
(1.98)
Such behavior is in agreement with exact results in one dimension up to logarithmic prefactors. At high frequency, one can show that $`\sigma (\omega )\omega ^{2K4}`$ in agreement with the RG result. The resulting conductivity is plotted on Figure 8.
#### 6.4 The fine prints
The variational method is thus an extremely efficient method for this type of disordered problems. It allows to get most of the physical properties in the localized phase, and give the qualitative features of the transition and the delocalized phase. Its physics is very similar to the one described for point like disorder in section 6.1. The variational method determines the “static” solution for the mode $`\omega _n=0`$ (but by taking into account the effects of all modes) without being restricted to a single Gaussian. It can then correctly compute the fluctuations (both in $`q`$ and $`\omega _n`$) around this solution, contrarily to the approximate method of section 3.1 which was impeded by the lack of knowledge of the static solution.
In $`d=1`$ an additional remarkable property can be seen. Although most of the properties are independent of the value of $`u_c`$, the conductivity is strongly dependent on it. As for the static solution one has continuous replica symmetry breaking for $`d>2`$ which goes continuously to the one step solution in $`d=2`$. In these dimensions taking the value of $`u_c`$ out of the variational equations gives the correct conductivity (no gap in $`\sigma (\omega )`$). In $`d=1`$, the $`u_c`$ obtained by minimizing the free energy would give an incorrect (gapped) conductivity. Another way to determine $`u_c`$ is to use the marginality condition that corresponds to the instability of the replicon mode, similar to the condition (1.63). This condition, more dynamical in nature, coincides with the free energy value of $`u_c`$ for $`d2`$ but is different in $`d=1`$. One can check that the marginality condition always gives the correct conductivity. Some arguments for why it is so were given in , but this point is not yet fully understood. This phenomenon has since been found to occur in other systems .
The success of the variational method is obviously linked to the fact that we are looking at small “quadratic” fluctuations around a certain (in our case highly disordered) solution. In this case the replacement
$$\mathrm{cos}(\varphi )\frac{1}{2}\varphi ^2$$
(1.99)
is very reasonable. Such an approximation is very good to compute correlation functions of the variable $`\varphi `$. This is the case for the density-density correlation and the optical conductivity. What is missing in the approximation (1.99) are the solitons that go from one minimum of the potential to another minimum. The energy cost of such excitations is grossly overestimated by the variational method as is shown on Figure 9.
This has important consequences for the calculation of correlation functions that involve the soliton creation operator $`e^{i\theta }`$. These correlation functions are found incorrectly to decay exponentially with distance at equal time and to be zero at unequal time. As a result, the variational method does not allow the calculation of Fermion Green’s function (they involve the operators $`e^{i(\theta \pm \varphi )}`$) nor superconducting correlations (which involve the operator $`e^{2i\theta }`$). In the case of a disordered XXX spin chain, the situation is even worse, since $`S_x\mathrm{cos}(\theta )`$, whereas $`S_z\frac{_x\varphi }{\pi }+()^{x/a}\mathrm{cos}2\varphi `$. Therefore, even in the presence of a disorder that preserves SU(2) symmetry (such as a random bond disorder), the variational method would lead to a spurious breaking of rotational symmetry. It would also give poor results for systems that include both the $`\theta `$ and the $`\varphi `$ field in the Hamiltonian. Examples of such theories include disordered Hubbard ladders , disordered spin ladders or XXX spin chain in a random fields .
The mishandling of soliton excitations also limits our knowledge of the transport properties at finite temperatures. Indeed the optical conductivity does not correspond to transport of charge but charge oscillations around the equilibrium positions, it is thus well described by our harmonic approximation as shown in Figure 9. On the contrary transport at finite $`T`$ involves real charge displacements. Since $`\rho (x)\varphi `$, displacing a charge amounts to make a solitonic excitation in the field $`\varphi `$ as shown in Figure 10. Indeed Mott’s arguments, to compute the conductivity of noninteracting electrons in presence of phonons $`\sigma (T)\mathrm{exp}\left[\left(\frac{T_0}{T}\right)^{\frac{1}{d+1}}\right]`$ is strongly reminiscent of an instanton calculation as shown in Figure 10.
A similar argument applies to Efros and Schklovskii calculations. An extension of the variational approach to finite temperatures would indeed still lead to $`\sigma (T,\omega =0)=0`$ proof that it is missing the excitations that are important at finite temperature. Unfortunately no way to treat such solitons has been found at present despite some attempts .
Despite these limitations, the variational method is up to now the only analytical method giving information for such localized systems in the localized phase. As with all variational approaches, some physical insight in the properties of the system under consideration is needed to determine whether the method as any chance of success. Clearly, one must apply this method only to systems that can be reasonably well understood qualitatively from their classical action. Fortunately many systems fall in this category, and we examine some of those in the following.
#### 6.5 Higher dimension: electronic crystals and classical systems
First the GVM can be used to study classical systems using the standard mapping $`\tau z`$. The action (1.47) and its extension to higher spatial dimension describes elastic objects (lines in this case as shown in Figure 11) pinned by columnar (i.e. time or $`z`$ independent) defects. This situation is realized for example in vortex in type II superconductors irradiated by heavy ions (creating the linear track of disorder). This system in $`2+1`$ dimension is equivalent to a $`d=2`$ quantum bose system in presence of pins. In a similar way than in $`d=1`$ (see section 6.3) such system has a pinned phase (the Bose glass) . The variational method can be used to describe the Bose glass phase . However contrarily to $`d=1`$ it cannot be used to go to the superfluid regime since to describe a two dimensional “melting” of the Bose glass phase dislocations are important (no dislocations exist in $`d=1`$) and for reasons explained above the GVM overestimates the energy cost of topological excitations. Another way to say it, is that in $`d>1`$ we loose the elastic description (1.39) of the Fermion or Boson operators. The GVM can thus only be used in phases where the particles are localized so that some elastic description can again be used.
We can thus use the variational method in higher dimensions to study electronic crystals. This included Charge density wave, but also the two dimensional Wigner crystal of electrons. In such a phase the electrons are confined by their repulsion (and a in some systems a magnetic field). An elastic description can be used. Some level of quantumness is hidden in the elastic parameters (“size” of the particles, quantization of the phonon modes of the crystal). For such systems the calculation of the optical conductivity is particularly useful since it is one of the few probes of such systems. Since the physics of such systems would deserve a review of its own we will not dwell further on it here but refer the reader to for details.
### 7 Commensurate systems
When the filling of the fermion system is commensurate the physics discussed above is modified in various ways since the backward scattering on disorder becomes real. If the forward scattering still exists, not much is changed. Two special cases will thus occur: (i) the forward scattering is absent. This occurs because of a symmetry of the system. This is the case for example for fermions at half filling with a random hopping or for spin chains with random exchange. (ii) The commensurate potential (either due to the lattice or due to electron-electron interaction) would open a gap. There will be a competition between Mott physics wanting to get a commensurate (gapped) insulator and the disorder that would like to destroy such a gap (push the system locally away from commensurability).
#### 7.1 The peculiar random exchange
For electrons at half filling with a random exchange the forward scattering does not exist and the disorder term is simply
$`H`$ $`=`$ $`{\displaystyle }dxV(x)i[\psi _+^{}\psi _{}\mathrm{H}.\mathrm{c}.]`$ (1.100)
$`=`$ $`{\displaystyle 𝑑x\frac{V(x)}{(2\pi \alpha )}\mathrm{sin}(2\varphi (x))}`$ (1.101)
Although this seems very similar to (1.13) one easily sees the difference on Figure 12. Contrarily to normal disorder where $`\varphi `$ follows the random phase of the random potential, here $`\varphi =\pm \pi /4`$ depending on the sign of the potential. Thus $`\varphi `$ is nearly gapped but for its kinks. The low energy properties will thus be dominated by the kinks in $`\varphi `$. Such kink structure between the doubly degenerate minima makes it unlikely that the GVM can be used for this problem.
In order to get an idea of the physics let use examine the case of free fermions $`K=1`$. For $`E=0`$ one can easily construct the eigenstates by solving
$`\psi _+`$ $``$ $`V(x)\psi _{}(x)=0`$ (1.102)
$`\psi _{}`$ $`+`$ $`iV(x)\psi _+(x)=0`$ (1.103)
to obtain the (unnormalized)
$$\psi _+(x)\psi _{}(x)e^{_0^x𝑑yV(y)}$$
(1.104)
which obviously decays as
$$\psi _\pm (x)e^\sqrt{x}$$
(1.105)
If one wants to view such a state as an exponentially localized state this means that the localization length diverges when $`E0`$. The divergence of the localization length is of course the signature that the physical properties will be quite different than the ones for the “normal” disorder studied in the previous sections. For the non-interacting case they have been studied using a variety of techniques such as replicas, Berezinskii methods, supersymmetry etc. . The density of states diverge at the Fermi level as
$$\rho (ϵ)\frac{1}{ϵ\mathrm{log}(1/ϵ)^3}$$
(1.106)
Quite astonishingly if one considers that all states are still localized the d.c. conductivity is now finite . This is surprising since it seems to violate the Mott phenomenological derivation. The conductivity is proportional to the absorbed power. To make a transition between two localized states one need one occupied state, one empty one separated by an energy $`\mathrm{}\omega `$. The number of such states is $`\rho (ϵ=0)\mathrm{}\omega `$. The absorbed power (and the conductivity) is thus (up to log correction)
$$\sigma (\omega )\rho (ϵ=0)(\mathrm{}\omega )^2$$
(1.107)
correctly giving back for a constant density of states at the Fermi level the $`\omega ^2`$ dependence of section 6.3. One would thus naively have expected a $`\sigma (\omega )\omega `$ in the case of a singular density of state (1.106). A qualitative understanding of the behavior of the d.c. conductivity is thus still lacking.
When interactions are included, the problem becomes more difficult to tackle. Going back to the spin chain version of (1.100) a real space renormalization procedure has been introduced. This procedure works beautifully and allows the calculation of most correlation functions. We refer the reader to for details. As for the normal disorder (see section 3.2) interaction were found to be irrelevant at the disordered fixed point. The divergence of the localization length also changes drastically the correlation functions. In particular the spin spin correlation functions now decay as power law instead of being exponential. It would be extremely interesting to have an equivalent derivation of this real space RG directly in the boson representation.
#### 7.2 Mott versus Anderson
A particularly interesting situation occurs when the non-disordered system possesses a gap. In that case the competition between this gap and the disorder is non trivial. This arises in a large number of systems such as disordered Mott insulators , systems with external (Peierls or spin-Peierls systems ) or internal commensurate potential (ladders or spin ladders , disordered spin 1 chains ). It is physically simple to see that in order to destroy the gap one needs a disorder comparable to the gap (although in some specific cases Imry-Ma effects can destroy the gap for infinitesimal disorder ). This makes the complete description of the gap closure and of the physics of the resulting phases is extremely difficult with the usual analytic techniques such as the perturbative renormalization group described in section 3.2, due to the absence of a weak coupling fixed point. Indeed both the commensurate Mott insulator (which can be described by a sine-Gordon Hamiltonian ) and the disordered Anderson insulator correspond to strong coupling fixed points.
To be concrete, but keep the technical details to a minimum let us consider again spinless fermions. The commensurability can be described by adding to the Hamiltonian (1.9) a periodic potential at the Fermi level
$$H_{\mathrm{com}}=g𝑑x(\psi _+^{}\psi _{}+\psi _{}^{}\psi _+)$$
(1.108)
which becomes in the boson representation
$$H_{\mathrm{com}}=𝑑x\frac{g}{\pi \alpha \mathrm{}}\mathrm{cos}2\varphi $$
(1.109)
The reason to consider here a periodic potential is specific to pathologies associated to the Mott gap for spinless fermions. Similar results are expected for a Mott insulator (which is due to the $`4k_F`$ potential of the lattice).
If this system is studied by RG then one is faced with the competition of two strong coupling fixed points since
$`{\displaystyle \frac{dg}{dl}}`$ $`=`$ $`(2K)g`$ (1.110)
$`{\displaystyle \frac{d𝒟}{dl}}`$ $`=`$ $`(32K)𝒟`$ (1.111)
Using the usual qualitative argument consisting in taking the most divergent operator, it was concluded that if disorder reaches strong coupling ($`𝒟(l^{})=1>g(l^{})`$) first, we will be in the Anderson phase, with a localization length $`l_0=ae^l^{}\left(\frac{1}{𝒟}\right)^{\frac{1}{32K}}`$, whereas if the commensurate potential reaches strong coupling first, we will be in the Mott phase with a correlation length (or soliton size) $`d\left(\frac{1}{g}\right)^{\frac{1}{2K}}`$. The phase transition between the two phases occurs for $`l_0d`$. This picture relies on the important assumption that there is no other stable fixed point than the Mott insulator and the Anderson insulator. Even if it was so, it would not be possible to determine the phase boundaries, nor determine what type of critical point separate the Mott and the Anderson insulator. In order to make progress on these issues and obtain a more complete picture of dirty one-dimensional Mott insulators, one needs to solve the problem non-perturbatively . The methods of section (6.3) are well adapted since the problem can be cast in the sine-Gordon form.
#### 7.3 Variational approach
For the Mott versus Anderson problem, the variational action reads:
$`{\displaystyle \frac{S_{\mathrm{rep}.}}{\mathrm{}}}`$ $`=`$ $`{\displaystyle \underset{a}{}}\left[{\displaystyle \frac{dxd\tau }{2\pi K}\left(v(_x\varphi _a)^2+\frac{(_\tau \varphi _a)^2}{v}\right)}{\displaystyle \frac{g}{\pi \alpha \mathrm{}}}{\displaystyle 𝑑x𝑑\tau \mathrm{cos}2\varphi _a}\right]`$ (1.112)
$``$ $`{\displaystyle \frac{𝒟}{(2\pi \alpha \mathrm{})^2}}{\displaystyle \underset{a,b}{}}{\displaystyle 𝑑x_0^\beta 𝑑\tau 𝑑\tau ^{}\mathrm{cos}\left(2(\varphi _a(x,\tau )\varphi _b(x,\tau ^{}))\right)}`$
The term $`\mathrm{cos}2\varphi _a`$ in (1.112) is responsible for the opening of a gap. We search for a saddle point with a form of the variational connected Green’s Function slightly generalized with respect to (1.84):
$$vG_c^1(q,\omega _n)=\frac{1}{\pi \overline{K}}(\omega _n^2+v^2q^2)+m^2+\mathrm{\Sigma }_1(1\delta _{n,0})+I(\omega _n)$$
(1.113)
Where the parameter $`\mathrm{\Sigma }_1`$ and the function $`I(\omega _n)`$ satisfy the equations (1.86) and (1.85). We give here the main steps of the solution and refer the reader to for more details. The variational self-energy satisfies the equation (1.87). Finally, $`m`$ satisfies the equation:
$$m^2=\frac{4gv}{\pi \alpha }e^{2\mathrm{}\stackrel{~}{G}(0,0)}$$
(1.114)
The important physical quantities are simply given such as conductivity and compressibility are simply given by (1.15) and (1.55).
Since we expect the physics to be continuous for small enough $`K`$ (i.e. repulsive enough interactions), one can gain considerable insight by considering (see sec. 3.1) the classical limit $`\mathrm{}0`$, $`K0`$ keeping $`\overline{K}=K/\mathrm{}`$ fixed. In this limit one can solve analytically the saddle point equations (1.84)–(1.87), (1.113)–(1.114) and compute $`m`$, $`\mathrm{\Sigma }_1`$ and $`I(\omega _n)`$. The resulting phase diagram is parameterized with two physical lengths (for $`K0`$): The correlation length (or soliton size) of the pure gapped phase
$$d=\left(\frac{4g\overline{K}}{\alpha v}\right)^{1/2}$$
(1.115)
and the localization (or pinning) in the absence of commensurability length
$$l_0=\left(\frac{(\alpha v)^2}{16𝒟\overline{K}^2}\right)^{1/3}$$
(1.116)
Contrarily to the naive direct transition predicted by the extrapolation of the RG, we find three phases as shown in Figure 13.
Their main characteristics in term of conductivity and compressibility are summarized in Figure 14.
Mott insulator
At weak disorder we find a replica symmetric solution with $`\mathrm{\Sigma }_1=0`$ but with $`m0`$ ($`m`$ depends on the disorder). $`m`$ is given by the simple equation:
$$m^2=\frac{4gv}{\pi \alpha }\mathrm{exp}\left[\frac{𝒟\overline{K}^{1/2}}{\alpha ^2\pi ^{3/2}v^{1/2}m^3}\right].$$
(1.117)
It is convenient to work with the physical lengths $`d`$ and $`l_0`$ and the correlation length $`\xi `$ where:
$$\xi ^2=\frac{v^2}{(\pi \overline{K}m^2)}.$$
(1.118)
One can then rewrite (1.117) as:
$$\frac{1}{\xi ^2}=\frac{1}{d^2}\mathrm{exp}\left[\frac{1}{16}\left(\frac{\xi }{l_0}\right)^3\right].$$
(1.119)
For $`l_0/d>\frac{1}{2}\left(\frac{3e}{4}\right)^{1/3}`$, (1.119) has a single physical solution. For $`l_0/d<\frac{1}{2}\left(\frac{3e}{4}\right)^{1/3}`$, (1.119) has no solution, which means that the Mott insulator becomes unstable. In fact due to another contraint (see below) the Mott Insulator becomes unstable at an even smaller disorder.
Having $`m0`$ leads to zero compressibility $`\kappa =0`$. Disorder reduces the gap created by the commensurate potential and thus increases $`\xi `$ compared to the pure case. Let us now examine the equation for $`I(\omega _n)`$ giving the transport properties. An expansion around $`\mathrm{}=0`$ in equation (1.85) gives the self-consistent equation for $`I(\omega _n)`$:
$$I(\omega _n)=\frac{8𝒟v}{(\pi \alpha )^2}_0^\beta \mathrm{}G_c(x=0,\tau )(1\mathrm{cos}(\omega _n\tau ))𝑑\tau $$
(1.120)
Introducing the scaling form (to be contrasted with (1.95)):
$$I(\omega _n)=m^2f\left(\frac{\omega _n}{\sqrt{\pi \overline{K}}m}\right),$$
(1.121)
(1.120) can be recast in the form:
$$f(x)=\lambda \left[1\frac{1}{\sqrt{1+x^2+f(x)}}\right]$$
(1.122)
where:
$$\lambda =\frac{4𝒟\overline{K}^{1/2}v}{\pi ^{3/2}\alpha ^2m^3}=\frac{1}{4}\left(\frac{\xi }{l_0}\right)^3$$
(1.123)
Let us note that for $`\lambda =2`$, the equation (1.122) reduces to the equation (1.96). For $`\lambda >2`$ (1.122) has no physical solution. Using (1.123), this condition becomes:
$$\frac{l_0}{d}<\frac{1}{2}e^{\frac{1}{4}}$$
(1.124)
For $`\lambda <2`$, there is a physical solution of (1.122) such that $`lim_{x\pm \mathrm{}}f(x)=1+\lambda `$ and for $`x1`$ , $`f(x)=1+\alpha x^2+o(x^2)`$ with $`\alpha =\lambda /(2\lambda )`$. The conductivity of the Mott insulator can be obtained from $`f`$ in the form:
$$\sigma (\omega )=\frac{\xi \overline{K}}{\pi }\frac{ıx}{(1+f(ıx)x^2)}$$
(1.125)
where $`x=\omega /\omega ^{}`$ and $`\omega ^{}=v/\xi `$ is the characteristic frequency associated with the correlation length $`\xi `$. The conductivity $`\sigma (\omega )`$ is zero if:
$$\omega <\omega _c=\omega ^{}\sqrt{1+\lambda 3\left(\frac{\lambda }{2}\right)^{2/3}}$$
(1.126)
The replica symmetric phase with $`m0`$ has thus a conductivity gap $`\mathrm{}\omega _c`$ and can be assimilated to a Mott insulator (MI). The behavior of the conductivity is plotted on figure (15). Close to the threshold, one has $`\sigma (\omega )(\omega \omega _c)^{1/2}`$.
However the gap in the conductivity $`\mathrm{}\omega _c`$ decreases when disorder increases, and closes for $`d/l_0=2e^{1/4}1.557`$ (see Figure 16). For $`d/l_02e^{1/4}`$, the conductivity gap vanishes linearly with $`\frac{d}{l_0}`$.
For $`d/l_0>2e^{1/4}`$ the MI solution becomes unphysical *even though* the mass $`m`$ remains finite at this transition point, i.e. the system remains incompressible. For stronger disorder one must break replica symmetry, as for the pure disorder case of section 6.3. Here, however *two* possibilities arise depending on whether the saddle point allows for $`m0`$ or not. In the presence of a breaking of replica symmetry, one extra equation is needed to determine the breakpoint. As in the case of the Anderson insulator, such equation is provided by the marginality of the replicon condition discussed in section 6.4. Two phases exists:
The Anderson Glass
For large disorder compared to the commensurate potential $`d/l_0>1.861`$, $`m=0`$ is the only saddle point solution. The saddle point equations then reduce to those (1.84)–(1.87) . Thus, we recover the Anderson glass with interactions of section (6.3). As we have seen in section (6.3), in such phase the conductivity starts as $`\sigma (\omega )\omega ^2`$ showing no gap and the compressibility is finite.
In the Anderson glass phase, the disorder washes out completely the commensurate potential. The MI and the AG were the only two phases accessible by renormalization techniques . Within the replica variational formalism however, we find that an intermediate phase exists between them.
The Mott Glass
For intermediate disorder $`2e^{1/4}<d/l_0<1.861`$ a phase with *both* $`\mathrm{\Sigma }_10`$ and $`m0`$ is obtained. We shall call this phase the Mott Glass (MG). We shall not discuss in full detail the one-step solution of the saddle point equation here. we will rather stress the salient features of our solution. First, as a result of the marginality of replicon mode condition, $`m^2+\mathrm{\Sigma }_1`$ remains constant in the MI and MG as disorder strength is increased . In the MG phase, $`I(\omega _n)`$ is still of the form (1.121) but $`m`$ is replaced by $`\sqrt{m^2+\mathrm{\Sigma }_1}`$. The reduced self-energy $`f(x)`$ satifies (1.122) but with $`\lambda =2`$ in the whole Mott Glass phase. This implies that (see (1.96)) that the a.c. conductivity of the Mott Glass is identical to the one of an Anderson glass. However, since $`m0`$ in the MG, the system is incompressible ($`\kappa =0`$) like a Mott Insulator. Thus, the Mott Glass is a new glassy phase (since it has Replica Symmetry Breaking) with characteristics intermediate between those of an Anderson Insulator and those of a Mott Insulator.
#### 7.4 Physical discussion
The existence of a phase with a compressibility gap but no conductivity gap is quite remarkable since by analogy with noninteracting electrons one is tempted to associate a zero compressibility to the absence of available states at the Fermi level and hence to a gap in the conductivity as well. Our solution shows this is not the case, when interactions are turned on the excitations that consists in adding one particle (the important ones for the compressibility) become quite different from the particle hole excitations that dominate the conductivity. A similar situation is obtained in the case of the one dimensional Wigner crystal , which has the conductivity of a perfect 1d metal, $`\sigma (\omega )\delta (\omega )`$ but a zero compressiblity since $`\chi =lim_{q0}\frac{1}{\mathrm{ln}q}`$. This argument suggest that the difference in one-particle and two-particle properties is a consequence of the strong repulsion in the system.
In addition to the variational method itself the Mott glass phase can also be obtained by two other independent methods. Higher dimensional extensions of the present problem, similar to the one made in section 6.1 can be treated around four spatial dimensions using a $`d=4ϵ`$ functional renormalization group method (totally different form the $`d=2`$ RG). Such study confirms the existence of the intermediate Mott glass phase. One can also analyse (1.112) for zero kinetic energy and obtain the MG phase . Although we have done the derivation of the Mott phase for fermions in one dimension we expect its physics to survive into higher dimension. This can be seen by looking at the atomic limit (zero hoping) of an interacting fermionic system (in any dimension). If the repulsion extends over at least one interparticle distance, leading to small values of $`K`$, particle hole excitations are lowered in energy by excitonic effects. For example for fermions with spins with both an onsite $`U`$ and a nearest neighbor $`V`$ the gap to add one particle is $`\mathrm{\Delta }=U/2`$. On the other hand the minimal particle-hole excitations would be to have the particle and hole on neighboring sites (excitons) and cost $`\mathrm{\Delta }_{\mathrm{p}.\mathrm{h}.}=UV`$, as shown in Figure 17.
When disorder is added the gaps decrease respectively as $`\mathrm{\Delta }\mathrm{\Delta }W`$ and $`\mathrm{\Delta }_{\mathrm{p}.\mathrm{h}.}\mathrm{\Delta }_{\mathrm{p}.\mathrm{h}.}2W`$. Thus the conductivity gap closes, the compressibility remaining zero (for bounded disorder). According to this physical picture of the MG, the low frequency behavior of conductivity is dominated by excitons (involving neighboring sites). This is at variance from the AG where the particle and the hole are created on distant sites. This may have consequences on the precise low frequency form of the conductivity such as logarithmic corrections. When hopping is restored, we expect the excitons to dissociate and the MG to disappear above a critical value $`K>K^{}`$. Since finite range is needed for the interactions, in all cases (fermions or bosons) $`K^{}<1`$. In addition we expect $`K^{}<1/2`$ for fermions with spins. One interesting question is the question of d.c. transport in the three phases, and whether the Mott glass has a d.c. transport closer to the Anderson or the Mott phase. Since the excitons are neutral, one simple guess would be that such excitations would not contribute to the d.c. transport. The d.c. conductivity in the Mott glass phase would thus be still exponentially activated just as in the Mott insulator. Of course more detailed studies would be needed to confirm this point.
### 8 Conclusions
Many disordered fermionic system can thus be successfully described by an elastic disordered theory. In one dimension, this situation is ubiquitous due to the importance of collective excitations. Most physical system, whether one starts with fermions, bosons or spins, can be represented in terms of bosonic excitations. In higher dimension such a description is valid in crystalline phases such as a Wigner crystal in which the quantum particles are strongly localized due to their interactions or charge density waves. Disorder then leads to rich physical phenomena coming from the competition between the elasticity, wanting a well ordered structure and the disorder that distorts the structure to gain pinning energy. This leads to the existence of many metastable configurations and to glassy properties. Using the various methods described in these notes, we now have a good description of the low energy excitations of such structures. This gives access to a host of physical properties such as the a.c. transport properties.
Clearly one of the most important open questions is the the issue of topological defects in such structures. Indeed, such defects are needed to describe the melting of these cristalline phases and will be necessary to go to more “liquid” phases in which the statistics (fermionic or bosonic) of the particles will play a much more crucial role. In addition d.c. transport is obviously dominated by such excitations. Unfortunately so far the methods able to tackle the properties in the localized phase such as the Gaussian Variational Method cannot handle such solitonic excitations, so radically new methods will need to be designed to handle them.
### Acknowledgements
The work presented in these notes results from many fruitful and enjoyable collaborations. We would like to thank R. Chitra, H. Maurey, B. S. Shastry and specially P. Le Doussal. Most importantly, nothing would have started without an initial collaboration with H. J. Schulz, to the memory of whom we would like to dedicate these lecture notes.
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# Isolated Horizons: Hamiltonian Evolution and the First Law
## I Introduction
The zeroth and first laws of black hole mechanics refer to equilibrium situations and small departures therefrom. The standard treatments restrict themselves to stationary space-times admitting event horizons and small perturbations from stationarity. While this simple idealization is a natural starting point, from physical considerations it seems overly restrictive. (See and especially for a detailed, critical discussion.) A framework, which is tailored to more realistic physical situations was introduced in and the zeroth and first laws were extended to it in . This analysis generalizes black hole mechanics in two directions. First, the notion of event horizons is replaced by that of ‘isolated horizons’. While the former can only be defined retroactively, requiring access to the entire space-time history, the latter can be defined quasi-locally. Second, the underlying space-time need not admit any Killing field; isolated horizons need not be Killing horizons. The static event horizons normally used in black hole mechanics and the cosmological horizons in de Sitter space-times are all special cases of isolated horizons. Furthermore, since space-times can now admit gravitational and matter radiation, there is a large class of other examples.
The framework developed in for generalizing black hole mechanics was based on two restrictive assumptions. First, only undistorted, non-rotating horizons were considered. That is, the boundary conditions used in implied that the intrinsic 2-metric of the horizon is spherically symmetric<sup>*</sup><sup>*</sup>*However, it allowed space-times, such as the Robinson-Trautman solutions, in which there is no space-time Killing field whatsoever in any neighborhood of the isolated horizon . and that the imaginary part of the Weyl tensor component $`\mathrm{\Psi }_2`$ — which encodes the angular momentum — vanishes. Second, while a rather general class of matter fields was allowed, it was assumed that the only relevant charges — i.e., hair — are the standard electromagnetic ones. The second assumption was weakened in which allowed dilaton couplings and Yang-Mills fields. In this paper, we allow for distortion and more general matter sources. Distortion plays an important role in several astrophysical situations, e.g., in problems involving black holes immersed in external fields, or surrounded by matter rings, and especially in the problem of black hole collisions. Post-Newtonian considerations suggest that, during black hole coalescence, individual horizons are distorted due to the Coulomb attraction even in the regime in which the black holes are sufficiently far from each other for the gravitational radiation falling into their horizons to be negligible. This phenomenon is also seen in numerical simulations.
The extensions which incorporated dilatonic and Yang-Mills charges did not involve a significant generalization of the basic framework developed in . The present paper, on the other hand, does. We begin with substantially weaker boundary conditions (formulated in terms of (real) tensor fields rather than the spinors used in ), and show that they imply constancy of surface gravity (and electro-static potential) on the horizon. This property turns out to be necessary and sufficient for the usual action principle of tetrad gravity to continue to be valid in presence of isolated horizons. The action leads to a covariant phase space, constructed from solutions to the field equations. Ref , by contrast, used the canonical phase space based on spinorial variables which is tailored for quantization but which contains technical complications that are unnecessary to the classical mechanics of isolated horizons. Up to this point, distortion and rotation are both incorporated. However, in the last step, i.e., in the discussion of the first law, we restrict ourselves to non-rotating horizons. (Rotation is incorporated in .)
To formulate the first law one must first define the energy $`E_\mathrm{\Delta }`$ associated with any isolated horizon $`\mathrm{\Delta }`$. Since there can be radiation in the spacetime outside isolated horizons, the ADM energy $`E_{\mathrm{ADM}}`$ is not a good measure of $`E_\mathrm{\Delta }`$ . Instead, as in the work of Brown and York , the strategy is to define the energy of the horizon using a Hamiltonian framework. Experience with the phase space formulation of general relativity suggests that, in the presence of boundaries, the Hamiltonian $`H_t`$ generating time-translation along a suitable vector field $`t^a`$ acquires surface terms. The idea is to define $`E_\mathrm{\Delta }`$ as the surface term at $`\mathrm{\Delta }`$ in the Hamiltonian. The key issue then is that of selecting the ‘appropriate’ time-translation $`t^a`$. Since one expects the volume term in the expression of $`H_t`$ to be a linear combination of constraints and thus vanish when evaluated on solutions, the problem reduces to that of specifying the boundary values of $`t^a`$ (or, equivalently, of the lapse and shift fields). The conditions at infinity are obvious and, in any case, will not affect the surface term at $`\mathrm{\Delta }`$. Thus, we need to focus only on the boundary value of $`t^a`$ on $`\mathrm{\Delta }`$.
In the non-rotating case it is clear that, at the horizon, $`t^a`$ should be proportional to the null normal to $`\mathrm{\Delta }`$. However, our boundary conditions do not select the null normal uniquely; there is a freedom to rescaleThis is not surprising since this freedom exists already on Killing horizons. If the space-time is asymptotically flat and admits a static Killing field globally, one can eliminate this freedom by restricting oneself to that Killing field which is unit at infinity. However, this strategy is not available if there is radiation in the exterior region, or, as in the static solutions representing distorted black holes, the metric fails to be asymptotically flat. the normal by a constant (on $`\mathrm{\Delta }`$) which can vary from one space-time to another. This freedom is physically important because, amongst other things, the surface gravity is sensitive to it. Suppose we fix this freedom by tying the boundary value of $`t^a`$ to fields on the horizon, e.g., by demanding that surface gravity be a specific function of the area and charges. We can then ask whether the time-evolution generated by this $`t^a`$ preserves the symplectic structure. It turns out that the answer is not always in the affirmative. On the horizon, the evolution vector field $`t^a`$ and the electromagnetic potential $`A_a`$ have to be tied to the horizon parameters appropriately. These conditions impose a constraint on the surface gravity $`\kappa _{(t)}`$ and the electric potential $`\mathrm{\Phi }_{(t)}`$. Somewhat surprisingly, the constraint is precisely the first law $`\delta E_\mathrm{\Delta }^t=(\kappa _{(t)}/8\pi G)\delta a_\mathrm{\Delta }+\mathrm{\Phi }_{(t)}\delta Q_\mathrm{\Delta }`$. Thus the evolution defined by $`t^a`$ is Hamiltonian if and only if the first law holds.In the undistorted context, while this role of the first law was known to the authors of , its importance was not fully appreciated. The importance was noticed independently in and used effectively in to extract physical information on spherical black holes with Yang-Mills hair. In this sense, the first law is even more fundamental than it is generally taken to be. Conceptually, this is perhaps the most striking feature of the present framework.
The requirement that the first law hold is not sufficient to fix $`t^a`$ uniquely. Although every $`t^a`$ must be a null normal to $`\mathrm{\Delta }`$, the rescaling of $`t^a`$ from one spacetime to another can depend on the horizon parameters and the first law does not fully determine this parameter dependence. There is an infinite family of parameter-dependent vector fields $`t^a`$ each defining a consistent Hamiltonian evolution and a horizon energy $`E_\mathrm{\Delta }^t`$. By contrast, at infinity all these vector fields must tend to a time-translation of an universal flat metric, used in the construction of the phase space. Hence, there is only the familiar 3-parameter freedom in the definition of $`E_{\mathrm{ADM}}^t`$, associated with the choice of a rest-frame. Furthermore, in any one space-time, we can eliminate it by simply going to the rest frame and thus extract the total mass $`M_{\mathrm{ADM}}`$ of the system. Although not necessary for mechanics, it is natural to ask if one can define a similar notion of mass of isolated horizons. The answer is in the affirmative in the Einstein-Maxwell theory. Let us require that $`t^a`$ should not only lead to a consistent Hamiltonian evolution but agree, on static solutions, with the static Killing field which is unit at infinity. Then the horizon value of $`t^a`$ is uniquely determined for all space-times in the phase-space. There is a preferred notion of time-translation, say $`t_o^a`$. We can set $`M_\mathrm{\Delta }=E_\mathrm{\Delta }^{t_o}`$ and regard $`M_\mathrm{\Delta }`$ as the mass of an isolated horizon $`\mathrm{\Delta }`$. In the earlier work on non-distorted horizons , the discussion of the first law was carried out only in the context of these preferred evolution vector fields $`t_o^a`$. That derivation is more closely tied to the traditional discussion of the laws in the static context but is not necessary from the more general perspective of isolated horizons. Nonetheless, the availability of a canonical definition of the mass $`M_\mathrm{\Delta }`$ is useful for other applications of this framework, e.g. to numerical relativity.
The paper is organized as follows. In Section II, we specify the boundary conditions defining general isolated horizons, allowing both distortion and rotation. We explain the role of these conditions, compare them with those used in and work out their consequences, including the zeroth law. In Section III we introduce the Lagrangian framework based on tetrads and (real) Lorentz connections and in Section IV, the covariant phase space. The first law is discussed in Section V. Upto this point, the focus is on the Einstein-Maxwell theory (although incorporation of the dilaton is straightforward). In Section VI we extend the framework to incorporate Yang-Mills fields. The horizon mass is introduced in Section VII and subtleties associated with the dilaton and Yang-Mills fields are discussed. For the convenience of readers who may not be familiar with distorted black holes, Appendix A presents a variety of examples and, for convenience of readers who work in the Newman-Penrose formalism, Appendix B summarizes the structure of isolated horizons in that framework.
We have attempted to make this paper self-contained in terms of methodology and technical details. However, the motivation behind isolated horizons and certain properties of our Hamiltonian are the same as those discussed in detail in . Since the inclusion of distortion does not add anything substantial to these issues, we have refrained from repeating that discussion in this paper.
## II Structure of isolated horizons and the zeroth law
In this section, we will introduce the basic definitions of isolated horizons and analyze their immediate consequences. The definitions will become progressively stronger. However, even the strongest boundary conditions are significantly weaker than requiring the horizon to be a Killing horizon for a local Killing vector field. By proceeding in steps, we will be able to keep track of the precise assumptions that are needed to obtain various results. Also, the availability of a hierarchy of definitions will be useful in other applications —such as numerical relativity and quantum gravity— which lie beyond the scope of the present paper.
Let us begin by introducing some notation. Throughout this paper we assume that all manifolds and fields under consideration are smooth. Let $``$ be a 4-manifold equipped with a metric $`g_{ab}`$ of signature $`(,+,+,+)`$. Let $`\mathrm{\Delta }`$ be a null hypersurface of $`(,g_{ab})`$. A future directed null normal to $`\mathrm{\Delta }`$ will be denoted by $`\mathrm{}`$. (In this paper, the term ‘null normal’ will always refer to a future directed null normal.) Let $`q_{ab}\widehat{}=g_{\begin{array}{c}ab\end{array}}`$ be the degenerate intrinsic metric on $`\mathrm{\Delta }`$.<sup>§</sup><sup>§</sup>§Equalities which hold only at $`\mathrm{\Delta }`$ will be denoted by ‘$`\widehat{}=`$ ’ and the pullback of a covariant index will be denoted by an arrow under that index, e.g. $`\omega _\begin{array}{c}a\end{array}`$ will denote the pullback of the $`1`$-form $`\omega _a`$ to $`\mathrm{\Delta }`$. A tensor $`q^{ab}`$ on $`\mathrm{\Delta }`$ will be called an ‘inverse’ of $`q_{ab}`$ if it satisfies $`q^{ab}q_{ac}q_{bd}\widehat{}=q_{cd}`$. Thus $`q^{ab}`$ is unique only up to addition of terms of the form $`\mathrm{}^{(a}V^{b)}`$ for some vector field $`V`$ tangential to $`\mathrm{\Delta }`$. The expansion $`\theta _{(\mathrm{})}`$ of a specific null normal $`\mathrm{}`$ is defined by $`\theta _{(\mathrm{})}=q^{ab}_a\mathrm{}_b`$, where $`_a`$ is the derivative operator compatible with $`g_{ab}`$. It is straightforward to check that $`\theta _{(\mathrm{})}`$ is independent of the choice of $`q^{ab}`$. With this structure at hand, we can now introduce our first definition.
### A Non-expanding Horizons
Definition 1: A $`3`$-dimensional sub-manifold $`\mathrm{\Delta }`$ of a space-time $`(,g_{ab})`$ is said to be a non-expanding horizon if it satisfies the following conditions:
* (i) $`\mathrm{\Delta }`$ is topologically $`S^2\times \text{I}\text{R}`$ and null;
* (ii) The expansion $`\theta _{(\mathrm{})}`$ of $`\mathrm{}`$ vanishes on $`\mathrm{\Delta }`$ for any null normal $`\mathrm{}`$;
* (iii) All equations of motion hold at $`\mathrm{\Delta }`$ and the stress-energy tensor $`T_{ab}`$ of matter fields at $`\mathrm{\Delta }`$ is such that $`T_b^a\mathrm{}^b`$ is future directed and causal for any future directed null normal $`\mathrm{}`$.
Note that if conditions (ii) and (iii) hold for one null normal $`\mathrm{}`$ they hold for all.
The role of these conditions is as follows. The restriction on topology is geared to the structure of horizons that result from gravitational collapse. However, it can be weakened. One can retain the requirement that the horizon have compact cross-sections but replace $`S^2`$ by a manifold with higher genus. Our main analysis will extend to this case in a straightforward manner. More generally, we can allow $`\mathrm{\Delta }`$ to have non-compact cross-sections, as for example in the case of certain acceleration horizons. The results presented in this section, including our derivation of the zeroth law will go through. However, since such horizons extend to infinity, our Hamiltonian framework will have to be modified appropriately. Finally, one could envisage incorporation of NUT charge. This extension would be even more subtle because, if all fields are smooth, $`\mathrm{\Delta }`$ would be topologically $`S^3`$ and $`\mathrm{}`$ would provide a Hopf fibration. In this case, $`\mathrm{\Delta }`$ would not admit any cross-sections which are everywhere transverse to $`\mathrm{}`$. This extension will be discussed elsewhere.
Requirement (iii) is analogous to the dynamical conditions one imposes at infinity. While at infinity one requires that the metric (and other fields) approach a specific solution to the field equations (the ‘classical vacuum’) , at the horizon we only ask that the field equations be satisfied. The energy condition involved is very weak; it is implied by the (much stronger) dominant energy condition that is typically used. Thus, the first and the last conditions are quite tame.
The key condition is (ii). It implies, in particular, that the horizon area is constant ‘in time’ and thus incorporates the idea that the horizon is isolated without having to assume the existence of a Killing field. We will denote the area by $`a_\mathrm{\Delta }`$ and refer to $`R_\mathrm{\Delta }`$ defined by $`a_\mathrm{\Delta }=4\pi R_\mathrm{\Delta }^2`$ as the horizon radius. All these conditions are satisfied on any Killing horizon (with 2-sphere sections) if gravity is coupled to physically reasonable matter (including perfect fluids, Klein-Gordon fields, Maxwell fields possibly with dilatonic coupling, Yang-Mills fields).
Although the conditions in the definition are quite weak, they have surprisingly rich consequences. We will now discuss them in detail. In some of this analysis it is convenient to use a null-tetrad and the associated Newman-Penrose quantities (see Appendix B and references therein). Given a specific null normal field $`\mathrm{}^a`$ to $`\mathrm{\Delta }`$, we can introduce a complex null vector field $`m^a`$ tangential to $`\mathrm{\Delta }`$ and a real, future directed null field $`n^a`$ transverse to $`\mathrm{\Delta }`$ so that the following relations hold: $`n\mathrm{}=1,m\overline{m}=1`$ and all other scalar products vanish. The quadruplet $`(\mathrm{},n,m,\overline{m})`$ constitutes a null-tetrad. There is of course an infinite number of null tetrads compatible with a given $`\mathrm{}`$, related to one another by restricted Lorentz rotations. Our conclusions will not be sensitive to this gauge-freedom.
(a) Properties of $`\mathrm{}`$: Since $`\mathrm{}^a`$ is a null normal to $`\mathrm{\Delta }`$, it is automatically twist free and geodesic. We will denote the acceleration of $`\mathrm{}^a`$ by $`\kappa _{(\mathrm{})}`$
$$\mathrm{}^a_a\mathrm{}^b\widehat{}=\kappa _{(\mathrm{})}\mathrm{}^b.$$
(II.1)
Note that the acceleration is a property not of the horizon $`\mathrm{\Delta }`$ itself, but of a specific null normal to it: if we replace $`\mathrm{}`$ by $`\mathrm{}^{}=f\mathrm{}`$, then the acceleration changes via $`\kappa _{(\mathrm{}^{})}=f\kappa _{(\mathrm{})}+_{\mathrm{}}f`$.
Since the twist of $`\mathrm{}`$ vanishes, the Raychaudhuri equation implies:
$`_{\mathrm{}}\theta _{(\mathrm{})}\widehat{}=\kappa _{(\mathrm{})}\theta _{(\mathrm{})}{\displaystyle \frac{1}{2}}\theta _{(\mathrm{})}^2\sigma \overline{\sigma }R_{ab}\mathrm{}^a\mathrm{}^b`$
where $`\sigma =m^am^b_a\mathrm{}_b`$ is the shear of $`\mathrm{}`$ in the given null tetrad. Since $`\theta _{(\mathrm{})}`$ vanishes on $`\mathrm{\Delta }`$, we conclude: $`\sigma \overline{\sigma }+R_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=0`$. The condition on the stress-energy tensor ensures that $`R_{ab}\mathrm{}^a\mathrm{}^b=8\pi GT_{ab}\mathrm{}^a\mathrm{}^b`$ is non-negative on $`\mathrm{\Delta }`$. Hence, we conclude:
$$\sigma \widehat{}=0,\mathrm{and}R_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=0.$$
(II.2)
Thus, in particular, every null normal $`\mathrm{}`$ is free of expansion, twist and shear.
(b) Conditions on the Ricci tensor: The second equation in (II.2) implies that the vector $`R_{}^{a}{}_{b}{}^{}\mathrm{}^b`$ is tangential to $`\mathrm{\Delta }`$. The energy condition and the field equations imply this vector must also be future causal. This means that $`R^a{}_{b}{}^{}\mathrm{}_{}^{b}`$ must be proportional to $`\mathrm{}^a`$ and hence, $`R_{\begin{array}{c}a\end{array}b}\mathrm{}^b=0`$. In the Newman-Penrose formalism this condition translates to:
$$\mathrm{\Phi }_{00}=\frac{1}{2}R_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=0\mathrm{and}\mathrm{\Phi }_{01}=\overline{\mathrm{\Phi }}_{10}=\frac{1}{2}R_{ab}\mathrm{}^am^b\widehat{}=0.$$
(II.3)
Since this statement is equivalent to $`R_{\begin{array}{c}a\end{array}b}\mathrm{}^b=0`$, it is gauge invariant, i.e. it does not depend upon the specific choice of null normal $`\mathrm{}`$ and $`m`$.
(c) A natural connection 1-form on $`\mathrm{\Delta }`$: Since $`\mathrm{}`$ is expansion, shear and twist free, there exists a one-form $`\omega _a`$ intrinsic to $`\mathrm{\Delta }`$ such that
$$_\begin{array}{c}a\end{array}\mathrm{}^b\widehat{}=\omega _a\mathrm{}^b$$
(II.4)
which in turn implies
$`_{\mathrm{}}q_{ab}\widehat{}=2\begin{array}{c}_a\mathrm{}_b\hfill \end{array}\widehat{}=0.`$
Thus, every null normal $`\mathrm{}`$ is a ‘Killing field’ of the degenerate metric on $`\mathrm{\Delta }`$. Furthermore, we will now show that
$${}_{}{}^{2}ϵ:=im\overline{m}$$
(II.5)
is also invariantly defined. Since $`_{\mathrm{}}q_{ab}=0`$, the space $`𝒮`$ of integral curves of $`\mathrm{}`$ is naturally equipped with a non-degenerate metric $`\underset{¯}{q}_{ab}`$ (so that $`q_{ab}`$ on $`\mathrm{\Delta }`$ is the pull-back of $`\underset{¯}{q}_{ab}`$). Denote by $`\underset{¯}{ϵ}_{ab}`$ the unique (up to orientation) unit alternating tensor on $`(𝒮,\underset{¯}{q}_{ab})`$. $`{}_{}{}^{2}ϵ`$ is the pull-back to $`\mathrm{\Delta }`$ of $`\underset{¯}{ϵ}`$. Although $`\mathrm{}`$ is a ‘Killing field’ of the intrinsic horizon geometry, the space-time metric $`g_{ab}`$ need not admit a Killing field in any neighborhood of $`\mathrm{\Delta }`$. Robinson-Trautman metrics and the Kastor-Traschen solutions provide explicit examples of this type.
The 1-form $`\omega `$ will play an important role throughout this paper. It has an interesting geometrical interpretation. We can regard $`\omega `$ as a connection on the line bundle $`T\mathrm{\Delta }^{}`$ over $`\mathrm{\Delta }`$ whose fibers are the 1-dimensional null normals to $`\mathrm{\Delta }`$. Under the rescalings $`\mathrm{}\stackrel{~}{\mathrm{}}=f\mathrm{}`$, of the null normal $`\mathrm{}`$, it transforms via:
$$\omega _a\stackrel{~}{\omega }_a=\omega _a+_\begin{array}{c}a\end{array}\mathrm{ln}f.$$
(II.6)
(d) Induced Connection on $`\mathrm{\Delta }`$: Each metric submanifold $`M`$ of $``$ admits a natural connection — one which is torsion-free and compatible with the induced metric on $`M`$. This connection is also canonically induced by the space-time connection $``$. However, since the induced metric $`q_{ab}`$ on $`\mathrm{\Delta }`$ is degenerate, there exist infinitely many connections compatible with it. A general null sub-manifold inherits a unique (torsion-free) derivative operator $`𝒟`$ from $``$ if and only if its null normal $`\mathrm{}`$ satisfies $`\begin{array}{c}_a\mathrm{}_b\hfill \end{array}=0`$. Therefore, the conditions imposed in Definition 1 guarantee that every non-expanding horizon has a unique intrinsic derivative operator $`𝒟`$. The action of $`𝒟`$ on a vector field $`X^a`$ tangent to $`\mathrm{\Delta }`$ and on a 1-form $`\eta _a`$ intrinsic to $`\mathrm{\Delta }`$ is given by
$`𝒟_aX^b\widehat{}=\begin{array}{c}_a\hfill \end{array}\stackrel{~}{X}^b\text{and}𝒟_a\eta _b\widehat{}=\begin{array}{c}_a\stackrel{~}{\eta }_b\hfill \end{array}.`$
where $`\stackrel{~}{X}^b`$ and $`\stackrel{~}{\eta }_b`$ are arbitrary extensions of $`X^b`$ and $`\eta _b`$ to the full space-time manifold $``$. It is easy to show that $`𝒟`$ is independent of the extensions.
The 1-form $`\omega `$ captures only part of the information in $`𝒟`$. The full connection $`𝒟`$ on $`\mathrm{\Delta }`$ plays an important role in extracting physics in the strong field regime near $`\mathrm{\Delta }`$ . However, it is not essential to the discussion of isolated horizon mechanics.
(e) Conditions on the Weyl tensor: Let us begin with the definition of the Riemann tensor, $`[_a_b_b_a]X^c=2R_{abd}^{}{}_{}{}^{c}X^d`$. If we set $`X^c=\mathrm{}^c`$ and pull back the indices $`a`$ and $`b`$, then using (II.4), we obtain:
$$[𝒟_a\omega _b𝒟_b\omega _a]\mathrm{}^c\widehat{}=2R_{\begin{array}{c}abd\end{array}}{}_{}{}^{c}\mathrm{}_{}^{d}\widehat{}=2C_{\begin{array}{c}abd\end{array}}{}_{}{}^{c}\mathrm{}_{}^{d}$$
(II.7)
where $`C_{abc}^{}{}_{}{}^{d}`$ is the Weyl tensor. The last equality follows from $`R_{\begin{array}{c}ab\end{array}}\mathrm{}^b\widehat{}=0`$. Thus, if $`v`$ is any $`1`$-form on $`\mathrm{\Delta }`$ satisfying $`v\mathrm{}\widehat{}=0`$, contracting the previous equation with $`v_c`$ we get
$`C_{\begin{array}{c}abd\end{array}}{}_{}{}^{c}v_{c}^{}\mathrm{}^d\widehat{}=0.`$
Let us choose a null tetrad and set $`v`$ to be $`m`$ or $`\overline{m}`$. Then
$$\mathrm{\Psi }_0:=C_{abcd}\mathrm{}^am^b\mathrm{}^cm^d\widehat{}=0\text{and}\mathrm{\Psi }_1:=C_{abcd}\mathrm{}^am^b\mathrm{}^cn^d=C_{abcd}\mathrm{}^am^b\overline{m}^cm^d\widehat{}=0,$$
(II.8)
where we have used the trace-free property of the Weyl tensor in the second equation. It is also clear that equations (II.8) are independent of which null normal $`\mathrm{}`$, and vector fields $`m`$ and $`\overline{m}`$ we choose to construct the null tetrad; equation (II.8) is gauge invariant.
(f) Curvature of $`\omega `$: Let us contract (II.7) with $`n_c`$ and use $`\mathrm{}^an_a=1`$. Then we have:
$$2𝒟_{[a}\omega _{b]}\widehat{}=C_{\begin{array}{c}abd\end{array}}^{}{}_{}{}^{c}\mathrm{}^dn_c\widehat{}=C_{\begin{array}{c}abc\end{array}d}\mathrm{}^cn^d.$$
(II.9)
Expanding the Weyl tensor in terms of the $`\mathrm{\Psi }`$’s, one obtains
$`C_{abcd}\mathrm{}^cn^d`$ $`\widehat{}=`$ $`4(\mathrm{Re}\left[\mathrm{\Psi }_2\right])n_{[a}l_{b]}+2\mathrm{\Psi }_3\mathrm{}_{[a}m_{b]}+2\overline{\mathrm{\Psi }}_3\mathrm{}_{[a}\overline{m}_{b]}`$ (II.11)
$`2\overline{\mathrm{\Psi }}_1n_{[a}m_{b]}2\mathrm{\Psi }_1n_{[a}\overline{m}_{b]}+4i(\mathrm{Im}\left[\mathrm{\Psi }_2\right])m_{[a}\overline{m}_{b]}.`$
where
$$\mathrm{\Psi }_2:\widehat{}=C_{abcd}\mathrm{}^am^b\overline{m}^cn^d\text{and}\mathrm{\Psi }_3:\widehat{}=C_{abcd}\mathrm{}^an^b\overline{m}^cn^d.$$
(II.12)
Substituting this expression into (II.9), pulling back on the two free indices and taking into account (II.8) and (II.5), we obtain
$$d\omega \widehat{}=2(\mathrm{Im}\left[\mathrm{\Psi }_2\right]){}_{}{}^{2}ϵ.$$
(II.13)
This relation will play an important role in what follows. Note that, because $`\mathrm{\Psi }_0`$ and $`\mathrm{\Psi }_1`$ vanish on $`\mathrm{\Delta }`$, $`\mathrm{\Psi }_2`$ is gauge invariant.
Remark: It is interesting to compare the structure of $`\mathrm{\Delta }`$ with that of null-infinity $``$ (in the usual conformal gauge, in which the conformal factor is chosen such that the null-normal to $``$ is divergence-free). Both are null surfaces and can be regarded as ‘line bundles’ over a base space $`𝒮`$ of the integral curves of null normals. (For brevity, we will ignore a caveat concerning completeness of fibers.) The null normals are Killing fields of the intrinsic degenerate metric so that this metric is the pull-back to the 3-surface of a positive-definite metric on $`𝒮`$. In both cases, the space-time connection induces an intrinsic connection on the 3-surface . These connections capture physically important information. However, there are a number of differences as well. Since $``$ is constructed by a conformal completion, the conformal freedom permeates all geometric structures at $``$. In particular, given a physical space-time, the intrinsic metric and the derivative operator are known only up to conformal transformations. On the other hand, since $``$ is at infinity, in some ways its structure is both more rigid and simpler. First, without loss of generality, we can assume that the metric on $`𝒮`$ is a 2-sphere metric; the issue of distortion is physically irrelevant at $``$. Second, the Weyl tensor vanishes identically at $``$ and the curvature of the intrinsic connection captures non-trivial information about the next order space-time curvature. By contrast, at $`\mathrm{\Delta }`$ only four components of the Weyl curvature vanish and four other components are coded in the curvature of the intrinsic connection on $`\mathrm{\Delta }`$. In spite of these differences, one can carry over some techniques from null infinity to extract physical information about isolated horizons. In particular, using the analogs of techniques which have been successful at $``$, one can introduce preferred cross sections of and Bondi-type expansions near $`\mathrm{\Delta }`$ .
This concludes our analysis of the consequences of the boundary conditions defining non-expanding horizons. Note that, even though $`\mathrm{}`$ is a Killing field for the intrinsic, degenerate metric $`q_{ab}`$ on $`\mathrm{\Delta }`$, it is not an infinitesimal symmetry for other geometrical fields such as the intrinsic connection $`𝒟`$ or components of the curvature tensor. In the next sub-section, we will make the structure more rigid by suitably restricting the choice of $`\mathrm{}`$.
### B Weakly Isolated horizons
The time-independence of the intrinsic metric $`q_{ab}`$ captures the idea that $`\mathrm{\Delta }`$ is isolated in a suitable sense. While this condition has rich consequences, the resulting structure is still not sufficient for physical applications. In particular, since $`\mathrm{}`$ can be rescaled by an arbitrary positive definite function, the acceleration $`\kappa _{(\mathrm{})}`$ is not necessarily constant on $`\mathrm{\Delta }`$. Therefore, we need to impose additional restrictions on the physical fields at $`\mathrm{\Delta }`$ to establish the zeroth law. Since $`\mathrm{}`$ is already a symmetry of the intrinsic metric, it is natural to require it also be a symmetry of the ‘extrinsic curvature’. However, the standard definition of the extrinsic curvature is not applicable to null surfaces. Nonetheless, given a null normal $`\mathrm{}`$, we can construct a tensor field $`K_a{}_{}{}^{b}:=𝒟_a\mathrm{}^b`$, defined intrinsically on $`\mathrm{\Delta }`$, which can be thought of as the analogue of the extrinsic curvature.We are grateful to Thibault Damour for pointing out that $`K_a^b`$ is called the Weingarten map and is analogous to extrinsic curvature. This comment suggested the above motivation for our condition on the connection 1-form $`\omega `$. For an alternate, and in a sense weaker, condition see the remark at the end of Section II D. From the viewpoint of intrinsic structures on $`\mathrm{\Delta }`$ discussed in Section II A, it is perhaps more natural to ask that $`\mathrm{}`$ be a symmetry of the full intrinsic connection $`𝒟`$ (see section II D and .) However, this stronger condition is not necessary for the laws of mechanics discussed here. Indeed, on a metric sub-manifold, if we replace $`\mathrm{}`$ by the unit normal, $`K_a^b`$ is precisely the extrinsic curvature. It is then natural to demand that, on an isolated horizon, $`K_a^b`$ also be time-independent: $`_{\mathrm{}}K_a{}_{}{}^{b}\widehat{}=0`$. As a consequence of (II.4), this is equivalent to imposing $`_{\mathrm{}}\omega \widehat{}=0`$.
Let us examine the above condition. As we will show at the end of this section, given a non-expanding horizon we can always find a null normal $`\mathrm{}^a`$ which satisfies $`_{\mathrm{}}\omega \widehat{}=0`$. The behavior of this condition under rescalings of $`\mathrm{}`$ is complicated by the fact that $`\omega `$ itself depends upon the choice of null normal (see equation (II.6)). However, under a constant rescaling $`\mathrm{}\stackrel{~}{\mathrm{}}=c\mathrm{}`$, the connection 1-form $`\omega `$ is unchanged. Therefore, if $`\mathrm{}`$ satisfies the condition $`_{\mathrm{}}\omega \widehat{}=0`$, so does any $`\stackrel{~}{\mathrm{}}`$ related to $`\mathrm{}`$ by constant rescaling. This suggests we introduce an equivalence relation: Two future-directed null normals $`\mathrm{}`$ and $`\stackrel{~}{\mathrm{}}`$ belong to the same equivalence class $`[\mathrm{}]`$ if and only if $`\stackrel{~}{\mathrm{}}=c\mathrm{}`$ for some positive constant $`c`$.
The above considerations lead us to the following definition:
Definition 2: A weakly isolated horizon $`(\mathrm{\Delta },[\mathrm{}])`$ consists of a non-expanding horizon $`\mathrm{\Delta }`$, equipped with an equivalence class $`[\mathrm{}]`$ of null normals to it satisfying
$$_{\mathrm{}}\omega \widehat{}=0\text{for all}\mathrm{}[\mathrm{}].$$
(II.14)
As pointed out above, if this last equation holds for one $`\mathrm{}`$, it holds for all $`\mathrm{}`$ in $`[\mathrm{}]`$.
A Killing horizon (with 2-sphere cross-sections) is automatically a weakly isolated horizon, (provided the matter fields satisfy the energy condition of Definition 1). Given a non-expanding horizon $`\mathrm{\Delta }`$, one can always find an equivalence class $`[\mathrm{}]`$ of null-normals such that $`(\mathrm{\Delta },[\mathrm{}])`$ is a weakly isolated horizon. However, condition (II.14) does not by itself single out the appropriate equivalence class $`[\mathrm{}]`$. As indicated in Section II D, one can further strengthen the boundary conditions and provide a specific prescription to select the equivalence class $`[\mathrm{}]`$ uniquely. However, for mechanics of isolated horizons, these extra steps are unnecessary. In particular, our analysis will not depend on how the equivalence class $`[\mathrm{}]`$ is chosen. The adverb ‘weakly’ in Definition 2 emphasizes this point.
The condition (II.14) has several consequences which are relevant for this paper.
(a) Surface gravity: In the case of Killing horizons $`\mathrm{\Delta }_\mathrm{K}`$, surface gravity is defined as the acceleration of the Killing field $`\xi `$ normal to $`\mathrm{\Delta }_\mathrm{K}`$. However, if $`\mathrm{\Delta }_\mathrm{K}`$ is a Killing horizon for $`\xi `$, it is also a Killing horizon for $`c\xi `$ for any positive constant $`c`$. Hence, surface gravity is not an intrinsic property of $`\mathrm{\Delta }_\mathrm{K}`$, but depends also on the choice of a specific Killing field $`\xi `$. (Of course the result that the surface gravity is constant on $`\mathrm{\Delta }_\mathrm{K}`$ is insensitive to this rescaling freedom.) In asymptotically flat space-times admitting global Killing fields, this ambiguity is generally resolved by selecting a preferred normalization in terms of the structure at infinity. For example, in the static case, one requires the Killing field $`\xi `$ to be unit at infinity. However, in absence of a global Killing field or asymptotic flatness, this strategy does not work and we simply have to accept the constant rescaling freedom in the definition of surface gravity. In the context of isolated horizons, then, it is natural to keep this freedom.
A weakly isolated horizon is similarly equipped with a preferred family $`[\mathrm{}]`$ of null normals, unique up to constant rescalings. Therefore, it is natural to interpret $`\kappa _{(\mathrm{})}`$ as the surface gravity associated with $`\mathrm{}`$. Under the permissible rescalings $`\mathrm{}\stackrel{~}{\mathrm{}}=c\mathrm{}`$, the surface gravity transforms via: $`\kappa _{(\stackrel{~}{\mathrm{}})}=c\kappa _{(\mathrm{})}`$. Thus, while $`\omega `$ is insensitive to the rescaling freedom in $`[\mathrm{}]`$, $`\kappa _{(\mathrm{})}`$ captures this freedom fully. One can, if necessary, select a specific $`\mathrm{}`$ in $`[\mathrm{}]`$ by demanding that $`\kappa _{(\mathrm{})}`$ be a specific function of the horizon parameters which are insensitive to this freedom, e.g., by setting $`\kappa _{(\mathrm{})}=1/2R_\mathrm{\Delta }`$, where $`R_\mathrm{\Delta }`$ is the horizon radius (related to the horizon area $`a_\mathrm{\Delta }`$ via $`a_\mathrm{\Delta }=4\pi R_\mathrm{\Delta }^2`$).
(b) Zeroth law: The boundary conditions of Definition 2 allow us to define surface gravity $`\kappa _{(\mathrm{})}`$ of a weakly isolated horizon $`(\mathrm{\Delta },[\mathrm{}])`$. We will now show that the surface gravity is constant on $`\mathrm{\Delta }`$. In other words, the zeroth law holds for weakly isolated horizons.
Recall from (II.13) that on a non-expanding horizon, $`d\omega \widehat{}=2\mathrm{Im}\left[\mathrm{\Psi }_2\right]{}_{}{}^{2}ϵ`$ for any choice of null normal $`\mathrm{}`$. Since $`{}_{}{}^{2}ϵ`$ is the pull-back to $`\mathrm{\Delta }`$ of the alternating tensor $`\underset{¯}{ϵ}`$ on the space $`𝒮`$ (of orbits of $`\mathrm{}`$), clearly $`\mathrm{}{}_{}{}^{2}ϵ\widehat{}=0`$. Therefore,
$`\mathrm{}d\omega \widehat{}=0`$
for every null normal $`\mathrm{}`$. In particular, on a weakly isolated horizon this equation holds for any $`\mathrm{}[\mathrm{}]`$. Moreover, each of these restricted null normals also satisfies
$`0\widehat{}=_{\mathrm{}}\omega \widehat{}=d(\mathrm{}\omega )+\mathrm{}d\omega `$
Hence, we conclude:
$`d(\mathrm{}\omega )\widehat{}=d(\kappa _{(\mathrm{})})\widehat{}=0,`$
where we have used the definition (II.1) of $`\kappa _{(\mathrm{})}`$. Thus, surface gravity is constant on $`\mathrm{\Delta }`$.
Although this proof of the zeroth law appears extremely simple, the argument is not as trivial as it might first appear since we have used a number of consequences of the boundary conditions derived in section II A. In contrast to earlier derivations we do not require the presence of a Killing field even in a neighborhood of $`\mathrm{\Delta }`$. Therefore the proof applies also to space-times such as the Robinson-Trautman solutions which do not admit a Killing field. Also, $`\mathrm{\Delta }`$ need not be ‘complete’ — it may be of finite affine length with respect to any $`\mathrm{}`$ — and may not admit the analog of a ‘bifurcate surface’ on which the Killing field vanishes. Finally, the field equations are used rather weakly; we only need to assume that $`(R^a{}_{b}{}^{}\frac{1}{2}R\delta ^a{}_{b}{}^{})\mathrm{}^b`$ is a future directed causal vector.
Surface gravity does not have a definite value on a weakly isolated horizon. The value of $`\kappa _{(\mathrm{})}`$ depends upon the choice of null normal $`\mathrm{}[\mathrm{}]`$. Since all the normals $`\mathrm{}`$ to $`\mathrm{\Delta }`$ are future directed, the rescaling constant $`c`$ is necessarily positive. Therefore, if the surface gravity is non-zero (respectively, zero) with respect to one $`\mathrm{}`$, it is non-zero (respectively, zero) with respect to any other $`\mathrm{}[\mathrm{}]`$. This rescaling freedom is the same as the one discussed above in the context of Killing horizons.
We will conclude this sub-section with three remarks.
i) Freedom in the choice of $`[\mathrm{}]`$: Given a non-expanding horizon $`\mathrm{\Delta }`$, it is natural to ask if one can always select an equivalence class $`[\mathrm{}]`$ of null normals such that $`(\mathrm{\Delta },[\mathrm{}])`$ is a weakly isolated horizon. As indicated earlier in this section, the answer is in the affirmative and, furthermore, there is a considerable freedom in the choice of $`[\mathrm{}]`$. Let us examine this issue in some detail.
Since $`\mathrm{}d\omega \widehat{}=0`$ for any null normal $`\mathrm{}`$ to a non-expanding horizon, it follows that a null normal $`\mathrm{}`$ satisfies $`_{\mathrm{}}\omega \widehat{}=0`$ if and only if $`d\kappa _{(\mathrm{})}\widehat{}=0`$. Thus, to find a family $`[\mathrm{}]`$ required in the definition of weakly isolated horizons, it is necessary and sufficient to find a null normal $`\mathrm{}`$ for which the surface gravity is constant. On a non-expanding horizon, the surface gravity transforms as follows.
$`\mathrm{If}\mathrm{}\stackrel{~}{\mathrm{}}\widehat{}=f\mathrm{},\mathrm{then}\kappa _{(\stackrel{~}{\mathrm{}})}\widehat{}=f\kappa _{(\mathrm{})}+_{\mathrm{}}f.`$
Hence, starting with any $`\mathrm{}`$, we can simply solve for $`f`$ by requiring that $`\kappa _{(\stackrel{~}{l})}`$ be constant on $`\mathrm{\Delta }`$. The solution is not unique. If $`\kappa _{\mathrm{}}`$ is constant, given any non-zero function $`g`$ satisfying $`_{\mathrm{}}g\widehat{}=0`$ and a constant $`\stackrel{~}{\kappa }`$, let us set
$`f\widehat{}=ge^{\kappa _{(\mathrm{})}v}+{\displaystyle \frac{\stackrel{~}{\kappa }}{\kappa _{(\mathrm{})}}}`$
where $`v`$ satisfies $`_{\mathrm{}}v\widehat{}=1`$. Then, we obtain an $`\stackrel{~}{\mathrm{}}[\mathrm{}]`$ for which $`\kappa _{(\stackrel{~}{\mathrm{}})}\widehat{}=\stackrel{~}{\kappa }`$. This is the only freedom if both $`\kappa _{\mathrm{}}`$ and $`\kappa _{(\stackrel{~}{\mathrm{}})}`$ are to be constant. Thus, each non-expanding horizon gives rise to an infinite family of weakly isolated horizons. Put differently, although one can easily obtain weakly isolated horizons from non-expanding ones by choosing appropriate null normals $`[\mathrm{}]`$, a specific weakly isolated horizon carries much more information than the non-expanding horizon it comes from. At the end of Section II D, we will indicate how one can further strengthen the boundary conditions to give a prescription for selecting a specific $`[\mathrm{}]`$. However, the analysis of this paper does not depend on how this selection is made.
ii) How does Definition 2 compare with that used in the undistorted, non-rotating case? As one would expect, the definition given in is significantly stronger. Furthermore, it was tied to a foliation from the beginning. More precisely, it assumed that there exists a foliation to which $`\omega `$ is normal, with $`\omega =\kappa n`$, and it required that the expansion $`\mathrm{Re}\left[\mu \right]`$ of the null normal $`n`$ to the leaves of the foliation be constant on $`\mathrm{\Delta }`$ (see appendix B for definitions of the NP spin coefficients). Although it was shown that the foliation is unique if it exists, the heavy reliance on the foliation right from the beginning made that definition less elegant and its invariant content less transparent. Also, since we now allow $`d\omega `$ to be non-zero and impose no conditions on $`\mathrm{Re}\left[\mu \right]`$, we can now incorporate rotation and distortion.
iii) Alternate boundary conditions: In the definition of isolated horizons, we required $`_{\mathrm{}}\omega \widehat{}=0`$, which in particular implies $`_{\mathrm{}}\kappa _{(\mathrm{})}_{\mathrm{}}(\mathrm{}\omega )\widehat{}=0`$. Thus, the definition itself assured us that $`\kappa _{(\mathrm{})}`$ is ‘time independent’ and to prove the zeroth law we had to show that it is also independent of ‘angles’. Could we have used another definition in which the ‘time dependence’ of $`\kappa _{(\mathrm{})}`$ was not explicitly required but followed from other conditions? The answer is in the affirmative: In place of $`_{\mathrm{}}\omega \widehat{}=0`$, we could have required that $`\mathrm{\Delta }`$ admit a foliation on which the expansion $`\mathrm{Re}\left[\mu \right]`$ of $`n`$ and the Newman-Penrose spin coefficient $`\pi `$ which carries the angular momentum information are ‘time independent’ . Then the field equations would have implied that $`\kappa _{(\mathrm{})}`$ is time-independent. Furthermore, all results of this paper go through (and were in fact first obtained) with these modified boundary conditions. Note however that the new condition is neither weaker nor stronger than the one we used. Both require that $`\pi `$ be time independent. In addition, the present definition of isolated horizons requires that $`\kappa _{(\mathrm{})}`$ be time independent while the alternative definition would have required, instead, that $`\mathrm{Re}\left[\mu \right]`$ be time independent. In this paper, we chose the present definition because it can be stated without reference to a foliation.
### C Electromagnetic Field
We shall now describe the form of the electromagnetic field at an isolated horizon and introduce a partial gauge fixing at the horizon which will allow us to introduce the notion of an electric potential. In the next three sections — where we discuss the action, phase space and first law — we will assume that the only matter field present at the isolated horizon are Maxwell fields. However, as our discussion will make it clear, this restriction is made primarily for simplicity. The overall framework is rather general and can accommodate matter for which there exists a well-defined action principle and a (covariant) Hamiltonian framework. In particular, in Section VI, we describe how to extend the formalism to include Yang Mills Fields.
The isolated horizon boundary conditions restrict matter primarily through conditions on the stress-energy tensor $`T_{ab}`$. Let us begin with $`T_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=0`$, a direct consequence of the boundary conditions and Raychaudhuri equation. (This restriction arises due to the fact that $`\mathrm{\Delta }`$ is a non-expanding horizon; the subsequent stronger boundary conditions do not further constrain $`𝐅`$). Although this condition is weak, it turns out to have interesting consequences on the form of the electromagnetic field, $`𝐅`$, at $`\mathrm{\Delta }`$. The stress-energy tensor for electromagnetism is given in terms of the field strength $`𝐅`$ as
$$𝐓_{ab}=\frac{1}{4\pi }[𝐅_{ac}𝐅_{b}^{}{}_{}{}^{c}\frac{1}{4}g_{ab}𝐅_{cd}𝐅^{cd}].$$
(II.15)
Let us contract this expression with $`\mathrm{}^a\mathrm{}^b`$ and examine the consequences for $`𝐅`$. This gives
$$0\widehat{}=𝐓_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=\mathrm{}^am^b𝐅_{ab}^2,$$
(II.16)
where, to obtain the last expression, we have used the anti-symmetry of $`𝐅`$ and the fact that the metric at the horizon can be expressed in terms of a null-tetrad as $`g_{ab}=2\mathrm{}_{(a}n_{b)}+2m_{(a}\overline{m}_{b)}`$. An immediate consequence of (II.16) is that
$$\begin{array}{c}\mathrm{}^a𝐅_{ab}\hfill \end{array}\widehat{}=0.$$
(II.17)
In order to obtain a similar expression for $`{}_{}{}^{}𝐅`$ recall that the stress energy tensor can be rewritten as $`𝐓_{ab}=\frac{1}{4\pi }[{}_{}{}^{}𝐅_{ac}^{}{}_{}{}^{}𝐅_{b}^{}{}_{}{}^{c}\frac{1}{4}g_{ab}{}_{}{}^{}𝐅_{cd}^{}{}_{}{}^{}𝐅_{}^{cd}]`$. Applying the same argument which led to (II.17), we obtain a similar restriction on $`{}_{}{}^{}𝐅`$, namely
$$\begin{array}{c}\mathrm{}^a^{}𝐅_{ab}\hfill \end{array}\widehat{}=0.$$
(II.18)
These two restrictions tell us there is no flux of electromagnetic radiation across the horizon.
It is straightforward to show that (II.17), (II.18) and the form of the metric at $`\mathrm{\Delta }`$ place further restrictions on $`𝐓_{ab}`$:
$`\begin{array}{cccccc}\hfill 𝐓_{ab}\mathrm{}^am^b& \widehat{}=& 0\hfill & \hfill 𝐓_{ab}\mathrm{}^a\overline{m}^b& \widehat{}=& 0\hfill \\ \hfill 𝐓_{ab}m^am^b& \widehat{}=& 0\hfill & \hfill 𝐓_{ab}\overline{m}^a\overline{m}^b& \widehat{}=& 0.\hfill \end{array}`$
The first two equations contain no new information since we already knew from general arguments (see equation (II.3)), that $`\mathrm{\Phi }_{10}\widehat{}=0`$ and $`\mathrm{\Phi }_{01}\widehat{}=0`$. However, the last two equations do place further restrictions on the stress energy tensor. Since the equations of motion are enforced at the boundary, we see immediately that they are equivalent to further restricting the Ricci tensor at the horizon by requiring:
$$\mathrm{\Phi }_{02}:=\frac{1}{2}R_{ab}m^am^b\widehat{}=0\text{and}\mathrm{\Phi }_{20}:=\frac{1}{2}R_{ab}\overline{m}^a\overline{m}^b\widehat{}=0.$$
(II.19)
This result need not hold for general matter fields; it relies on the properties of the electromagnetic stress-energy tensor (II.15).
The next task is to define the electric and magnetic charges of the horizon. Since the horizon is an inner boundary of spacetime, the normal to a 2-sphere cross section of the horizon will naturally be inward pointing. Bearing this in mind, we define the electric and magnetic charges of the horizon as
$$Q_\mathrm{\Delta }:\widehat{}=\frac{1}{4\pi }_{S_\mathrm{\Delta }}{}_{}{}^{}𝐅\text{and}P_\mathrm{\Delta }:\widehat{}=\frac{1}{4\pi }_{S_\mathrm{\Delta }}𝐅.$$
(II.20)
For these definitions to be meaningful, the values of $`Q_\mathrm{\Delta }`$ and $`P_\mathrm{\Delta }`$ should be independent of the cross section of the horizon $`S_\mathrm{\Delta }`$. We will now show that $`\mathrm{\Delta }`$ being a non-expanding horizon guarantees this is the case. Let us first evaluate
$`_{\mathrm{}}\begin{array}{c}𝐅\hfill \end{array}\widehat{}=\mathrm{}\begin{array}{c}d𝐅\hfill \end{array}+\begin{array}{c}d(\mathrm{}𝐅)\hfill \end{array}.`$
The first term on the right hand side vanishes due to Maxwell’s equations on $`\mathrm{\Delta }`$, while the second term is zero due to the previous restriction on $`𝐅`$, (II.17). Therefore we conclude that $`𝐅`$ is Lie dragged by $`\mathrm{}`$. An identical argument for $`{}_{}{}^{}𝐅`$ leads to the analogous conclusion. Therefore we obtain
$$_{\mathrm{}}\begin{array}{c}𝐅\hfill \end{array}\widehat{}=0\text{and}_{\mathrm{}}\begin{array}{c}^{}𝐅\hfill \end{array}\widehat{}=0.$$
(II.21)
This result, along with (II.17) and (II.18), guarantees that $`Q_\mathrm{\Delta }`$ and $`P_\mathrm{\Delta }`$ are independent of the choice of cross section $`S_\mathrm{\Delta }`$ of the horizon. Note that this result was obtained using only the boundary conditions; equations of motion in the bulk are not needed.
Finally, let us examine the remaining freedom in the electromagnetic field. The boundary conditions do not restrict $`𝐅_{ab}n^am^b`$ and $`{}_{}{}^{}𝐅_{ab}^{}n^am^b`$ at all. These components describe the electromagnetic radiation flowing along the horizon. Therefore, isolated horizon boundary conditions allow electromagnetic radiation arbitrarily close to —and even at— the horizon, provided none crosses it.
So far we have confined ourselves to the field strengths $`𝐅`$ and $`{}_{}{}^{}𝐅`$. However, in the action principle and the Hamiltonian framework we have to consider also the Maxwell potential $`𝐀`$. Now, if the magnetic charge is non-zero, either one has to allow ‘wire singularities’ in the vector potentials or regard $`𝐀`$ as a connection on a non-trivial $`U(1)`$-bundle. (If we regard it as a connection on a $`\text{I}\text{R}^+`$-bundle, the magnetic charge is necessarily zero.) Since we wish to deal only with smooth fields, we will not allow ‘wire singularities’ in the potentials. If we work with bundles, the magnetic charge is quantized whence the space of histories has several disconnected components. Thus, in the first law, we will not be able to consider variations $`\delta `$ of fields with $`\delta P0`$. As far as mechanics of isolated horizons is concerned, there is essentially no loss of generality if we restrict ourselves to the case $`P_\mathrm{\Delta }=0`$. Therefore, in the next three sections, while working with Maxwell fields, we will do so. As usual, our final results can be formally extended to the case of non-vanishing magnetic charge by performing a duality rotation on $`𝐅`$. (As discussed in Section VI, the situation is rather different in the case of Yang-Mills fields.)
Recall that the first law in the Einstein-Maxwell case involves the electro-static potential $`\mathrm{\Phi }`$. In static space-times, one typically sets $`\mathrm{\Phi }=\xi ^a𝐀_a`$ where $`\xi `$ is the static Killing field and the gauge is chosen such that the vector potential $`𝐀`$ tends to zero at infinity and satisfies $`_\xi 𝐀=0`$ everywhere in space-time. We now need a strategy to define the electric potential $`\mathrm{\Phi }`$ without reference to a Killing field or infinity. To this end, we introduce the following definition:
Definition 3: The electromagnetic potential $`𝐀`$ will be said to be in a gauge adapted to the weakly isolated horizon $`(\mathrm{\Delta },[\mathrm{}])`$ if it satisfies
$$_{\mathrm{}}\begin{array}{c}𝐀\hfill \end{array}\widehat{}=0.$$
(II.22)
Mathematically, this restriction is analogous to this condition $`_{\mathrm{}}\omega \widehat{}=0`$ imposed on the gravitational field in Definition 2. However, while the condition on $`\omega `$ is a physical restriction on the form of the gravitational field at $`\mathrm{\Delta }`$, the condition on $`𝐀`$ is a gauge choice; it can always be imposed without physically constraining the electromagnetic field strength.
Given an electromagnetic potential $`𝐀`$ in a gauge adapted to $`(\mathrm{\Delta },[\mathrm{}])`$, we can now define the scalar potential $`\mathrm{\Phi }_{(\mathrm{})}`$ at the horizon in an obvious fashion:
$`\mathrm{\Phi }_{(\mathrm{})}:\widehat{}=\mathrm{}𝐀.`$
The key question now is whether our boundary conditions are strong enough to ensure that $`\mathrm{\Phi }_{(\mathrm{})}`$ is constant on $`\mathrm{\Delta }`$. Only then can we hope to use this notion of the scalar potential in the first law. Note that this question is rather similar to the one we asked of surface gravity $`\kappa _{(\mathrm{})}`$ in Section II B. By using arguments completely analogous to those that led us to the zeroth law, we will now show that the answer to the present question is also in the affirmative. In a gauge adapted to the horizon,
$`0\widehat{}=_{\mathrm{}}\begin{array}{c}𝐀\hfill \end{array}\widehat{}=\begin{array}{c}\mathrm{}𝐅\hfill \end{array}\begin{array}{c}d\mathrm{\Phi }\hfill \end{array}_{(\mathrm{})}.`$
As we saw above, the boundary conditions imply $`\begin{array}{c}\mathrm{}𝐅\hfill \end{array}\widehat{}=0`$ (equation (II.17)). Hence, it follows immediately that $`\mathrm{\Phi }_{(\mathrm{})}`$ is constant on the horizon. We can regard this result as the ‘electromagnetic part’ of the zeroth law of isolated horizon mechanics.
We will see in the next section that these zeroth laws play a key role in making the gravitational and the electromagnetic action principles well-defined in presence of isolated horizons. As with surface gravity $`\kappa _{(\mathrm{})}`$, the functional dependence of $`\mathrm{\Phi }_{(\mathrm{})}`$ on the horizon parameters varies with the choice of $`\mathrm{}[\mathrm{}]`$. We will see in Section V that the Hamiltonian framework constrains these dependencies in an interesting fashion.
### D Other definitions and remarks
In this sub-section, we introduce two new definitions which are important to the general framework of isolated horizons.
The first is concerned with rotation. From one’s experience with the Newman-Penrose framework, one expects the gravitational contribution to angular momentum to be coded in the imaginary part of $`\mathrm{\Psi }_2`$. This expectation will be shown to be correct in . Therefore, in the Einstein-Maxwell theory, we introduce the following definition:
Definition 4: A weakly isolated horizon $`(\mathrm{\Delta },[\mathrm{}])`$ will be said to be non-rotating if $`\mathrm{Im}\left[\mathrm{\Psi }_2\right]`$ vanishes on $`\mathrm{\Delta }`$.
If $`(\mathrm{\Delta }_\mathrm{K},[\xi ])`$ is a Killing horizon and $`\xi `$ is a hypersurface orthogonal, time-like vector field near $`\mathrm{\Delta }_\mathrm{K}`$, on physical grounds one would expect the horizon to be non-rotating. Is this expectation compatible with our definition? The answer is in the affirmative. For, in this case, one can show that $`B_{ab}:=C_{acbd}\xi ^c\xi ^d`$ vanishes in the region where $`\xi ^a`$ is time-like. Hence, by continuity, it also vanishes on $`\mathrm{\Delta }`$ forcing $`\mathrm{Im}\left[\mathrm{\Psi }_2\right]`$ to vanish there. Similarly, if the space-time admits a hypersurface orthogonal, rotational Killing field $`\phi ^a`$ in a neighborhood of $`\mathrm{\Delta }`$ $`\mathrm{Im}\left[\mathrm{\Psi }_2\right]`$ again vanishes on $`\mathrm{\Delta }`$. The definition is again compatible with one’s intuition that the horizon should be non-rotating in this case. In this paper, while we allow for presence of rotation in the first four sections, we will restrict ourselves to non-rotating horizons in the proof of the first law in Section V.
Finally, for completeness, let us introduce a stronger notion of ‘isolation’ by strengthening the boundary conditions of Definition 2.
Definition 5: A weakly isolated horizon $`(\mathrm{\Delta },[\mathrm{}])`$ is said to be isolated if
$$[_{\mathrm{}},𝒟]V\widehat{}=0$$
(II.23)
for all vector fields $`V`$ tangential to $`\mathrm{\Delta }`$ and all $`\mathrm{}[\mathrm{}]`$
As before, if any one $`\mathrm{}`$ satisfies this condition, so do all $`\mathrm{}[\mathrm{}]`$. However, unlike (II.14), condition (II.23) is a genuine restriction in the sense that it can not always be met by a judicious choice of null normals. Generically it does suffice to single out the equivalence class $`[\mathrm{}]`$ uniquely . In particular, in the Kerr family, the only $`[\mathrm{}]`$ which satisfies (II.23) is the one containing constant multiples of the globally defined Killing field which is orthogonal to the horizon. Every Killing horizon is of course an isolated horizon. Thus, even though (II.23) is a stronger condition than (II.14), it is still very weak compared to conditions normally imposed. For most physical applications, e.g., to numerical relativity, it is appropriate to work with isolated horizons. For mechanics of isolated horizons, however, we can —and will— work with the larger class of weakly isolated horizons.
The following consequences of the boundary conditions defining weakly isolated horizons, derived in this section, will play an important role in the subsequent discussion:
$`_\begin{array}{c}a\end{array}\mathrm{}_b\widehat{}=\omega _a\mathrm{}^b`$;
$`\kappa _{(\mathrm{})}\widehat{}=\mathrm{}^a\omega _a`$, the surface gravity defined by the null normal $`\mathrm{}^a`$, is constant on $`\mathrm{\Delta }`$;
There is a natural (area) 2-form $`{}_{}{}^{2}ϵ`$ on $`\mathrm{\Delta }`$ satisfying $`_{\mathrm{}}{}_{}{}^{2}ϵ\widehat{}=0`$ and $`{}_{}{}^{2}ϵ_{ab}^{}\mathrm{}^b\widehat{}=0`$.
The electromagnetic potential $`𝐀`$ is chosen to satisfy $`_{\mathrm{}}𝐀\widehat{}=0`$ and in this gauge the scalar potential $`\mathrm{\Phi }_{(\mathrm{})}:=\mathrm{}^a𝐀_a`$ is constant on $`\mathrm{\Delta }`$; and
The electromagnetic field satisfies $`\begin{array}{c}\mathrm{}^a𝐅_{ab}\hfill \end{array}\widehat{}=0`$; and $`\begin{array}{c}\mathrm{}^a^{}𝐅_{ab}\hfill \end{array}\widehat{}=0`$.
## III Action
In this paper we use the first order formulation of general relativity in terms of tetrads and connections. Since tetrads are essential to incorporate spinorial matter, it is natural to base the framework on tetrads from the beginning. The use of a first order formalism, on the other hand, is motivated primarily by mathematical simplicity. In the first order framework, the action and the Hamiltonians can be expressed entirely in terms of differential forms which significantly simplify the variational calculations. The previous paper which dealt with undistorted horizons used spinors and self-dual connections, while here we choose to use orthonormal tetrads and real, Lorentz connections. For analyzing mechanics of isolated horizons, there are two advantages to this. First, the Hamiltonian and symplectic structure are now manifestly real which simplifies evaluation of the boundary terms at the horizon. Second, the analysis can now be extended to other space-time dimensions in a straightforward manner. However, these simplifications, come with a price. Since, at present, the self dual variables appear to be indispensable for non-perturbative quantization, the results obtained here will have to be re-expressed in terms of self-dual variables in order to extend the analysis of the quantum horizon geometry and black hole entropy to include rotation.
### A Preliminaries
Let us begin with the first order action for Einstein-Maxwell theory on a 4-dimensional manifold $``$ which is topologically $`M\times \text{I}\text{R}`$, where $`M`$ is an oriented Riemannian 3-manifold without boundary (the complement of a compact set of) which is diffeomorphic to (the complement of a compact set of) $`\text{I}\text{R}^3`$. Thus, topological complications of $`M`$, if any, are confined to a compact set. In this subsection we shall only give the relevant formulae. For details, see e.g. . Our basic fields will consist of a triplet $`(e_a^I,A_{aI}{}_{}{}^{J},𝐀_a)`$ defined on $``$ where $`e_a^I`$ denotes a co-tetrad, $`A_{aI}^J`$ the gravitational (Lorentz) connection and $`𝐀_a`$ the electro-magnetic connection. Here, lower case latin letters refer to the tangent space of $``$ while the upper case letters $`I,J`$ etc. refer to an internal four dimensional vector space $`V`$ with a fixed metric $`\eta _{IJ}`$ of signature $`(+++)`$. The co-tetrad $`e_a^I`$ is an isomorphism between the tangent space $`T_p()`$ at any point $`p`$ and the internal space $`V`$. Using it, we define a metric on $``$ by $`g_{ab}:=e_a^Ie_b^J\eta _{IJ}`$ which also has signature $`(+++)`$. The Lorentz connection $`A_{aI}^J`$ acts only on internal indices and defines a derivative operator
$`D_ak_I:=_ak_I+A_{aI}^{}{}_{}{}^{J}k_J,`$
where $``$ is a fiducial derivative operator which, as usual, will be chosen to be flat and torsion free. Finally, $`𝐀_a`$ is the $`U(1)`$ electromagnetic connection 1-form on $``$. (As noted in Section II C, we will assume that the magnetic charge is zero.) All fields will be assumed to be smooth and satisfy the standard asymptotic conditions at infinity.
The 2-forms $`\mathrm{\Sigma }^{IJ}`$
$`\mathrm{\Sigma }_{IJ}:={\displaystyle \frac{1}{2}}ϵ_{IJKL}e^Ke^L`$
constructed from the co-tetrads will play an important role throughout our calculations. In particular, the action for an asymptotically flat space-time (with no internal boundary) is given by (see e.g. )
$$S(e,A,𝐀)=\frac{1}{16\pi G}_{}\mathrm{\Sigma }^{IJ}F_{IJ}+\frac{1}{16\pi G}_\tau _{\mathrm{}}\mathrm{\Sigma }^{IJ}A_{IJ}\frac{1}{8\pi }_{}𝐅{}_{}{}^{}𝐅.$$
(III.1)
Here $`F`$ and $`𝐅`$ are the curvatures of the gravitational and electromagnetic connections $`A`$ and $`𝐀`$ respectively:
$`F_I{}_{}{}^{J}=dA_I{}_{}{}^{J}+A_I{}_{}{}^{K}A_K{}_{}{}^{J}\text{and}𝐅=d𝐀,`$
$`{}_{}{}^{}𝐅_{ab}^{}=\frac{1}{2}ϵ_{ab}{}_{}{}^{cd}𝐅_{cd}^{}`$ is the dual of $`𝐅`$ defined using $`e_a^I`$, and $`\tau _{\mathrm{}}`$ is the time-like cylinder at infinity. The boundary term at $`\tau _{\mathrm{}}`$ ensures the differentiability of the action.
Let us briefly examine the equations of motion arising from the action. Varying the action with respect to the connection, one obtains
$`D\mathrm{\Sigma }=0.`$
This condition implies that the connection $`D`$ defined by $`A`$ has the same action on internal indices as the unique connection $``$ compatible with the co-tetrad, i.e., satisfying $`_ae_b^I=0`$. When this equation of motion is satisfied, the curvature $`F`$ is related to the Riemann curvature $`R`$ of $``$ by
$`F_{ab}^{}{}_{}{}^{IJ}=R_{ab}^{}{}_{}{}^{cd}e_c^Ie_d^J.`$
Varying the action with respect to $`e_a^I`$ and taking into account the above relation between curvatures, one obtains Einstein’s equations
$`G_{ab}=8\pi GT_{ab}2G\left(𝐅_{ac}𝐅_{bd}g^{cd}{\displaystyle \frac{1}{4}}g_{ab}𝐅_{cd}𝐅^{cd}\right)`$
where $`G_{ab}`$ is the Einstein tensor and $`T_{ab}`$ the electromagnetic stress energy tensor. Finally, variation with respect to the electromagnetic connection, $`𝐀`$, yields Maxwell’s equation
$`d{}_{}{}^{}𝐅=0.`$
### B Internal boundary $`\mathrm{\Delta }`$
Let us now consider the variational principle for asymptotically flat histories which admit a weakly isolated horizon $`\mathrm{\Delta }`$ as their internal boundary. The manifold $``$ under consideration has an internal boundary $`\mathrm{\Delta }`$, topologically $`S^2\times \text{I}\text{R}`$. As before, $``$ is topologically $`M\times \text{I}\text{R}`$, where $`M`$ is now an oriented manifold with an internal, 2-sphere boundary, whose topological complications are again confined to a compact region. Space-time is bounded to the future and past by two (partial Cauchy) surfaces $`M^\pm `$, extending to spatial infinity (see Figure 1).
Following Definition 2 of weakly isolated horizons, we will equip $`\mathrm{\Delta }`$ with a fixed equivalence class of vector fields $`[\mathrm{}]`$ which are transversal to its 2-sphere cross-sections (where, as before, $`\mathrm{}\mathrm{}^{}`$ if and only if $`\mathrm{}\widehat{}=c\mathrm{}^{}`$ for a constant $`c`$). It is also convenient to fix an internal null tetrad $`(\mathrm{}^I,n^I,m^I,\overline{m}^I)`$ on $`\mathrm{\Delta }`$, each element of which is annihilated by the fiducial, flat internal connection $``$.
The permissible histories consist of smooth triplets $`(e,A,𝐀)`$ on $``$ satisfying boundary conditions at infinity and on $`\mathrm{\Delta }`$. The boundary conditions at infinity are, as before, the standard ones which ensure asymptotic flatness. Since the asymptotic behavior and boundary integrals at infinity play only a secondary role in our analysis, we shall not spell out the precise fall-off requirements. At $`\mathrm{\Delta }`$, the histories are subject to three conditions: i) the tetrads $`e`$ should be such that the vector field $`\mathrm{}^a:=\mathrm{}^Ie_I^a`$ defined by each history belongs to the equivalence class $`[\mathrm{}]`$ fixed on $`\mathrm{\Delta }`$; ii) the tetrad $`e`$ and the gravitational connection $`A`$ should be such that $`(\mathrm{\Delta },[\mathrm{}])`$ is a weakly isolated horizon for the history; and, iii) the electromagnetic potential $`𝐀`$ is in a gauge adapted to the horizon, i.e., $`_{\mathrm{}}𝐀\widehat{}=0`$.
Remark: In space-time, we have the freedom to perform a local, internal Lorentz rotation on the tetrad $`e_I^a`$ (and the gravitational connection $`A_{aI}^J`$). All these tetrads define the same Lorentzian metric $`g_{ab}`$. Since $`\mathrm{}^a`$ is required to be a null normal to $`\mathrm{\Delta }`$, the permissible gauge rotations are reduced on $`\mathrm{\Delta }`$ to the sub-group $`(\text{I}\text{R}^+\times E^2)_{\mathrm{loc}}`$ of local null rotations preserving the null direction field $`\mathrm{}`$. (Here $`\text{I}\text{R}^+`$ is the group of rescalings of $`\mathrm{}^a,n^a`$ which leaves $`m^a`$ fixed and $`E^2`$ is the 3-dimensional Euclidean group consisting of rotations in the $`\mathrm{}`$-$`m`$, $`\mathrm{}`$-$`\overline{m}`$ and $`m`$-$`\overline{m}`$ planes.) Condition i) above — dictated by the existence of a preferred equivalence class $`[\mathrm{}]`$ in Definition 2 — further reduces the internal gauge freedom to $`\text{I}\text{R}^+\times (E_{\mathrm{loc}}^2)`$, i.e., reduces the group $`\text{I}\text{R}_{\mathrm{loc}}^+`$ of local $`\mathrm{}`$-$`n`$ rescalings to the group $`\text{I}\text{R}^+`$ of global rescalings. Thus, while any one space-time $`(,g_{ab})`$, still defines infinitely many histories due to the freedom of tetrad-rotations, this freedom is somewhat reduced at $`\mathrm{\Delta }`$ because of the structure fixed by the boundary conditions.Nonetheless, from a space-time perspective, the multiplicity of histories can still be rather surprising. For example, if $`g_{ab}`$ is the Schwarzschild metric with mass $`M>0`$, there is a history in which the surface gravity $`\kappa _{(\mathrm{})}`$ is positive and another in which it is zero. This redundancy can be eliminated by working with isolated, rather than weakly isolated horizons.
Given any tetrad $`e_I^a`$, the internal null vectors $`(\mathrm{}^I,n^I,m^I,\overline{m}^I)`$ fixed on $`\mathrm{\Delta }`$ trivially provide a null tetrad $`(\mathrm{}^a,n^a,m^a,\overline{m}^a)`$. In terms of these vectors, we can express $`\begin{array}{c}\mathrm{\Sigma }\hfill \end{array}^{IJ}`$ as:
$$\begin{array}{c}\mathrm{\Sigma }\hfill \end{array}^{IJ}\widehat{}=2\mathrm{}^{[I}n^{J]}{}_{}{}^{2}ϵ+2n(im\mathrm{}^{[I}\overline{m}^{J]}i\overline{m}\mathrm{}^{[I}m^{J]}),$$
(III.2)
where, as before, $`{}_{}{}^{2}ϵ=im\overline{m}`$ is the pull-back to $`\mathrm{\Delta }`$ of the natural alternating tensor on the 2-sphere $`𝒮`$ of integral curves of $`\mathrm{}^a`$ associated with the given history. The weak isolation of $`(\mathrm{\Delta },[\mathrm{}])`$ restricts the form of the connection $`A`$ at $`\mathrm{\Delta }`$. To see this, recall that one of the equations of motion requires the connection $`D`$ defined by $`A`$ to have the same action on internal indices as $``$. Hence, $`_\begin{array}{c}a\end{array}\mathrm{}_I\widehat{}=_a\mathrm{}_I+A_{\begin{array}{c}a\end{array}I}{}_{}{}^{J}\mathrm{}_{J}^{}\widehat{}=A_{\begin{array}{c}a\end{array}I}{}_{}{}^{J}\mathrm{}_{J}^{}`$ where, in the second step we have used the fact that the flat derivative operator, $``$, has been chosen to annihilate the internal tetrad on $`\mathrm{\Delta }`$. Since $`_ae_b^I=0`$ by definition of $``$, and $`_\begin{array}{c}a\end{array}\mathrm{}^b\widehat{}=\omega _a\mathrm{}^b`$ (see II.4) it follows that $`\begin{array}{c}A\hfill \end{array}_{IJ}\mathrm{}^J\widehat{}=\omega \mathrm{}_I`$. Hence, on $`\mathrm{\Delta }`$, $`A`$ has the form:
$$\begin{array}{c}A\hfill \end{array}_{IJ}\widehat{}=2l_{[I}n_{J]}\omega +C_{IJ},$$
(III.3)
where the 1-form $`C_{IJ}`$ satisfies $`C_{IJ}\mathrm{}^J\widehat{}=0`$.
With this background material at hand, we are now ready to consider variations of the action (III.1) in the presence of an inner boundary representing a weakly isolated horizon. A key question is whether a new surface term at the horizon is necessary to make the variational principle well-defined. We will show that, thanks to the zeroth law, such a term is not needed.
A simple calculation yields:
$$\delta S(e,A,𝐀)=_{}\text{Equations of Motion}\delta \varphi \frac{1}{16\pi G}_\mathrm{\Delta }\mathrm{\Sigma }^{IJ}\delta A_{IJ}\frac{1}{4\pi }_\mathrm{\Delta }\delta 𝐀{}_{}{}^{}𝐅.$$
(III.4)
where, in the first term, $`\varphi `$ stands for the basic fields $`(e,A,𝐀)`$ in the action. Note that, as in the case without an internal boundary, the variation of the boundary term at infinity precisely cancels the contribution arising from the variation of the bulk terms.
In order to show that the action principle is viable, it is necessary to show that the terms at the horizon vanish due to the boundary conditions imposed there. Let us begin with the gravitational term. Using (III.2) and (III.3) it can be re-expressed as
$$\frac{1}{8\pi G}_\mathrm{\Delta }\delta \omega {}_{}{}^{2}ϵ.$$
(III.5)
Since $`{}_{}{}^{2}ϵ`$ is the pull-back to $`\mathrm{\Delta }`$ of the alternating tensor on the 2-sphere $`𝒮`$ of integral curves of $`\mathrm{}`$, it follows that $`_{\mathrm{}}{}_{}{}^{2}ϵ\widehat{}=0`$. Furthermore, the weak isolation of the horizon ensures $`_{\mathrm{}}\omega \widehat{}=0`$ and, since the null normal $`\mathrm{}^a`$ defined by any tetrad belongs to the fixed equivalence class $`[\mathrm{}]`$ at the horizon, we have $`\delta \mathrm{}\widehat{}=c_\delta \mathrm{}`$ for some constant $`c_\delta `$. These two facts imply $`_{\mathrm{}}\delta \omega \widehat{}=0`$. Thus the entire integrand is Lie dragged by $`\mathrm{}`$. In the variational principle, however, all fields are fixed on the initial and final hypersurfaces, say $`M^\pm `$. In particular, $`\delta \omega `$ necessarily vanishes on the initial and final cross sections of the horizon. Therefore, the integrand in (III.5) vanishes on the initial and final cross sections and is Lie dragged by $`\mathrm{}`$. This immediately implies (III.5) is zero.
Let us now consider the electromagnetic term. Since every $`𝐀`$ is in a gauge adapted to the isolated horizon, $`_{\mathrm{}}𝐀\widehat{}=0`$. Furthermore, $`\delta \mathrm{}^a\widehat{}=c_\delta \mathrm{}^a`$, so we conclude $`_{\mathrm{}}\delta 𝐀\widehat{}=0`$. Next, (II.21) ensures $`_{\mathrm{}}{}_{}{}^{}𝐅\widehat{}=0`$. Thus, the integrand of the electromagnetic surface term is Lie dragged by $`\mathrm{}^a`$. An identical argument to the one presented above implies that the electromagnetic surface term in (III.4) also vanishes. Therefore, the variation of the action (III.1) continues to yield Einstein-Maxwell equations in spite of the presence of an inner boundary representing a weakly isolated horizon.
It is instructive to re-examine the key step in the above argument. Suppose we only had a non-expanding horizon. Then, the gravitational surface term could still be reduced to (III.5), and $`{}_{}{}^{2}ϵ`$ and $`{}_{}{}^{}𝐅`$ would still be Lie-dragged by $`\mathrm{}`$. However, in this case, we could not argue that $`\omega `$ and $`𝐀`$ are also Lie-dragged. As we saw in Sections II B and II C, these conditions are equivalent, respectively, to the constancy of the surface gravity $`\kappa _{(\mathrm{})}`$ and the electromagnetic potential $`\mathrm{\Phi }_{(\mathrm{})}`$ on $`\mathrm{\Delta }`$. In this sense, given a non-expanding horizon as the inner boundary, the gravitational and electromagnetic zeroth laws are the necessary and sufficient conditions one must impose for the viability of the standard, first order, tetrad action principle.
Remark: Note that (III.1) is not the unique viable action for the problem: as usual, there is freedom to add suitable boundary terms without affecting the viability. Specifically, we are free to add any horizon boundary term which is composed entirely of fields which are Lie dragged by $`\mathrm{}`$, for example the intrinsic horizon metric $`q_{ab}`$ and fields $`\omega `$, $`{}_{}{}^{2}ϵ`$, and $`𝐀`$ . Then, due to the argument given above, the new action would also be viable. However, as usual, this freedom will not affect the definition of the symplectic structure which underlies the Hamiltonian treatment of the next section.
## IV Covariant Phase space
Let us now construct the phase space of space-times containing weakly isolated horizons. In the next section, we will use this framework to construct Hamiltonians generating suitable time translations and define the energy of an isolated horizon. In , the phase space was constructed by performing a Legendre transform of the action. This procedure leads to a ‘canonical’ framework in which the phase space consists of configuration and momentum variables defined on a spatial hypersurface. With the self-dual connections used in , the gravitational configuration variable turns out to be a connection and its conjugate momentum, a 2-form so that the Hamiltonian description can again be given in terms of forms. With the full Lorentz connections now under consideration, the situation turns out to be more complicated. Specifically, one encounters certain second class constraints and, when these are solved, one ends up with the same canonical phase space that one would have obtained through a second order formalism. In the Hamiltonian framework, then, the simplicity we encountered in Section III is lost. More specifically, constraint functions and Hamiltonians now contain terms involving second derivatives of the basic canonical variables which make variations rather complicated. (For details, see chapters 3 and 4 in .) Therefore, in this section we will not use a Legendre transform. Instead, we will construct the ‘covariant phase space’ from the space of solutions to field equations (see, e.g., ). As in Section III, all expressions will now involve only the basic form-fields and their exterior derivatives and variational calculations will continue to be simple.
To specify the phase space, let us begin as in Section III by fixing a manifold $``$ with an internal boundary $`\mathrm{\Delta }`$ (see figure 1). As before, we will equip $`\mathrm{\Delta }`$ with an equivalence class $`[\mathrm{}]`$ of vector fields transverse to its 2-sphere cross-sections. To evaluate the symplectic structure and Hamiltonians, we will often use a partial Cauchy surface $`M`$ in the interior of $``$ which intersects $`\mathrm{\Delta }`$ in a 2-sphere $`S`$. Points of the covariant phase space $`\mathrm{\Gamma }`$ will consist of histories considered in Section III which satisfy field equations. More explicitly, $`\mathrm{\Gamma }`$ consists of asymptotically flat solutions $`(e,A,𝐀)`$ to the field equations on $``$ such that i) the vector field $`\mathrm{}^a:=\mathrm{}^Ie_I^a`$ belongs to the equivalence class $`[\mathrm{}]`$ fixed on $`\mathrm{\Delta }`$, ii) in each solution, $`(\mathrm{\Delta },[\mathrm{}])`$ is a weakly isolated horizon; and, iii) the electromagnetic potential $`𝐀`$ is in a gauge adapted to the horizon, i.e., $`_{\mathrm{}}𝐀\widehat{}=0`$.
Our next task is to use the action (III.1) to define the symplectic structure $`\mathrm{\Omega }`$ on $`\mathrm{\Gamma }`$. It is convenient to make a brief detour and first introduce two new fields which can be regarded as ‘potentials’ for the surface gravity $`\kappa _{(\mathrm{})}`$ and the electric potential $`\mathrm{\Phi }_{(\mathrm{})}`$. Given any point $`(e,A,𝐀)`$ in the phase space $`\mathrm{\Gamma }`$, let us define scalar fields $`\psi `$ and $`\chi `$ on $`\mathrm{\Delta }`$ as follows:
i) $`_{\mathrm{}}\psi \widehat{}=(\mathrm{}\omega )\widehat{}=\kappa _{(\mathrm{})}`$ and $`_{\mathrm{}}\chi \widehat{}=(\mathrm{}𝐀)\widehat{}=\mathrm{\Phi }_{(\mathrm{})}`$; and
ii) $`\psi `$ and $`\chi `$ vanish on $`S^{}`$, the intersection of $`M^{}`$ with $`\mathrm{\Delta }`$.<sup>\**</sup><sup>\**</sup>\**Condition ii) serves only to fix the freedom to add constants to $`\psi `$ and $`\chi `$. One could envisage replacing it by a different condition. Our results will be insensitive to this choice.
Note that these conditions associate with each point of $`\mathrm{\Gamma }`$ a unique pair $`(\psi ,\chi )`$ on $`\mathrm{\Delta }`$ and in the ‘extremal’ case $`\kappa _{(\mathrm{})}=0`$, $`\psi `$ vanishes identically.
We wish to use the standard procedure involving second variations of the action to define the symplectic structure.<sup>††</sup><sup>††</sup>††Actually this procedure provides a pre-symplectic structure, i.e., a closed 2-form on the phase-space which, however, is generally degenerate. The vectors in its kernel represent infinitesimal ‘gauge transformations’. The physical phase space is obtained by quotienting the space of solutions by gauge transformations and inherits a true symplectic structure from the pre-symplectic structure on the space of solutions. The 2-form $`\mathrm{\Omega }`$ introduced below is indeed degenerate. However, for simplicity, we will abuse the notation somewhat and refer to $`\mathrm{\Omega }`$ as the symplectic structure. Let us recall the main steps of this procedure. One first constructs the symplectic current $`J`$: Given a point $`\gamma `$ in the phase space $`\mathrm{\Gamma }`$ and two tangent vectors $`\delta _1`$ and $`\delta _2`$ at that point, $`J`$ provides a closed 3-form $`J(\gamma ;\delta _1,\delta _2)`$ on $``$. Integrating $`dJ`$ over the part $`\stackrel{~}{}`$ of space-time under consideration, one obtains
$`0={\displaystyle _\stackrel{~}{}}𝑑J(\gamma ;\delta _1,\delta _2)={\displaystyle _\stackrel{~}{}}J.`$
Now, if there is no internal boundary, one can choose $`\stackrel{~}{}`$ to be a region bounded by any two Cauchy surfaces $`M_1`$ and $`M_2`$ so that the boundary is given by $`\stackrel{~}{}=M_1M_2\tau _{\mathrm{}}`$, where $`\tau _{\mathrm{}}`$ is the time-like ‘cylinder at infinity’. In simple cases, the asymptotic conditions ensure that the integral $`_\tau _{\mathrm{}}J(\gamma ;\delta _1,\delta _2)`$ vanishes. Then, taking orientations into account, it follows that $`_MJ(\gamma ;\delta _1,\delta _2)`$ is independent of the choice of Cauchy surface $`M`$. One then sets the symplectic structure to be
$`\mathrm{\Omega }|_\gamma (\delta _1,\delta _2)={\displaystyle _M}J(\gamma ;\delta _1,\delta _2).`$
In our case, the second variation of the action (III.1) yields the following symplectic current:
$$J(\gamma ;\delta _1,\delta _2)=\frac{1}{16\pi G}[\delta _1\mathrm{\Sigma }^{IJ}\delta _2A_{IJ}\delta _2\mathrm{\Sigma }^{IJ}\delta _1A_{IJ}]\frac{1}{4\pi }[\delta _1{}_{}{}^{}𝐅\delta _2𝐀\delta _2{}_{}{}^{}𝐅\delta _1𝐀].$$
(IV.1)
Using the fact that the fields $`\gamma (e,A,𝐀)`$ satisfy the field equations and the tangent vectors $`\delta _1,\delta _2`$ satisfy the linearized equations off $`\gamma `$, one can directly verify that $`J(\gamma ;\delta _1,\delta _2)`$ is in fact closed as guaranteed by the general procedure involving second variations. It is now natural to choose $`\stackrel{~}{}`$ to be a part of our space-time $``$ bounded by two partial Cauchy surfaces $`M_1,M_2`$, the time-like cylinder $`\tau _{\mathrm{}}`$ and a part $`\stackrel{~}{\mathrm{\Delta }}`$ of the isolated horizon bounded by $`M_1`$ and $`M_2`$. Again, the asymptotic conditions ensure that the integral of $`J`$ over $`\tau _{\mathrm{}}`$ vanishes. Hence,
$`({\displaystyle _{M_1}}+{\displaystyle _{M_2}}+{\displaystyle _{\stackrel{~}{\mathrm{\Delta }}}})J(\gamma ;\delta _1,\delta _2)=0.`$
However, this does not immediately provide us the conserved symplectic structure because the integral of $`J`$ over $`\stackrel{~}{\mathrm{\Delta }}`$ does not vanish in general. Since the isolation of the horizon implies that there are no fluxes of physical quantities across $`\mathrm{\Delta }`$, one might expect that, although non-zero, the integral over $`\stackrel{~}{\mathrm{\Delta }}`$ would be ‘controllable’. This is indeed the case. Using the forms (III.2) and (III.3) of $`\mathrm{\Sigma }`$ and $`A`$ on $`\mathrm{\Delta }`$ and the definitions of the potentials $`\psi `$ and $`\chi `$, it is easy to verify that the pull-back of the symplectic current to $`\mathrm{\Delta }`$ is itself exact:
$`\begin{array}{c}J\hfill \end{array}(\gamma ;\delta _1,\delta _2)\widehat{}=dj(\gamma ,\delta _1,\delta _2)`$
where the 2-form $`j`$ on $`\mathrm{\Delta }`$ is given by:
$`j(\gamma ,\delta _1,\delta _2)={\displaystyle \frac{1}{8\pi G}}[\delta _1\psi \delta _2({}_{}{}^{2}ϵ)\delta _2\psi \delta _1({}_{}{}^{2}ϵ)]+{\displaystyle \frac{1}{4\pi }}[\delta _1\chi \delta _2{}_{}{}^{}𝐅\delta _2\chi \delta _1{}_{}{}^{}𝐅].`$
Hence, if $`M_1`$ and $`M_2`$ intersect $`\stackrel{~}{\mathrm{\Delta }}`$ in 2-spheres $`S_1`$ and $`S_2`$ respectively, we have:
$`{\displaystyle _{\stackrel{~}{\mathrm{\Delta }}}}J(\gamma ;\delta _1,\delta _2)=({\displaystyle _{S_1}}+{\displaystyle _{S_2}})j(\gamma ;\delta _1,\delta _2).`$
The negative sign appearing in the above expression is due to the choice of orientation of $`S_\mathrm{\Delta }`$, which is induced from $`M`$ rather than from $`\mathrm{\Delta }`$. Using these results we can define the symplectic structure as:
$`\mathrm{\Omega }|_\gamma (\delta _1,\delta _2)=`$ (IV.4)
$`{\displaystyle \frac{1}{16\pi G}}{\displaystyle _M}[\delta _1\mathrm{\Sigma }^{IJ}\delta _2A_{IJ}\delta _2\mathrm{\Sigma }^{IJ}\delta _1A_{IJ}]+{\displaystyle \frac{1}{8\pi G}}{\displaystyle _S}[\delta _1({}_{}{}^{2}ϵ)\delta _2\psi \delta _2({}_{}{}^{2}ϵ)\delta _1\psi ]`$
$`{\displaystyle \frac{1}{4\pi }}{\displaystyle _M}[\delta _1{}_{}{}^{}𝐅\delta _2𝐀\delta _2{}_{}{}^{}𝐅\delta _1𝐀]+{\displaystyle \frac{1}{4\pi }}{\displaystyle _S}[\delta _1{}_{}{}^{}𝐅\delta _2\chi \delta _2{}_{}{}^{}𝐅\delta _1\chi ]`$
Again, using field equations one can directly verify that the right side of (IV.4) is independent of the choice of the partial Cauchy surface; the symplectic structure is ‘conserved’. We will use $`(\mathrm{\Gamma },\mathrm{\Omega })`$ as our covariant phase space.
Note that, even though the action did not contain a surface term at the horizon, the symplectic structure does. So, the overall situation is the same as in the undistorted, non-rotating case considered in . Finally, our discussion of the action principle and our construction of the covariant phase space is applicable to all weakly isolated horizons $`\mathrm{\Delta }`$; nowhere did we have to restrict ourselves to the non-rotating case.
## V Hamiltonian evolution and the first law
To discuss the first law, we must first define horizon energy, which in turn requires a time evolution field $`t^a`$ on $``$. Given a vector field $`t^a`$ satisfying appropriate boundary conditions, $`\delta _t:=(_te,_tA,_t𝐀)`$ satisfies the linearized equations for any $`\gamma :=(e,A,𝐀)`$ in $`\mathrm{\Gamma }`$ and thus defines a vector field on $`\mathrm{\Gamma }`$.<sup>‡‡</sup><sup>‡‡</sup>‡‡In the Lie-derivatives, the internal indices are treated as scalars; thus $`_te_a^I=t^b_be_a^I+e_b^I_at^b`$. To make $`\delta _t`$ a well-defined vector field on $`\mathrm{\Gamma }`$, we now exclude $`M^\pm `$ from $``$ and let $``$ and $`\mathrm{\Delta }`$ be without future and past boundaries. Whenever needed, these boundaries, $`M^\pm `$ and $`S^\pm `$, can be added by taking the obvious closure of $``$. This $`\delta _t`$ can be interpreted as the infinitesimal generator of time evolution on the covariant phase space. It is then natural to ask if this vector field is a phase space symmetry, i.e., if $`_{\delta _t}\mathrm{\Omega }`$ vanishes everywhere on $`\mathrm{\Gamma }`$. The necessary and sufficient condition for this to happen is that there exist a function $`H_t`$ —the Hamiltonian generating the $`t`$-evolution— such that
$$\delta H_t=\mathrm{\Omega }(\delta ,\delta _t)$$
(V.1)
for all vector fields $`\delta `$ to $`\mathrm{\Gamma }`$. On general grounds, one expects $`H_t`$ to contain a surface term $`E_{\mathrm{ADM}}^t`$ at infinity representing the total (i.e. ADM) energy, and a surface term $`E_\mathrm{\Delta }^t`$ at the horizon which can be interpreted as the horizon energy, both tied to the evolution field $`t^a`$.
A key question then is to specify the appropriate boundary conditions on $`t^a`$. It is clear that, at infinity, $`t^a`$ should be an asymptotic time-translation, i.e., should approach a time-translation Killing field of the flat metric used to specify the boundary conditions. At the horizon, on the other hand, the metric is not universal and the space-time defined by a generic point $`\gamma `$ of the covariant phase space does not admit any Killing field near $`\mathrm{\Delta }`$. Therefore, specification of the boundary conditions at $`\mathrm{\Delta }`$ is not as straightforward as that at infinity. It is for this reason that we now assume that $`(\mathrm{\Delta },[\mathrm{}])`$ is a non-rotating, weakly isolated horizon for all points $`\gamma (e,A,𝐀)`$ of the phase space. The problem of specifying the appropriate boundary conditions on $`t^a`$ in the rotating case is more complicated . However, it has been addressed successfully and will be discussed in .
Recall that the internal boundary $`\mathrm{\Delta }`$ of $``$ is equipped with a specific equivalence class $`[\mathrm{}]`$ of vector fields. As discussed in Section II, these $`\mathrm{}`$ are the isolated horizon analogs of constant multiples of Killing fields on the Killing horizons in static space-times. Therefore, in the non-rotating case, it is natural to demand that, on $`\mathrm{\Delta }`$, the evolution vector field $`t^a`$ should belong to the equivalence class $`[\mathrm{}]`$. This automatically ensures that $`(_te,_tA,_t𝐀)`$ satisfy the appropriate boundary conditions to define a tangent vector at each point of the phase space $`\mathrm{\Gamma }`$. However, unlike at infinity, the geometry at the horizon is not fixed once and for all. Therefore, it is natural to allow the precise value of the evolution vector field $`t^a`$ on $`\mathrm{\Delta }`$ to vary from one point of the phase space to another. In more familiar terms, this corresponds to allowing the (boundary values of) lapse and shift fields to depend on dynamical fields $`(e,A,𝐀)`$ themselves, a procedure routinely used in numerical relativity and gauge-fixed calculations in canonical gravity. Following the current terminology in numerical relativity, we will refer to such $`t^a`$ as live evolution vector fields. The use of live fields turns out to be necessary to ensure that $`\delta _t`$ is a phase space symmetry, i.e., yields a Hamiltonian evolution on $`(\mathrm{\Gamma },\mathrm{\Omega })`$.
Let us fix a live $`t^a`$ whose restriction to the horizon belongs to the equivalence class $`[\mathrm{}]`$ at all points of the phase space. To analyze if $`\delta _t`$ is a Hamiltonian vector field, it is simplest to compute the 1-form $`X_t`$ on $`\mathrm{\Gamma }`$ defined by
$$X_t(\delta )=\mathrm{\Omega }(\delta ,\delta _t).$$
(V.2)
Now $`\delta _t`$ is Hamiltonian — i.e., $`_{\delta _t}\mathrm{\Omega }=0`$ on $`\mathrm{\Gamma }`$ — if and only if $`X_t`$ is closed, i.e.,
$`\mathrm{d}\mathrm{d}X_t=0`$
where $`\mathrm{d}\mathrm{d}`$ denotes the exterior derivative on (the infinite dimensional) phase space $`\mathrm{\Gamma }`$. If this is the case then, up to an additive constant, the Hamiltonian is given by
$`\mathrm{d}\mathrm{d}H_t=X_t.`$
To calculate the right side of (V.2), it is useful to note the following identities from differential geometry:
$$\begin{array}{cccccc}\hfill _tA& =& tF+D(tA)\hfill & \hfill _t\mathrm{\Sigma }& =& tD\mathrm{\Sigma }+D(t\mathrm{\Sigma })[(tA),\mathrm{\Sigma }]\hfill \\ \hfill _t𝐀& =& t𝐅+d(t𝐀)\hfill & \hfill _t{}_{}{}^{}𝐅& =& t(d{}_{}{}^{}𝐅)+d(t{}_{}{}^{}𝐅)\hfill \end{array}$$
(V.3)
Using these, the field equations satisfied by $`(e,A,𝐀)`$ and the linearized field equations for $`\delta `$, we obtain the required expression of $`X_t`$:
$`X_t(\delta )`$ $`:=`$ $`\mathrm{\Omega }(\delta ,_t)`$ (V.4)
$`=`$ $`{\displaystyle \frac{1}{16\pi G}}{\displaystyle _M}\mathrm{Tr}\left[(tA)\delta \mathrm{\Sigma }(t\mathrm{\Sigma })\delta A\right]{\displaystyle \frac{1}{4\pi }}{\displaystyle _M}(t𝐀)\delta ({}_{}{}^{}𝐅)(t{}_{}{}^{}𝐅)\delta 𝐀.`$ (V.5)
Note that the expression involves integrals only over the 2-sphere boundaries $`S_{\mathrm{}}`$ and $`S_\mathrm{\Delta }`$ of $`M`$, the partial Cauchy surface used in the evaluation of the symplectic structure; there is no volume term.
The integrals at infinity can be evaluated easily by making use of the fall-off conditions. As one would expect, the electromagnetic term vanishes (because $`𝐀`$ falls off at least as $`1/r`$ while $`𝐅`$ falls off as $`1/r^2`$) while the gravitational term yields precisely the ADM energy $`E_{\mathrm{ADM}}^t`$ associated with the asymptotic time-translation defined by $`t^a`$. At the horizon, we can use equations (III.2) and (III.3) to show that $`t\mathrm{\Sigma }`$ contracted on internal indices with $`\delta A`$ vanishes and (II.18) implies $`t{}_{}{}^{}𝐅=0`$ leaving
$`X_t(\delta )`$ $`=`$ $`{\displaystyle \frac{1}{8\pi G}}{\displaystyle _{S_\mathrm{\Delta }}}(t\omega )\delta ({}_{}{}^{2}ϵ){\displaystyle \frac{1}{4\pi }}{\displaystyle _{S_\mathrm{\Delta }}}(t𝐀)\delta ({}_{}{}^{}𝐅)+\delta E_{\mathrm{ADM}}^t`$ (V.6)
$`=`$ $`{\displaystyle \frac{1}{8\pi G}}\kappa _{(t)}\delta a_\mathrm{\Delta }\mathrm{\Phi }_{(t)}\delta Q_\mathrm{\Delta }+\delta E_{\mathrm{ADM}}^t`$ (V.7)
where, in the last step, we have used the fact that both $`t\omega =\kappa _{(t)}`$ and $`t𝐀=\mathrm{\Phi }_{(t)}`$ are constant on the horizon and the definition (II.20) of electric charge.
The necessary and sufficient condition for the existence of a Hamiltonian is that $`X_t`$ be closed. Clearly, this is equivalent to
$$\frac{1}{8\pi G}\mathrm{d}\mathrm{d}\kappa _{(t)}\mathrm{d}\mathrm{d}a_\mathrm{\Delta }+\mathrm{d}\mathrm{d}\mathrm{\Phi }_{(t)}\mathrm{d}\mathrm{d}Q_\mathrm{\Delta }=0,$$
(V.8)
where $``$ denotes the antisymmetric tensor product on $`\mathrm{\Gamma }`$. Now, (V.8) trivially implies that the surface gravity $`\kappa _{(t)}`$ and the electric potential $`\mathrm{\Phi }_{(t)}`$ at the horizon defined by $`t^a`$ can depend only upon the area and charge of the horizon. Other factors, such as the ‘shape’ of the distorted horizon, can not affect the values of $`\kappa _{(t)}`$ or $`\mathrm{\Phi }_{(t)}`$. Finally, (V.8) is the necessary and sufficient condition that there exists a function $`E_\mathrm{\Delta }^t`$ , also only of $`a_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$ such that
$$\delta E_\mathrm{\Delta }^t=\frac{1}{8\pi G}\kappa _{(t)}\delta a_\mathrm{\Delta }+\mathrm{\Phi }_{(t)}\delta Q_\mathrm{\Delta }.$$
(V.9)
Since $`E_\mathrm{\Delta }^t`$ is a function only of $`a_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$, it is a function of fields defined locally at the horizon. As noted before, it is natural to interpret $`E_\mathrm{\Delta }^t`$ as the horizon energy defined by the time translation $`t^a`$. The total Hamiltonian is given by:
$$H_t=E_{\mathrm{ADM}}^tE_\mathrm{\Delta }^t.$$
(V.10)
Let us summarize. Eq (V.9) is a necessary and sufficient condition for the 1-form $`X_t`$ to be closed. Therefore, the vector field $`\delta _t`$ on $`\mathrm{\Gamma }`$ defined by the space-time evolution field $`t^a`$ is Hamiltonian if and only if the first law (V.9) holds. Thus (V.9) is a restriction on the choice of the live vector field $`t^a`$: While any $`t^a`$ (which preserves the boundary conditions) defines an evolution flow on the phase space, it is only when
$$\frac{1}{8\pi G}\kappa _{(t)}\mathrm{d}\mathrm{d}a_\mathrm{\Delta }+\mathrm{\Phi }_{(t)}\mathrm{d}\mathrm{d}Q_\mathrm{\Delta }$$
is an exact 1-form on $`\mathrm{\Gamma }`$ that this flow is Hamiltonian (i.e., preserves the symplectic structure). At first, this restriction seems somewhat surprising because, in absence of internal boundaries, every vector field $`t^a`$ (which tends to a fixed Killing field of the flat metric at infinity) defines a Hamiltonian evolution. However, even in this context, there is no a priori reason to expect this tight correspondence to hold if one allows general, live vector fields $`t^a`$ whose boundary values at infinity can change from one space-time to another. Finally, we will see in Section VII that every space-time belonging to the phase space $`\mathrm{\Gamma }`$ admits an infinite family of vector fields $`t^a`$ for which $`X_t`$ is closed. Therefore, in particular, the first law does not restrict the ‘background’ space-times (or the variations $`\delta `$) in any way. Indeed, for any space-time in our phase space, there is an infinite family of first laws, one associated with each permissible $`t^a`$.
We will conclude this section with a few remarks.
i) Form of $`H_t`$: The Hamiltonian (V.10) contains only surface terms. This may seem surprising because, in the canonical framework, the familiar Hamiltonian contains a volume integral consisting of a linear combination of constraints. While the volume term vanishes ‘on shell’ and does not contribute to numerical value of the canonical Hamiltonian on physical states, it is nonetheless crucial for obtaining the correct evolution equations since derivatives of the Hamiltonian transverse to the constraint surface are needed to construct the Hamiltonian vector field. The covariant phase space, by contrast, consists only of solutions to the field equations whence the issue of taking ‘off shell’ derivatives never arises. In diffeomorphism invariant theories, the Hamiltonian on the covariant phase space is always made of surface terms.<sup>\**</sup><sup>\**</sup>\**In particular, therefore, Hamiltonians generating diffeomorphism which have support away from the boundaries vanish identically. Unlike their counterparts on the canonical phase space, the infinitesimal phase-space motions induced by such space-time vector fields are in the kernel of the covariant symplectic structure. If space-time has several asymptotic regions, the boundary term in each region defines the standard energy corresponding to that region. Therefore, in the present case, it was natural to interpret $`E_\mathrm{\Delta }^t`$ as the horizon-energy defined by the $`t^a`$ evolution. Finally, we should emphasize that we used a covariant phase space only for simplicity. The final results go through (and, in fact, were first obtained) in a canonical framework as well.
ii)Comparisons: As noted in the Introduction, all treatments of the first law for non-rotating but possibly distorted horizons available in the literature refer to static space-times. The isolated horizon framework, by contrast, does not refer to a Killing field at all and thus allows a significantly larger class of physically interesting situations. On the other hand, since it relies on a Hamiltonian framework, we cannot incorporate phenomenological matter if it does not admit a phase space description. Other treatments based on Hamiltonian methods generally restrict themselves to static space-times with a non-zero surface gravity. This assumption is essential there because those treatments use ‘bifurcate’ surfaces in an important way and these do not exist in the extremal static solutions where the surface gravity vanishes. In contrast, the results of this section, do not refer to a bifurcate surface and go through irrespective of whether $`\kappa _{(t)}`$ is non-zero or zero. In a realistic collapse, the physical space-time is not expected to have the bifurcate surface. The present analysis uses only the portion of the physical space-time in which the horizon has settled down with no further in-going radiation, rather than an analytical continuation of the near horizon geometry used in certain approaches. Finally, in contrast to other treatments, we have an infinite family of first laws, one for each evolution field $`t^a`$ for which $`\delta _t`$ is a phase space symmetry.
iii)Non-uniqueness of energy: Each permissible, live $`t^a`$ defines a horizon energy $`E_\mathrm{\Delta }^t`$. At first it seems surprising that there is so much freedom in the notion of energy. Let us compare the situation at $`^+`$, which, like $`\mathrm{\Delta }`$, is null. There, we only have a 3-parameter freedom which, furthermore, can be eliminated simply by fixing a rest frame. How does this difference arise? Recall that energy is (the numerical value of) the generator of an unit time translation. At infinity, all 4-metrics in the phase space approach the same flat metric. Hence, we can simply fix a unit time-translation Killing field $`t_o^a`$ of that flat metric near infinity and use its restriction to $`^+`$ as the unit time-translation for all metrics in the phase space. By contrast, there is no fixed 4-metric near $`\mathrm{\Delta }`$ to which all the metrics in our phase space approach. Hence, we do not have the analog of $`t_o^a`$; only the equivalence class $`[\mathrm{}]`$ is now common to all the metrics. If, for a given metric $`\stackrel{~}{g}_{ab}`$ in our collection, we select the time translation represented by a specific $`\stackrel{~}{\mathrm{}}^a`$ in $`[\mathrm{}]`$, a priori we do not know which vector field $`l^a`$ in $`[\mathrm{}]`$ would represent the ‘same’ time-translation for another geometry $`g_{ab}`$. One might imagine using the seemingly simplest strategy: just fix a $`\mathrm{}_o^a`$ in $`[\mathrm{}]`$ and demand that $`t^a`$ approach that $`\mathrm{}_0^a`$ for all points $`\gamma `$ in the phase space. Unfortunately, the strategy is not viable because such a $`t^a`$ fails to define a Hamiltonian evolution in $`\mathrm{\Gamma }`$.<sup>\*†</sup><sup>\*†</sup>\*†As we saw above, a necessary condition for $`\delta _t`$ to be a Hamiltonian vector field on $`\mathrm{\Gamma }`$ is that $`\kappa _{(t)}`$ is a function only of $`a_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$. Therefore, if we can find a tangent vector $`\delta `$ in the phase space with $`\delta a_\mathrm{\Delta }=\delta Q_\mathrm{\Delta }=0`$ but $`\delta \kappa _{(t)}0`$, $`t^a`$ can not define a Hamiltonian evolution. It is easy to find such a tangent vector $`\delta `$ for this $`t^a`$. Finally, if we restrict ourselves to globally static space-times, we can overcome this difficulty by always working with the Killing field which is unit at infinity. However, in absence of global Killing fields, the behavior of the evolution vector field $`t^a`$ near the horizon is unrelated to its behavior near infinity. Nonetheless, as we shall show in Section VII, if one has sufficient control on the space of static solutions of the theory under consideration, it is possible to select a preferred energy function on the phase space and use it as the mass of the isolated horizon. In the Einstein-Maxwell case, all static solutions with horizons are explicitly known whence the strategy is viable.
## VI Yang-Mills Field
In the previous three sections we restricted our attention to Einstein-Maxwell theory. We will now indicate how Yang-Mills fields can be included. This section is divided into three parts. In the first, we discuss restrictions on the Yang-Mills fields due to the horizon boundary conditions and introduce the notion of a ‘Yang-Mills gauge adapted to the horizon’. In the second part, we discuss the action principle and construct the covariant phase space for Einstein-Yang-Mills theory. Using this formalism, in the third subsection, we introduce a Hamiltonian generating time evolution and extend the first law to the Yang-Mills case.
### A Preliminaries
We will restrict ourselves to compact gauge groups $`G`$ and Yang-Mills connections defined on trivial bundles. Since the bundle is trivial, the connection gives rise to a smooth, globally defined Lie algebra valued 1-form, $`𝐀`$. As usual, the Yang-Mills derivative operator $`𝐃`$ will be defined as $`𝐃\lambda =\lambda +[𝐀,\lambda ]`$, where $``$ is a flat Yang-Mills connection, and the field strength, $`𝐅`$, via
$$𝐅:=d𝐀+𝐀𝐀.$$
(VI.1)
The stress energy tensor, $`𝐓`$, is given in terms of the field strength as
$$𝐓_{ab}=\frac{1}{4\pi }[𝐅_{ac}^{}{}_{}{}^{i}𝐅_{b}^{}{}_{}{}^{c}{}_{i}{}^{}\frac{1}{4}g_{ab}𝐅_{cd}^{}{}_{}{}^{i}𝐅_{}^{cd}{}_{i}{}^{}],$$
(VI.2)
where the label $`i`$ runs over the internal indices in the Lie algebra of the group G.
Let us begin by examining how the isolated horizon boundary conditions restrict the form of the field strength, $`𝐅`$, on $`\mathrm{\Delta }`$. Since the Yang-Mills stress energy tensor has the same form as the Maxwell one, (II.15), the analysis is completely analogous to that of Section II C. Therefore, we shall not include derivations of the results, but instead highlight the differences.
Recall that, on a non-expanding horizon, $`R_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=0`$ whence $`T_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=0`$. This has several consequences for the Yang-Mills field. In particular, one concludes
$$\begin{array}{c}\mathrm{}^a𝐅_{ab}\hfill \end{array}\widehat{}=0\text{and}\begin{array}{c}\mathrm{}^a^{}𝐅_{ab}\hfill \end{array}\widehat{}=0.$$
(VI.3)
These two restrictions guarantee there is no flux of Yang-Mills field across the horizon. Making use of the specific form of the stress energy tensor, we also conclude
$$\mathrm{\Phi }_{02}=\frac{1}{2}R_{ab}m^am^b\widehat{}=0\text{and}\mathrm{\Phi }_{20}=\frac{1}{2}R_{ab}\overline{m}^a\overline{m}^b\widehat{}=0.$$
(VI.4)
Our next task is to define the Yang-Mills equivalents of the electric and magnetic charges of the horizon. Naively, one might consider integrating $`𝐅`$ and $`{}_{}{}^{}𝐅`$ over a 2-sphere cross section of the horizon as was done in the Maxwell theory. However, these 2-forms now have a free internal index and are only gauge covariant rather than gauge invariant. Since there is no preferred internal basis at the horizon, the integrals would fail to be well-defined. Therefore, we must look for 2-forms which are gauge invariant. A natural quantity to consider is the norm of $`𝐅`$, defined by the Killing-Cartan form $`K_{ij}`$ on the Lie-algebra of $`G`$ and the (contravariant) $`{}_{}{}^{2}ϵ`$ on the horizon:
$$𝐅:=\left[({}_{}{}^{2}ϵ𝐅)^i({}_{}{}^{2}ϵ𝐅)^jK_{ij}\right]^{\frac{1}{2}}$$
(VI.5)
(Although the contravariant $`{}_{}{}^{2}ϵ`$ is ambiguous up to terms of the type $`\mathrm{}^{[a}V^{b]}`$ where $`V^a`$ is any vector field tangential to $`\mathrm{\Delta }`$, this ambiguity does not affect $`𝐅`$ because of (VI.3). The norm of $`{}_{}{}^{}𝐅`$ is defined analogously. These two quantities are gauge invariant and allow us to define the electric and magnetic Yang-Mills charges of the horizon:
$$Q_\mathrm{\Delta }^{YM}:\widehat{}=\frac{1}{4\pi }_{S_\mathrm{\Delta }}{}_{}{}^{}𝐅{}_{}{}^{2}ϵ\text{and}P_\mathrm{\Delta }^{YM}:\widehat{}=\frac{1}{4\pi }_{S_\mathrm{\Delta }}𝐅{}_{}{}^{2}ϵ.$$
(VI.6)
Recall that the unusual signs in the definitions of the charges arise due to the orientation of the $`S_\mathrm{\Delta }`$ — the normal to the two sphere is inward pointing. In Maxwell theory, the magnetic charge is zero unless we consider either connections on non-trivial bundles or allow ‘wire singularities’. As is well known, this is not true for Yang-Mills theory: the magnetic charge can be non-zero even if we restrict attention to smooth fields on a trivial bundle.
We would now like to verify that the charges defined in (VI.6) are independent of the cross section of the horizon $`S_\mathrm{\Delta }`$ on which the integration is performed. The isolated horizon boundary conditions guarantee this is the case. First, recall the geometric identity
$$_{\mathrm{}}\begin{array}{c}𝐅\hfill \end{array}\widehat{}=\mathrm{}\begin{array}{c}\mathrm{𝐃𝐅}\hfill \end{array}[(\mathrm{}𝐀),\begin{array}{c}𝐅\hfill \end{array}]+\begin{array}{c}𝐃(\mathrm{}𝐅)\hfill \end{array}.$$
(VI.7)
A similar expression for $`{}_{}{}^{}𝐅`$ is also true. The first term on the right hand side vanishes due to the field equations and the third term is zero due to the previous restriction on $`𝐅`$, (VI.3). Therefore at the isolated horizon,
$$_{\mathrm{}}\begin{array}{c}𝐅\hfill \end{array}\widehat{}=[(\mathrm{}𝐀),𝐅]\text{and}_{\mathrm{}}\begin{array}{c}^{}𝐅\hfill \end{array}\widehat{}=[(\mathrm{}𝐀),{}_{}{}^{}𝐅].$$
(VI.8)
In the Maxwell case, $`𝐅`$ and $`{}_{}{}^{}𝐅`$ are Lie dragged by $`\mathrm{}`$. However, for non-Abelian fields, this is not a gauge invariant statement; the terms on the right hand sides of (VI.8) are necessary for gauge invariance. Although the field strength and its dual are not Lie dragged along $`\mathrm{}^a`$, recalling that $`_{\mathrm{}}{}_{}{}^{2}ϵ\widehat{}=0`$ and using the cyclic property of the trace, it is straightforward to demonstrate that their norms are:
$$_{\mathrm{}}𝐅\widehat{}=0\text{and}_{\mathrm{}}{}_{}{}^{}𝐅\widehat{}=0.$$
(VI.9)
This result, along with (VI.3), guarantees that the charges $`Q_\mathrm{\Delta }^{YM}`$ and $`P_\mathrm{\Delta }^{YM}`$ are independent of the choice of cross section $`S_\mathrm{\Delta }`$ of the horizon.
Let us consider the remaining components of the Yang-Mills field. The boundary conditions place no restrictions on $`𝐅_{ab}n^am^b`$ and $`{}_{}{}^{}𝐅_{ab}^{}n^am^b`$ at all. As in the electromagnetic case, these components describe the radiation flowing along the horizon. The isolated horizon boundary conditions allow radiation arbitrarily close to — and even at — the horizon, provided none crosses it.
We have so far restricted our attention to the field strength and its dual. However, in the action principle and phase space, the basic variable will be the Yang-Mills connection $`𝐀`$. Let us begin with the definition of the Yang-Mills equivalent of the electric potential. Recall that, given an $`\mathrm{}`$, the electric potential was defined in Section II C as $`\mathrm{\Phi }_{(\mathrm{})}\widehat{}=(\mathrm{}𝐀)`$. This definition is not appropriate in the Yang-Mills case since the resulting potential has a free internal index and is therefore not gauge invariant. Instead, we define the Yang-Mills potential, $`\mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}`$ to be negative the norm of $`(\mathrm{}𝐀)`$:
$$\mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}:\widehat{}=(\mathrm{}𝐀).$$
(VI.10)
This gives us a gauge invariant potential at the horizon.
As in Maxwell theory, we need to constrain the form of $`𝐀`$ at the horizon. Several considerations motivate our choice of these boundary conditions. First they must be chosen so that the action principle is well defined. Second, if the gauge group is U(1) the boundary conditions should reduce to those given in Section II C for the electromagnetic field. Finally, we should be able to show that the Yang-Mills electric potential is constant on the horizon. These considerations suggest the following definition:
Definition 6: The connection $`𝐀`$ will be said to be in a gauge adapted to the isolated horizon $`(\mathrm{\Delta },[\mathrm{}])`$ if it satisfies the following two conditions
* (i) The Yang-Mills potential is constant on the horizon
$$d\mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}\widehat{}=0$$
(VI.11)
* (ii) The dual of the field strength $`({}_{}{}^{}𝐅)`$ and $`(\mathrm{}𝐀)`$ point in the same Lie algebra direction,
$$(\mathrm{}𝐀)^i({}_{}{}^{2}ϵ{}_{}{}^{}𝐅)^i$$
(VI.12)
on the horizon.
These boundary conditions satisfy the requirements discussed above. First, it is straightforward to show that in the U(1) case, condition (i) is equivalent to requiring $`_{\mathrm{}}𝐀\widehat{}=0`$ and (ii) is redundant. Second, as we shall see in Section VI B, these boundary conditions are also sufficient to make the variational principle well defined.
It is not difficult to show that these conditions can always be satisfied. The remaining gauge freedom is simply $`𝐀g^1𝐀g+g^1g`$, where $`g`$ satisfies $`_{\mathrm{}}g\widehat{}=0`$.
### B Action and Phase Space
In this section we will consider the first order action for Einstein-Yang-Mills theory on the manifold $``$ described in Section III B. The basic fields will consist of the triplet $`(e_a^I,A_{aI}{}_{}{}^{J},𝐀_a^i)`$, where $`e_a^I`$ and $`A_{aI}^{}{}_{}{}^{J}`$ are the tetrad and Lorentz connection and $`𝐀_a^i`$ is the Yang-Mills connection. The gravitational fields, $`e_a^I`$ and $`A_{aI}^{}{}_{}{}^{J}`$ satisfy the same boundary conditions (at $`\mathrm{\Delta }`$ and infinity) as in Section III B. Furthermore, we require the Yang-Mills fields to be in a gauge adapted to the horizon and assume they fall off sufficiently fast at infinity.<sup>\*‡</sup><sup>\*‡</sup>\*‡More specifically we require the fall off conditions on the Yang-Mills connection to be such that all integrals, in particular the symplectic structure, are finite and yet the asymptotic electric and magnetic charges are not forced to vanish. While it is not trivial to meet these conditions (for examples the conditions used in appear not to lead to a well-defined symplectic structure) they can be met. However, as in the rest of the paper, for brevity, we will not spell out the boundary conditions at infinity in detail. The Einstein-Yang-Mills action is
$$S(e,A,𝐀)=\frac{1}{16\pi G}_{}\mathrm{\Sigma }^{IJ}F_{IJ}+\frac{1}{16\pi G}_\tau _{\mathrm{}}\mathrm{\Sigma }^{IJ}A_{IJ}\frac{1}{8\pi }_{}\mathrm{Tr}[𝐅{}_{}{}^{}𝐅].$$
(VI.13)
The gravitational part of the action has previously been discussed in Section III, therefore we shall only consider in detail the Yang-Mills terms and verify the variational principle is well defined. Taking into account the results of Section III, a variation of the action can be expressed as
$$\delta S(e,A,𝐀)=_{}\text{Equations of Motion}\delta \varphi \frac{1}{4\pi }_\mathrm{\Delta }\mathrm{Tr}[\delta 𝐀{}_{}{}^{}𝐅].$$
(VI.14)
We must demonstrate the boundary term at $`\mathrm{\Delta }`$ vanishes due to the conditions imposed on the Yang-Mills fields. Using (VI.3) and (VI.12), one can show that the trace in (VI.14) can be replaced by a product of norms:
$`\mathrm{Tr}[\delta 𝐀{}_{}{}^{}𝐅]=(\delta \mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}c_\delta \mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}){}_{}{}^{}𝐅{}_{}{}^{\mathrm{\Delta }}ϵ,`$
where $`{}_{}{}^{\mathrm{\Delta }}ϵ=n{}_{}{}^{2}ϵ`$ is the volume form on $`\mathrm{\Delta }`$ and $`\delta \mathrm{}=c_\delta \mathrm{}`$.
In the action principle, variations are performed keeping data fixed at the initial and final slices. In particular $`\delta \mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}`$ and $`c_\delta `$ vanish there. However, the boundary conditions guarantee that $`\mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}`$, and hence its variation, is constant on $`\mathrm{\Delta }`$. Since, $`\delta \mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}`$ and $`c_\delta `$ vanish on the initial cross section of the horizon and are constant, it follows that $`\delta \mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}\widehat{}=0`$ and $`c_\delta \widehat{}=0`$. Therefore, the Yang-Mills horizon boundary term vanishes; the action principle is well defined in the presence of Yang-Mills fields. As in the Einstein-Maxwell case, boundary conditions played a crucial role in demonstrating the viability of the action.
We now wish to construct the covariant phase space and symplectic structure. As before, points in covariant phase space $`\mathrm{\Gamma }`$ will consist of histories which satisfy the Einstein-Yang-Mills field equations, appropriate falloff conditions at infinity and the isolated horizon boundary conditions at $`\mathrm{\Delta }`$. Before proceeding further, we shall once again need to introduce an additional field at the horizon. This can be regarded as ‘potential’ for the Yang-Mills potential $`\mathrm{\Phi }^{\mathrm{YM}}`$. Given any point $`(e,A,𝐀)`$ in the phase space $`\mathrm{\Gamma }`$, let us define the scalar field $`\chi `$ on $`\mathrm{\Delta }`$ as follows: i) $`_{\mathrm{}}\chi =\mathrm{\Phi }_{(\mathrm{})}^{\mathrm{YM}}`$; and ii) $`\chi `$ vanishes on $`S^{}`$, the intersection of $`M^{}`$ with $`\mathrm{\Delta }`$. These conditions are identical to those imposed in the Maxwell case.
Once again, we take second variations of the action in order to obtain a symplectic structure. Since the gravitational terms are exactly the same in Section III, we shall only describe in detail the Yang-Mills part of the symplectic structure. The second variation of the action (VI.13) yields the following symplectic current:
$$J(\gamma ;\delta _1,\delta _2)=J_{\mathrm{grav}}\frac{1}{4\pi }\mathrm{Tr}[\delta _1{}_{}{}^{}𝐅\delta _2𝐀\delta _2{}_{}{}^{}𝐅\delta _1𝐀].$$
(VI.15)
Using the fact that the field equations and linearized field equations are satisfied, one can directly verify that $`J(\gamma ;\delta _1,\delta _2)`$ is indeed a closed 3-form. We again choose the spacetime region of interest, $`\stackrel{~}{}`$, to be that part of the spacetime $``$ bounded by $`M_1`$, $`M_2`$, infinity and a portion $`\stackrel{~}{\mathrm{\Delta }}`$ of the isolated horizon. Integrating $`dJ`$ over $`\stackrel{~}{}`$ and using asymptotic falloff conditions, we obtain
$`({\displaystyle _{M_1}}+{\displaystyle _{M_2}}+{\displaystyle _{\stackrel{~}{\mathrm{\Delta }}}})J(\gamma ;\delta q_1,\delta _2)=0.`$
The integral of $`J`$ over $`\stackrel{~}{\mathrm{\Delta }}`$ does not vanish but, as in Section III, the pull back of $`J`$ to $`\mathrm{\Delta }`$ is exact. Therefore, we can express the integral over $`\stackrel{~}{\mathrm{\Delta }}`$ of $`J`$ as integrals over the initial and final 2-spheres $`S_1`$ and $`S_2`$. Using these results, and keeping track of orientations, we obtain the symplectic structure:
$`\mathrm{\Omega }|_\gamma (\delta _1,\delta _2)=`$ (VI.18)
$`{\displaystyle \frac{1}{16\pi G}}{\displaystyle _M}[\delta _1\mathrm{\Sigma }^{IJ}\delta _2A_{IJ}\delta _2\mathrm{\Sigma }^{IJ}\delta _1A_{IJ}]+{\displaystyle \frac{1}{8\pi G}}{\displaystyle _S}[\delta _1{}_{}{}^{2}ϵ\delta _2\psi \delta _2{}_{}{}^{2}ϵ\delta _1\psi ]`$
$`{\displaystyle \frac{1}{4\pi }}{\displaystyle _M}[\delta _1{}_{}{}^{}𝐅\delta _2𝐀\delta _2{}_{}{}^{}𝐅\delta _1𝐀]+{\displaystyle \frac{1}{4\pi }}{\displaystyle _S}[\delta _1({}_{}{}^{}𝐅{}_{}{}^{2}ϵ)\delta _2\chi \delta _2({}_{}{}^{}𝐅{}_{}{}^{2}ϵ)\delta _1\chi ].`$
Full use of the isolated horizon boundary conditions has been made in obtaining this symplectic structure. In particular, to obtain the given form of the Yang-Mills surface term, we have used the fact that $`(\mathrm{}𝐀)`$ and $`({}_{}{}^{2}ϵ{}_{}{}^{}𝐅)`$ point in the same direction in the Lie algebra. Using field equations one can directly verify that the right side of (VI.18) is independent of the choice of the partial Cauchy surface; the symplectic structure is ‘conserved’. We will use $`(\mathrm{\Gamma },\mathrm{\Omega })`$ as our covariant phase space.
### C Hamiltonian and First Law
In this subsection we will generalize the arguments of Section V to obtain an expression for the energy of the horizon in Einstein-Yang-Mills theory. To do so, we must specify a time evolution vector field $`t^a`$. As before, we require $`t^a`$ to be a member of the preferred equivalence class $`[\mathrm{}^a]`$ at the horizon (this again requires restriction to non-rotating isolated horizons) and approach unit time translation asymptotically. Given $`t^a`$ we can calculate the infinitesimal generator of time evolution, $`\delta _t=(_te,_tA,_t𝐀)`$, and determine whether it is Hamiltonian. Recall that $`\delta _t`$ is Hamiltonian if and only if the 1-form $`X_t`$ on the phase space defined by
$$X_t(\delta ):=\mathrm{\Omega }(\delta ,\delta _t)$$
(VI.19)
is closed. Let us calculate $`X_t`$. The gravitational part will be identical to the expression obtained in Section V, therefore we shall concentrate on the Yang-Mills terms. As with the Maxwell field, the Lie derivatives of the Yang-Mills fields can be re-expressed using the following identities:
$$_t𝐀=t𝐅+𝐃(t𝐀)_t{}_{}{}^{}𝐅=t(𝐃{}_{}{}^{}𝐅)[\mathrm{}𝐀,{}_{}{}^{}𝐅]+𝐃(t{}_{}{}^{}𝐅).$$
(VI.20)
Making use of these expressions, the field equations satisfied by $`(e,A,𝐀)`$ and the linearized field equations for $`\delta `$, we obtain the required expression for $`X_t`$:
$$X_t(\delta )=\frac{1}{16\pi G}_M(tA)\delta \mathrm{\Sigma }(t\mathrm{\Sigma })\delta A\frac{1}{4\pi }_M\mathrm{Tr}[(t𝐀)\delta ({}_{}{}^{}𝐅)(t{}_{}{}^{}𝐅)\delta 𝐀].$$
(VI.21)
As before the expression involves integrals only over the 2-sphere boundaries $`S_{\mathrm{}}`$ and $`S_\mathrm{\Delta }`$ of $`M`$; there is no volume term. The gravitational terms yield $`\delta E_{\mathrm{ADM}}^t`$ at infinity and $`(1/8\pi G)\kappa _{(t)}\delta a_\mathrm{\Delta }`$ at the horizon. The Yang Mills term at infinity vanishes due to fall-off conditions, therefore we need only calculate the Yang-Mills contribution at the horizon. This is composed of two terms, the second of which vanishes due to the restriction (VI.3) which guarantees $`\mathrm{}\begin{array}{c}^{}𝐅\hfill \end{array}\widehat{}=0`$. Since we are in a gauge adapted to the horizon, $`(\mathrm{}𝐀)`$ and $`{}_{}{}^{}𝐅`$ point in the same internal direction. This allows us to replace the trace in the first term by norms:
$`\mathrm{Tr}[(t𝐀)\delta ({}_{}{}^{}𝐅)]=\mathrm{\Phi }_{(t)}^{\mathrm{YM}}\delta ({}_{}{}^{}𝐅{}_{}{}^{2}ϵ)`$
and guarantees that $`\mathrm{\Phi }_{(t)}^{\mathrm{YM}}`$ is constant on the horizon. Making use of the definition of Yang-Mills electric charge, (VI.6), we obtain
$`X_t(\delta )={\displaystyle \frac{1}{8\pi G}}\kappa _{(t)}\delta a_\mathrm{\Delta }\mathrm{\Phi }_{(t)}^{\mathrm{YM}}\delta Q_\mathrm{\Delta }^{\mathrm{YM}}+\delta E_{\mathrm{ADM}}^t.`$
Recall that the necessary and sufficient condition for the existence of a Hamiltonian is that $`X_t`$ be closed. Clearly, this is equivalent to
$$\frac{1}{8\pi G}\mathrm{d}\mathrm{d}\kappa _{(t)}\mathrm{d}\mathrm{d}a_\mathrm{\Delta }+\mathrm{d}\mathrm{d}\mathrm{\Phi }_{(t)}^{\mathrm{YM}}\mathrm{d}\mathrm{d}Q_\mathrm{\Delta }^{YM}=0,$$
(VI.22)
Once again, we conclude that the surface gravity $`\kappa _{(t)}`$ and the Yang-Mills potential $`\mathrm{\Phi }_{(t)}^{\mathrm{YM}}`$ at the horizon defined by $`t^a`$ can depend only upon the area $`a_\mathrm{\Delta }`$ and charge $`Q_\mathrm{\Delta }^{\mathrm{YM}}`$ of the horizon. Finally, (VI.22) is also the necessary and sufficient condition that there exist a function $`E_\mathrm{\Delta }^t`$, also only of $`a_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }^{\mathrm{YM}}`$ such that
$$\delta E_\mathrm{\Delta }^t=\frac{1}{8\pi G}\kappa _{(t)}\delta a_\mathrm{\Delta }+\mathrm{\Phi }_{(t)}^{\mathrm{YM}}\delta Q_\mathrm{\Delta }^{\mathrm{YM}}.$$
(VI.23)
As before, $`E_\mathrm{\Delta }^t`$ is interpreted as the horizon energy defined by the time translation $`t^a`$. We conclude that the vector field $`\delta _t`$ is Hamiltonian if and only if the first law, (VI.23), holds.
We will conclude this section with a few remarks.
i) The derivation of the first law and its final form are completely analogous to those in the Einstein-Maxwell theory. By contrast, in the discussion of the first law for undistorted isolated horizons of , certain restrictions were imposed on the permissible variations $`\delta `$ in the Einstein-Yang-Mills case. In our treatment, subtleties arise only in the definition of a canonical horizon mass (see Section VII C) rather than the discussion of the first law itself.
ii) Although the Yang-Mills magnetic charge $`P_\mathrm{\Delta }^{\mathrm{YM}}`$ will generically not be zero, no term involving $`\delta P_\mathrm{\Delta }^{\mathrm{YM}}`$ arises in the first law.
iii) How does our result compare with those previously available ? In the first law for Yang-Mills fields is proved for globally stationary spacetimes and small perturbations from one such space-time to another. Assuming the Yang-Mills fields fall off sufficiently fast at infinity, (in the non-rotating case) the first law of then reads
$`\delta M={\displaystyle \frac{1}{8\pi G}}\kappa \delta a+{\displaystyle _{\mathrm{Hor}}}\mathrm{Tr}[\varphi \delta {}_{}{}^{}F].`$
Here, $`M`$ is the ADM mass evaluated at infinity, while all terms on the right hand side are evaluated at the horizon. Due to a different gauge choice at the horizon, the authors define a Lie algebra valued potential $`\varphi `$ and leave the ‘$`\mathrm{\Phi }\delta Q`$’ term inside an integral. However, the general form of this first law is the same as ours. In this sense, our framework generalizes the results of to non-static contexts.
In , the first law is proved for globally stationary spacetimes and arbitrary small departures therefrom. However, there are a number of important differences between these results and the ones obtained in this paper. In the non-rotating case, the first law of reads
$`\delta M_{\mathrm{ADM}}+V\delta Q={\displaystyle \frac{1}{8\pi G}}\kappa \delta a`$
where $`V`$ and $`Q`$ are the Yang-Mills potential and charge evaluated at infinity while $`\kappa `$ and $`a`$ are of course evaluated at the horizon. Because of the non-Abelian nature of the Yang-Mills field, unlike in the Maxwell case, the charge $`Q`$ evaluated at infinity is now quite different from the charge evaluated at the horizon and, as in the Maxwell theory, the potential $`V`$ evaluated at infinity has no direct bearing on the potential at the horizon. Furthermore, that calculation makes an essential use of the bifurcation 2-sphere and all fields are required to be smooth there. This restriction implies that the Yang-Mills potential at the horizon vanishes. (The same is true if one restricts the analysis of to the Maxwell case.)
The first law derived in the isolated horizon framework is valid also in presence of radiation in the exterior space-time region and makes no reference to the bifurcation 2-sphere. (Although we restricted ourselves to the non-rotating case, rotation has been incorporated in this framework in .) Furthermore, it has the aesthetically pleasing feature that all quantities that appear in (VI.23) — including the energy $`E_\mathrm{\Delta }^t`$, the potential $`\mathrm{\Phi }_{(t)}^{\mathrm{YM}}`$ and the charge $`Q_\mathrm{\Delta }^{\mathrm{YM}}`$ — are evaluated at the horizon. In particular, one can now meaningfully consider the physical process version in which one does an experiment at the horizon by dropping a test particle/field and changing the horizon charge infinitesimally. More generally (VI.23) is genuinely a law governing the mechanics of the horizon.
## VII Horizon Mass
For notational simplicity, we will say that a (live) vector field $`t^a`$ is permissible if it gives rise to a Hamiltonian evolution. We saw in Sections V and VI that each permissible vector field $`t^a`$ defines a horizon energy $`E_\mathrm{\Delta }^t`$. In the phase space framework, $`E_\mathrm{\Delta }^t`$ has a direct interpretation: it is the surface term at the horizon in the expression of the Hamiltonian generating the $`t^a`$-evolution. However, in many physical applications — such as the study of black hole mergers — one is interested in properties of a specific space-time, rather than the full phase space. Then, it is useful to have at one’s disposal a canonical notion of energy, the analog of the ADM energy in the rest frame at infinity. This quantity could then be interpreted as the horizon mass. In this section, we will introduce this notion in detail. The discussion is divided into three parts. In the first, we consider the Einstein-Maxwell theory; in the second, we discuss dilatonic couplings ; and, in the third, we analyze the Einstein-Yang-Mills system.
### A Einstein-Maxwell theory
In Section V we showed that $`t^a`$ is permissible if and only if (V.8) holds on the phase space. We will now construct a large family of permissible evolution fields $`t^a`$. Fix any regular function $`\kappa _o`$ of two variables $`a_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$. Then, given any point $`\gamma (e,A,𝐀)`$ of $`\mathrm{\Gamma }`$, we define (the boundary value of) the vector field $`t^a`$ as follows. Consider the vector field $`\mathrm{}^a`$ on $`\mathrm{\Delta }`$ defined by the tetrad, $`\mathrm{}^a=e_I^a\mathrm{}^I`$, and denote by $`\kappa _{(\mathrm{})}`$ the surface gravity associated with it. Then, $`\kappa _o=c\kappa _{(\mathrm{})}`$ for some constant $`c`$. Let us set $`t^a=c\mathrm{}^a`$. Repeating this procedure at each phase-space point $`\gamma `$, we obtain a live vector field $`t^a`$ with $`\kappa _{(t)}=\kappa _o`$. (The resulting $`c`$ will be constant on $`\mathrm{\Delta }`$ but a function on the phase space.) Next, consider the electro-magnetic potential, which is guaranteed to be constant on $`\mathrm{\Delta }`$ by our boundary conditions but whose value at any phase space point is so far completely free. We will now use (V.8) to fix it. Equation (V.8) implies:
$`{\displaystyle \frac{\kappa _{(t)}}{Q_\mathrm{\Delta }}}={\displaystyle \frac{\mathrm{\Phi }_{(t)}}{a_\mathrm{\Delta }}}.`$
Since $`\kappa _{(t)}=\kappa _o`$ is known, we can simply integrate the equation for $`\mathrm{\Phi }_{(t)}`$ as a function of $`a_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$. Furthermore, the solution is unique if we impose the physical condition that $`\mathrm{\Phi }_{(t)}`$ should vanish whenever $`Q_\mathrm{\Delta }=0`$. Thus, starting from any regular function $`\kappa _o`$ of $`a_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$, we have obtained a permissible evolution field $`t^a`$. Conversely, it is easy to verify that every permissible vector field arises via this construction. There is clearly a very large family of such live vector fields.
An obvious question is if there is a ‘canonical’ or ‘natural’ choice of $`t^a`$? We will now show that the answer is in the affirmative. Recall that, in the Einstein-Maxwell theory, there is precisely a 2-parameter family of globally static solutions admitting horizons: the Reissner-Nordström family. (Since $`𝐀`$ is required to be a globally defined connection on a trivial $`U(1)`$ bundle, the magnetic charge is zero on the entire phase space.) Let us focus on this family. Denote by $`\xi ^a`$ the static Killing field which is unit at infinity. Its surface gravity is a specific function of $`a_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$:
$`\kappa _{(\xi )}={\displaystyle \frac{1}{2R_\mathrm{\Delta }}}\left(1{\displaystyle \frac{GQ_\mathrm{\Delta }^2}{R_\mathrm{\Delta }^2}}\right).`$
As before, $`R_\mathrm{\Delta }`$ is the horizon radius, defined by $`a_\mathrm{\Delta }=4\pi R_\mathrm{\Delta }^2`$. We can therefore use $`\kappa _{(\xi )}`$ in place of $`\kappa _o`$ in the above construction. The resulting permissible, live vector field $`t_o^a`$ agrees with $`\xi ^a`$ on the horizon of every static solution. This property is satisfied only if we set $`\kappa _o=\kappa _{(\xi )}`$.
Next, we can ‘integrate’ (V.9) to obtain the horizon energy $`E_\mathrm{\Delta }^{t_o}`$. Although a priori there is the freedom to add a constant, we can fix it by requiring that the energy vanish as $`a_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$ tend to zero. Indeed, we have no choice in this since one cannot construct a quantity with dimensions of mass from the fundamental constants that appear in the Einstein-Maxwell theory. (Einstein-Yang-Mills theory does admit such a constant and we will see in Section VII C that it leads to an interesting modification of the situation discussed here.) Let us define the horizon mass via
$`M_\mathrm{\Delta }=E_\mathrm{\Delta }^{t_o}.`$
To justify this definition, let us begin by restricting to static solutions. In each static solution, we are free to extend $`t_o^a`$ such that it coincides with the Killing field $`\xi ^a`$. General symplectic arguments imply that, on any connected component of the space of static solutions, the numerical value of the total Hamiltonian, generating evolution along $`\xi ^a`$, must be constant (see, e.g., .) In the Einstein-Maxwell case, there is a single connected component and, by the dimensional argument given above, the numerical value of the Hamiltonian must vanish on it. Hence, from (V.10) it follows that, on any static solution,
$`H_{t_o}=M_{\mathrm{ADM}}M_\mathrm{\Delta }=0.`$
On a general solution, of course, $`M_{\mathrm{ADM}}`$ would be greater than $`M_\mathrm{\Delta }`$, the difference being equal to the energy in radiation. If the horizon is complete in the future and time-like infinity $`i^+`$ satisfies certain regularity conditions, as in one can argue that the difference is precisely the total energy radiated across $`^+`$ and hence $`M_\mathrm{\Delta }`$ equals the future limit of the Bondi mass. These considerations support our interpretation of $`M_\mathrm{\Delta }`$ as the horizon mass.
Finally, since we now have a canonical evolution field $`t_o^a`$, we can drop the suffix $`t`$ on surface gravity and electromagnetic potential and write the first law (V.9) in the more familiar form:
$`\delta M_\mathrm{\Delta }={\displaystyle \frac{1}{8\pi G}}\kappa \delta a_\mathrm{\Delta }+\mathrm{\Phi }\delta Q_\mathrm{\Delta }.`$
In contrast to treatments based on static space-times, the quantities that enter this law are all defined at the horizon. Therefore, as pointed out in , it is now possible to interpret this law also in the ‘active’ sense where one considers physical processes which increase the area and the charge of a given horizon. To our knowledge, the standard proofs of this physical version are not applicable to processes in which the background has non-zero electric charge and the process changes it infinitesimally.
We will conclude with a few remarks.
i) In the above discussion, the permissible evolution field $`t_o^a`$ was constructed by setting $`t_o^a=c\mathrm{}^a`$ where $`c`$ is given by $`\kappa _o(1/2R_\mathrm{\Delta })(1G(Q_\mathrm{\Delta }/R_\mathrm{\Delta })^2)=c\kappa _{(\mathrm{})}`$. For $`c`$ to be well-defined, it is necessary that $`\kappa _o`$ vanishes whenever $`\kappa _{(\mathrm{})}`$ does. Therefore, for the mass to be well-defined, we must excise those points from the phase space at which $`\kappa _{(\mathrm{})}`$ vanishes but $`\kappa _o`$ does not. However, this is not a serious limitation. In particular, we still retain all static solutions including the extremal ones at which $`\kappa _o`$ vanishes.
ii) Since we have a specific $`\kappa _o`$, we can use (V.8) to obtain the corresponding electrostatic potential: $`\mathrm{\Phi }_o=Q_\mathrm{\Delta }/R_\mathrm{\Delta }`$. Furthermore, by integrating (V.9) it is easy to express $`M_\mathrm{\Delta }`$ explicitly in terms of the horizon parameters:
$$M_\mathrm{\Delta }=\frac{1}{4\pi G}\kappa a_\mathrm{\Delta }+\mathrm{\Phi }Q_\mathrm{\Delta }=\frac{R_\mathrm{\Delta }}{2G}\left(1+\frac{GQ_\mathrm{\Delta }^2}{R_\mathrm{\Delta }^2}\right).$$
(VII.1)
Thus, the functional dependence of $`M_\mathrm{\Delta }`$ on the horizon parameters at any point of the phase space is the same as in static space-times. Note that this is a result of the framework, not an assumption. Its derivation involved two distinct steps. First, and most importantly, the first law (V.9) arose as a necessary and sufficient condition for the existence of a consistent Hamiltonian framework. Second, the freedom in $`t^a`$ was exploited in order to construct the preferred, permissible evolution field $`t_o^a`$. It is quite significant that $`M_\mathrm{\Delta }`$ can be expressed so simply using just the parameters defined locally at the horizon even when there is radiation arbitrarily close to it. This fact is likely to play an important role in the problem of extracting physics in the strong field regimes from numerical simulations of black hole collisions . It is important to notice that although we made use of our knowledge of static solutions to arrive at a canonical $`t_o^a`$ and the mass function $`M_\mathrm{\Delta }`$, the final result (VII.1) makes no reference to these solutions. $`M_\mathrm{\Delta }`$ is a simple function of the parameters which can be directly computed from the geometry of any one isolated horizon.
iii) In the earlier work on undistorted horizons, one restricted oneself to the preferred evolution field $`t_o^a`$ from the very beginning (although this vector field was selected using a different but equivalent procedure).This strategy seems to have generated a misunderstanding (see, e.g.,) that the first law was obtained in merely by identifying the parameters labeling a general isolated horizon with those of static horizons and then using the Smarr formulas available in the static context. This was not the case. Rather, static solutions were used only to select the appropriate normalization of the evolution vector field $`t^a`$ at the horizon. The Hamiltonian framework was then used to define the horizon mass without any reference to Smarr formulas. As in this section, the mass was then shown to reproduce the Smarr-type formulas on general horizons. The a priori freedom in the choice of a permissible $`t^a`$ was not discussed and the first law appeared only in the more familiar form, given above.
### B Dilatonic coupling
The Einstein-Maxwell-dilaton system was studied in some detail in the undistorted case in . We will revisit it here in the more general context considered in this paper because it brings out a subtlety in the definition of the horizon mass $`M_\mathrm{\Delta }`$ and the associated first law.
The dilaton is a scalar field $`\varphi `$ which can couple to the Maxwell field in a non-standard fashion. The coupling is governed by a constant $`\alpha `$. If $`\alpha `$=$`0`$, one obtains the standard Einstein-Maxwell-Klein-Gordon theory and the situation then is completely analogous to the Einstein-Maxwell theory considered above. If $`\alpha `$=$`1`$, the theory represents the low energy limit of string theory. In this case, there are some interesting differences from the Einstein-Maxwell theory considered in this paper. To bring out these differences, in this subsection we will set $`\alpha `$=$`1`$. (The situation for a general value of $`\alpha `$ is discussed in where one can also find details on the material summarized below.)
In the standard formulation, the theory has three charges, all defined at infinity; the ADM mass $`M_{\mathrm{ADM}}`$, the usual electric charge $`Q_{\mathrm{}}`$ and another charge $`\stackrel{~}{Q}_{\mathrm{}}`$:
$`Q_{\mathrm{}}={\displaystyle \frac{1}{4\pi }}{\displaystyle _S_{\mathrm{}}}{}_{}{}^{}𝐅\mathrm{and}\stackrel{~}{Q}_{\mathrm{}}={\displaystyle \frac{1}{4\pi }}{\displaystyle _S_{\mathrm{}}}e^{2\varphi }{}_{}{}^{}𝐅.`$
$`\stackrel{~}{Q}`$ is conserved in space-time (i.e. its value does not change if the 2-sphere of integration is deformed) while $`Q`$ is not. From the perspective of the isolated horizons, it is more useful to use $`a_\mathrm{\Delta },Q_\mathrm{\Delta },\stackrel{~}{Q}_\mathrm{\Delta }`$ as the basic chargesIn the undistorted case, the dilaton is constant on $`\mathrm{\Delta }`$ and hence we can replace $`Q_\mathrm{\Delta }`$ by $`\varphi _\mathrm{\Delta }`$ as in .:
$`Q_\mathrm{\Delta }={\displaystyle \frac{1}{4\pi }}{\displaystyle _{S_\mathrm{\Delta }}}{}_{}{}^{}𝐅,\mathrm{and}\stackrel{~}{Q}_\mathrm{\Delta }={\displaystyle \frac{1}{4\pi }}{\displaystyle _{S_\mathrm{\Delta }}}e^{2\varphi }{}_{}{}^{}𝐅.`$
Although the standard electric charge is not conserved in space-time, it is conserved along $`\mathrm{\Delta }`$ whence $`Q_\mathrm{\Delta }`$ is well-defined.
It is straightforward to extend the construction of the phase space to include the dilaton. The only difference is that the charge $`Q`$ in equations (V.6) - (V.9) is replaced by $`\stackrel{~}{Q}`$. With this minor change, the discussion of the first part of Section VII A is also unaffected. Thus, given any function $`\kappa _o`$ of $`a_\mathrm{\Delta }`$ and $`\stackrel{~}{Q}_\mathrm{\Delta }`$, we can construct a permissible, (live) evolution field $`t^a`$.
The difference arises in the next step where we constructed a preferred $`t_o^a`$. With the dilatonic coupling, the theory has a unique three parameter family of static solutions which can be labeled by $`(a_\mathrm{\Delta },Q_\mathrm{\Delta },\stackrel{~}{Q}_\mathrm{\Delta })`$. As in the Reissner Nordström family, these solutions are spherically symmetric. In terms of these parameters, the surface gravity $`\kappa _{(\xi )}`$ of the static Killing field which is unit at infinity is given by:
$`\kappa _{(\xi )}={\displaystyle \frac{1}{2R_\mathrm{\Delta }}}\left[1+2G{\displaystyle \frac{Q_\mathrm{\Delta }\stackrel{~}{Q}_\mathrm{\Delta }}{R_\mathrm{\Delta }^2}}\right]\left[12G{\displaystyle \frac{Q_\mathrm{\Delta }\stackrel{~}{Q}_\mathrm{\Delta }}{R_\mathrm{\Delta }^2}}\right]^{\frac{1}{2}}.`$
The problem in the construction of the preferred $`t_o^a`$ is that we need a function $`\kappa _o`$ which depends only on $`a_\mathrm{\Delta }`$ and $`\stackrel{~}{Q}_\mathrm{\Delta }`$. Therefore, we can no longer set $`\kappa _o=\kappa _{(\xi )}`$ on the entire phase space because $`\kappa _{(\xi )}`$ depends on all three horizon parameters.
To extract the mass function $`M_\mathrm{\Delta }`$ on the phase space, we can proceed as follows. Let us foliate $`\mathrm{\Gamma }`$ by $`Q_\mathrm{\Delta }=\mathrm{const}`$ surfaces. On each leaf, $`\kappa _{(\xi )}`$ trivially depends only on $`a_\mathrm{\Delta }`$ and $`\stackrel{~}{Q}_\mathrm{\Delta }`$ and so we can set $`\kappa _o=\kappa _{(\xi )}`$. Therefore, by the procedure outlined in Section VII A, we obtain a (live) vector field $`t_o^a`$ and can define the mass $`M_\mathrm{\Delta }(\gamma )=E_\mathrm{\Delta }^{t_o}(\gamma )`$ for all points $`\gamma `$ on this leaf. Repeating this procedure for each leaf, we obtain a live vector field $`t_o^a`$ and a mass function $`M_\mathrm{\Delta }`$ everywhere on $`\mathrm{\Gamma }`$. However, the surface gravity $`\kappa _{(t_o)}`$ now depends on all three parameters, rather than just $`a_\mathrm{\Delta }`$ and $`\stackrel{~}{Q}_\mathrm{\Delta }`$. Therefore, the first law (V.9) cannot hold for arbitrary variations $`\delta `$ and consequently $`\delta _{t_o}`$ fails to be a Hamiltonian vector field. Put differently, although there is a multitude of permissible, live vector fields, each leading to a first law, none of them can coincide with the Killing field $`\xi ^a`$ (which is unit at infinity) on all static solutions. This is a significant departure from the Einstein-Maxwell case considered above.
Nonetheless, (modulo the caveat discussed in the first remark at the end of Section VII A) the above procedure does provide us with a well-defined mass function $`M_\mathrm{\Delta }`$ on the entire phase space which can be expressed in terms of the horizon parameters as
$`M_\mathrm{\Delta }={\displaystyle \frac{R_\mathrm{\Delta }}{2G}}\left[12G{\displaystyle \frac{Q_\mathrm{\Delta }\stackrel{~}{Q}_\mathrm{\Delta }}{R_\mathrm{\Delta }^2}}\right]^{\frac{1}{2}}.`$
It equals $`M_{\mathrm{ADM}}`$ in static space-times and has other properties which motivated our interpretation of $`M_\mathrm{\Delta }`$ as the horizon mass in the Einstein-Maxwell case. Since this function is well-defined on the entire phase space, we can simply vary it and express the result in terms of the horizon parameters. The result is:
$`\delta M_\mathrm{\Delta }={\displaystyle \frac{1}{8\pi G}}\kappa \delta a_\mathrm{\Delta }+\widehat{\mathrm{\Phi }}\delta \widehat{Q}_\mathrm{\Delta }`$
where $`\kappa =\kappa _{(t_o)}`$, $`\widehat{\mathrm{\Phi }}^2=(Q_\mathrm{\Delta }\stackrel{~}{Q}_\mathrm{\Delta }/R_\mathrm{\Delta }^2)`$ and $`\widehat{Q}_\mathrm{\Delta }^2=Q_\mathrm{\Delta }\stackrel{~}{Q}_\mathrm{\Delta }`$. Thus, although there is still a first law in terms of $`t_o^a`$ and $`M_\mathrm{\Delta }`$, it does not have the canonical form (V.9) because $`t_o^a`$ is not a permissible vector field. More generally, in theories with multiple scalar fields , if one focuses only on static sectors, one obtains similar ‘non-standard’ forms of the first law with work terms involving scalar fields. This reflects the fact that there is no permissible vector field $`t^a`$, defined for all points of the phase space, which coincides with the properly normalized Killing field on all static solutions. In the undistorted case, the analysis was carried out only in terms of the vector field $`t_o^a`$ and the horizon mass $`M_\mathrm{\Delta }`$ . The resulting first law had the above form.
Alternatively, one can restrict oneself to variations $`\overline{\delta }`$ which are tangential to the leaves of the phase space foliation constructed above. Since $`t_o^a`$ is a permissible vector field for any one leaf, we obtain the standard first law
$`\overline{\delta }M_\mathrm{\Delta }=(1/8\pi G)\kappa _{(t_o)}\overline{\delta }a_\mathrm{\Delta }+\mathrm{\Phi }_{(t_o)}\overline{\delta }\stackrel{~}{Q}_\mathrm{\Delta }`$
for the restricted variations. The idea of using such restricted variations was suggested in in the context of Yang-Mills fields (although the foliations and other details were not spelled out there).
To summarize, because there is now a three parameter family of static solutions rather than two —or, more precisely, because the standard surface gravity $`\kappa _{(\xi )}`$ in static space-times depends on $`a_\mathrm{\Delta },\stackrel{~}{Q}_\mathrm{\Delta }`$ and $`Q_\mathrm{\Delta }`$— a canonical, permissible evolution field is no longer available. However, there is still a multitude of permissible evolution fields and corresponding first laws. Furthermore, one can still define a canonical mass function $`M_\mathrm{\Delta }`$ on the entire phase space.
### C Yang-Mills fields
In Einstein-Maxwell theory, with and without the dilaton, there is no way to construct a quantity with the dimensions of mass from the fundamental constants in the theory. The situation is different for Einstein-Yang-Mills theory because the coupling constant $`g`$ has dimensions $`(LM)^{1/2}`$. The existence of a quantity with units of mass has interesting consequences which we will now discuss.
Let us begin with a summary of the known static solutions in Yang-Mills theory. First, the Reissner-Nordström family constitutes a continuous two parameter set of static solutions of the Einstein-Yang-Mills theory, labelled by $`(a_\mathrm{\Delta },Q_\mathrm{\Delta }^{\mathrm{YM}})`$. In addition, there is a 1-parameter family of ‘embedded Abelian solutions’ with (a fixed) magnetic charge $`P_\mathrm{\Delta }^o`$, labelled by $`(a_\mathrm{\Delta },P_\mathrm{\Delta }^o)`$. Finally, there are families of ‘genuinely non-Abelian solutions’. For these, the analog of the Israel theorem for Einstein-Maxwell theory fails to hold ; the theory admits static solutions which need not be spherically symmetric. In particular, an infinite family of solutions labelled by two integers $`(n_1,n_2)`$ is known to exist. All static, spherically symmetric solutions are known and they correspond to the infinite sub-family $`(n_1,n_2=0)`$, labelled by a single integer. However, the two parameter family is obtained using a specific ansätz, so there may well exist other static solutions. Although the available information on static solutions is quite rich, in contrast to the Einstein-Maxwell-Dilaton system, one is still rather far from having complete control of the static sector of the Einstein-Yang-Mills theory.
The zeroth and first laws do hold in the Einstein-Yang-Mills case. At present, however, we can only hope to repeat the strategy used in the last two sub-sections to define a canonical mass function $`M_\mathrm{\Delta }`$ on portions of the phase space. In order to define it on the full phase space, of the ‘uniqueness’ and ‘completeness’ conjectures of will have to hold (possibly with a suitable modification.<sup>\*∥</sup><sup>\*∥</sup>\*∥For example it may be appropriate to restrict oneself to the class of space-times admitting isolated horizons which are complete in the future. Physically, this is the most interesting case since such horizons would represent the asymptotic geometry resulting from a gravitational collapse or black hole mergers. Nevertheless, new insight into the static solutions can be obtained by restricting attention to certain leaves of the phase space. The basic idea is taken from but applied in a slightly different manner to the more general context of distorted horizons.
Consider a connected component of the known static solutions, labelled by $`\stackrel{}{n}(n_1,n_2)`$. This is a 1-dimensional sub-space of the phase space which we denote $`S_\stackrel{}{n}`$. Each point in $`S_\stackrel{}{n}`$ can be labelled by the value of the horizon area $`a_\mathrm{\Delta }`$. Calculate the surface gravity $`\kappa _{(\xi )}`$ for this family, where $`\xi ^a`$ is the static Killing field which it unit at infinity and set $`\kappa _o=\kappa _{(\xi )}`$ in the construction sketched in Section VII A. We then obtain a live vector field $`t_o^a`$, and the corresponding first law
$`\delta E_\mathrm{\Delta }^{t_o}={\displaystyle \frac{1}{8\pi G}}\kappa _{(t_o)}\delta a_\mathrm{\Delta }`$
on the full phase space.
When restricted to $`S_\stackrel{}{n}`$, we can interpret $`E_\mathrm{\Delta }^{t_o}`$ as the horizon mass $`M_\mathrm{\Delta }^{(\stackrel{}{n})}`$ and replace $`\kappa _{(t_o)}`$ by the function $`\beta _\stackrel{}{n}(a_\mathrm{\Delta })`$ used in the literature: $`\beta _\stackrel{}{n}=2\kappa _{(t_o)}R_\mathrm{\Delta }`$. Then, by integrating the first law along $`S_\stackrel{}{n}`$, one obtains:
$`M_\mathrm{\Delta }^{(\stackrel{}{n})}={\displaystyle \frac{1}{2G}}{\displaystyle _0^{R_\mathrm{\Delta }}}\beta _\stackrel{}{n}(x)𝑑x`$
where, we have used the fact that, since $`E_\mathrm{\Delta }^{(t_o)}`$ is a surface integral at $`\mathrm{\Delta }`$, it vanishes as the horizon area goes to zero. Thus, the horizon mass is completely determined by $`\beta _\stackrel{}{n}(a_\mathrm{\Delta })`$.
Next, we use the fact that the Hamiltonian given by $`H_{t_o}=E_{\mathrm{ADM}}^{t_o}E_\mathrm{\Delta }^{t_o}`$ (see (V.10)) is constant on each connected, static sector if $`t_o^a`$ coincides with the static Killing field on the entire sector. By construction, our $`t_o^a`$ has this property. In the Einstein-Maxwell case, since there is no constant with the dimension of energy, it follows that the restriction of $`H_{\mathrm{ADM}}^{t_o}`$ to the static sector must vanish. The situation is quite different in Einstein-Yang-Mills theory where the Yang-Mills coupling constant $`g`$ provides a scale. In $`c=1`$ units, $`(g\sqrt{G})^1\mathrm{Mass}`$. Therefore, we can only conclude
$`M_{\mathrm{ADM}}^{(\stackrel{}{n})}=M_\mathrm{\Delta }^{(\stackrel{}{n})}+(g\sqrt{G})^1C^{(\stackrel{}{n})}`$
for some $`\stackrel{}{n}`$-dependent constant $`C^{(\stackrel{}{n})}`$. As the horizon radius shrinks to zero, the static solution under consideration tends to the solitonic solution with the same ‘quantum numbers’ $`\stackrel{}{n}`$. Hence, by taking this limit, we conclude $`(g\sqrt{G})^1C^{(\stackrel{}{n})}=M_{\mathrm{ADM}}^{\mathrm{soliton},(\stackrel{}{n})}`$. Therefore, we have the following interesting relation between the black hole and solitonic solutions:
$`M_{\mathrm{ADM}}^{\mathrm{BH},(\stackrel{}{n})}={\displaystyle \frac{1}{2G}}{\displaystyle _0^{R_\mathrm{\Delta }}}\beta _\stackrel{}{n}(x)𝑑x+M_{\mathrm{ADM}}^{\mathrm{soliton},(\stackrel{}{n})}`$
where the integral of $`\beta _\stackrel{}{n}`$ is evaluated on the 1-dimensional ‘parameter space’ of $`S_\stackrel{}{n}`$ (given by the horizon radius). Furthermore, as is clear from the above discussion, the ADM mass of the soliton is a multiple of $`(g\sqrt{G})^1`$. Thus, somewhat surprisingly, the derivation of the first law in the isolated horizon framework has led to an interesting relation between the ADM masses of black holes and their solitonic analogs in the static sector.
## VIII Discussion
In the first part of this paper, we introduced the notions of weakly isolated and isolated horizons. In contrast with earlier work , the definitions allow for the possible presence of distortion and rotation at the horizon. In addition, the present definitions are more geometric and intrinsic; in particular, they never refer to a foliation.
The notion of an isolated horizon, unlike that of an event horizon, is completely quasi-local. One can test if a given 3-surface in space-time is (weakly) isolated or not simply by examining space-time geometry at the surface. Furthermore, space-times admitting an isolated horizon $`\mathrm{\Delta }`$ need not admit any Killing field even in a neighborhood of $`\mathrm{\Delta }`$. In particular, they can admit radiation in the exterior region. Therefore, such space-times can serve as more realistic models of late stages of a gravitational collapse or black hole merger. In the near $`\mathrm{\Delta }`$ geometry of vacuum solutions is examined in detail using similar techniques to those used at null infinity. The resulting structure — presence of a preferred ‘rest frame’, constrains on possible isometries, Bondi-type expansions of the metric — should be useful to extract physics in the strong field regime of general relativity, especially in the problem of binary black hole collisions.
This paper, however, focused on another aspect of isolated horizons: extensions of the zeroth and first laws of black hole mechanics. All previous discussions of these laws were restricted to perturbations of stationary black holes. Using Lagrangian and Hamiltonian frameworks, we extended these laws to arbitrary space-times admitting non-rotating isolated horizons in Einstein-Maxwell-dilaton and Einstein-Yang-Mills theory. Furthermore, the analysis suggests that it should be rather easy to incorporate other forms of matter, provided they admit Lagrangian and Hamiltonian descriptions.
The generalization of black hole mechanics presented in this paper has several interesting features. First, all quantities that enter the first laws are defined locally at the horizon $`\mathrm{\Delta }`$. In standard treatments, some quantities such as area and surface gravity are defined at the horizon. Others, like energy and sometimes even the Yang-Mills/Maxwell charge and potential, are evaluated at infinity. In part because of this ‘mismatch’, to our knowledge the ‘physical process version’ of the first law had not previously been established for processes which change the charge of the black hole. Since all quantities in the present treatment are defined locally at the horizon, it is now straightforward to establish the law for such processes . Secondly, other treatments based on a Hamiltonian framework often critically use the bifurcate 2-surface which does not exist in the extremal case. Therefore, extremal black holes are often excluded from the first law. The present analysis never makes reference to bifurcate surfaces (which do not exist in physical space-times resulting from gravitational collapse). Therefore, our discussion of the first law holds also in the extremal case. Thirdly, with obvious modifications of boundary conditions at infinity, our analysis includes cosmological horizons where thermodynamic considerations are also applicable .
Finally, and perhaps most importantly, our analysis sheds new light on the ‘origin’ of the first law: it arose as a necessary and sufficient condition for the existence of a Hamiltonian generating time evolution. A new feature of our framework is the existence of an infinite family of first laws corresponding to the infinite family of ‘permissible’ vector fields $`t^a`$. (A vector field $`t^a`$ is permissible if it is Hamiltonian, that is, induces canonical transformations on the phase space.) In theories where we have sufficient control on the space of static solutions, such as Einstein-Maxwell, one can select a natural evolution field $`t_o^a`$. Corresponding to evolution along this $`t_o^a`$, there is a canonical notion of energy which can be interpreted as the mass of the isolated horizon. There exist also preferred values of surface gravity and electric potential and a canonical first law. This additional structure is extremely useful in other applications of the framework, such as extraction of physical information from numerical simulations of black-hole collisions. However, it is not essential to the discussion of mechanics: our derivation of the first law in Sections V and VI does not require any knowledge of the static sector of the theory.
The Hamiltonian approach to black hole mechanics has appeared in the literature before, most notably in the work of Brown and York . The spirit of the Brown-York approach is similar to ours. In particular, they do not restrict themselves to stationary situations. However, in that work, the focus is on an outer, time-like boundary whereas our focus is on the inner, null boundary representing the isolated horizon. Conserved quantities in presence of internal boundaries were recently discussed also by Julia and Silva in a more general context of theories with gauge symmetries. As in our framework, their treatment exploits the simplifications that occur in a first order formalism and the final surface-integral expressions of conserved charges are dictated by the precise boundary conditions imposed at the internal boundaries. Their treatment is based on superpotentials and is thus complements the Hamiltonian methods used here and in .
In this paper, the Lagrangian and Hamiltonian frameworks are based on real tetrads and Lorentz connections. It is therefore quite straightforward to extend our analysis to any space-time dimension. Indeed, it has already been extended to 2+1 dimensions in . However, our phase space — and especially the explicit symplectic structure used here — is tailored to the Einstein-matter system. While it should be possible to extend it to higher derivative theories of gravity as in , that task would not be as simple.
## Acknowledgments
We are grateful to Chris Beetle and Jerzy Lewandowski for countless discussions. We thank Alejandro Corichi and Daniel Sudarsky for stimulating correspondence, Jiri Bicak and Werner Israel for pointing out references on distorted black holes and Thibault Damour, Sean Hayward and Bob Wald for helpful comments. This work was supported in part by the NSF grants PHY94-07194, PHY95-14240, INT97-22514 and by the Eberly research funds of Penn State. SF was supported in part by a Braddock Fellowship.
## A Examples of distorted horizons
Because of the no hair theorems in the Einstein-Maxwell theory, distorted horizons have received a rather limited attention in the literature. Therefore, in this appendix we will discuss a few explicit examples in Einstein-Maxwell theory. For a general construction and an existence result, see .
To obtain explicit solutions, one has to impose symmetries. All solutions considered in this section will be static and axi-symmetric. As one would expect from the no-hair theorems, they fail to be asymptotically flat, whence they fail to represent isolated black holes in the standard sense. Nonetheless, they all satisfy the isolated horizon boundary conditions. That framework also serves to ‘explain’ the otherwise surprising feature that the surface gravity of these solutions depends only on the area and the charge and is insensitive to the distortion parameters.
In the literature on static, distorted black holes, it is generally assumed that the solution is valid only in a finite region around the horizon and its distant behavior is suitably modified by the far-away matter which causes the distortion. For undistorted isolated horizons, Robinson-Trautman space-times offer interesting examples of vacuum, asymptotically flat solutions which admit isolated horizons but no Killing fields whatsoever. Distorted analogs of these solutions are not known but presumably exist. It would be interesting to find them.
### 1 Black Hole in a magnetic universe
Let us begin with a simple example: an uncharged black hole in an ‘external magnetic field’ which distorts the horizon. The specific solution we wish to consider is static and axisymmetric and was first obtained in the Ernst-potential framework . The magnetic field is non-zero on the horizon. Thus, one has to consider the full set of Einstein-Maxwell equations on the horizon.
The space-time has topology $`S^2\times \text{I}\text{R}^2`$ and the metric is given by
$$ds^2=F^2\left[\left(1\frac{2M}{r}\right)dt^2+\frac{dr^2}{12M/r}+r^2d\theta ^2\right]+\frac{r^2\mathrm{sin}^2\theta }{F^2}d\varphi ^2$$
(A.1)
where
$`F=1+{\displaystyle \frac{1}{4}}B_0^2r^2\mathrm{sin}^2\theta .`$
$`B_0`$ is a constant and on the axis the magnetic field is given by $`B=B_0dr`$. Because $`F`$ diverges at infinity, the metric fails to be asymptotically flat. For $`M=0`$, the metric reduces to that of the Melvin universe and for $`M0`$ it admits a Killing horizon at $`r=2M`$. To examine the behavior at the horizon, let us first cast the metric in the Eddington-Finkelstein coordinates $`(v,r,\theta ,\varphi )`$ where $`dv=dt+(1\frac{2M}{r})^1dr`$:
$$ds^2=F^2\left(1\frac{2M}{r}\right)dv^2+2F^2dvdr+F^2r^2d\theta ^2+\frac{r^2\mathrm{sin}^2\theta }{F^2}d\varphi ^2.$$
(A.2)
Since the metric is not asymptotically flat, the standard procedure of normalizing the Killing field to be unit at infinity is not applicable. Thus, we only have an equivalence class $`[\mathrm{}]`$ of (preferred) null normals to the horizon, $`\mathrm{}\frac{}{v}`$. Let $`\mathrm{\Delta }`$ be the Killing horizon and assume the Killing field $`\frac{}{v}`$ is a member of the equivalence class $`[\mathrm{}]`$. It follows trivially that $`(\mathrm{\Delta },[\mathrm{}])`$ is a non-rotating isolated horizon. If $`B_00`$, the scalar curvature $`{}_{}{}^{2}R`$ of the horizon 2-metric has $`\theta `$-dependence; the horizon is distorted. However, an explicit calculation shows that, as in the Schwarzschild space-time, the surface gravity $`\kappa `$ is given by $`1/2r1/2R_\mathrm{\Delta }`$ and the electrostatic potential $`\mathrm{\Phi }`$ vanishes on $`\mathrm{\Delta }`$. At first, it is quite surprising that while the presence of distortion affects $`\mu `$, $`\mathrm{\Phi }_{11}`$, $`\mathrm{\Psi }_2`$ and $`{}_{}{}^{2}R`$, it does not affect $`\kappa `$ or $`\mathrm{\Phi }`$. However, as we saw in Section V, this result is to be expected from the general Hamiltonian considerations.
### 2 Distorted black holes as special cases of Weyl solutions
In this sub-section, we will review the construction of a large family of distorted black holes starting from Weyl solutions and a recent generalization of these results to include electric charge .
A general class of static, axisymmetric spacetimes was found by Weyl in 1917 . The metric for such a spacetime can be cast in the following form:
$$ds^2=e^{2\psi }dt^2+e^{2(\gamma \psi )}(d\rho ^2+dz^2)+e^{2\psi }\rho ^2d\varphi ^2$$
(A.3)
where $`\psi `$ and $`\gamma `$ are smooth functions of $`\rho `$ and $`z`$. Einstein’s vacuum equations expressed in terms of $`\psi `$ and $`\gamma `$ take a particularly simple form. The equation for $`\psi `$,
$$\psi _{,\rho \rho }+\frac{\psi _{,\rho }}{\rho }+\psi _{,zz}=0,$$
(A.4)
is simply the Laplace equation in flat space with cylindrical coordinates $`(\rho ,z,\varphi )`$. (In addition, $`\psi `$ has to be independent of the angular coordinate $`\varphi `$.) Given a solution for $`\psi `$, the function $`\gamma `$ can be determined by simple integration:
$`\gamma _{,\rho }`$ $`=`$ $`\rho [\psi _{,\rho }^2\psi _{,z}^2]`$ (A.5)
$`\gamma _{,z}`$ $`=`$ $`2\rho [\psi _{,\rho }\psi _{,z}]`$ (A.6)
The Schwarzschild metric is of course a particular solution to these equations and corresponds to choosing for $`\psi `$ and $`\gamma `$:
$$\psi =\psi _\mathrm{S}:=\frac{1}{2}\mathrm{ln}\left(\frac{LM}{L+M}\right)\text{and}\gamma =\gamma _\mathrm{S}:=\frac{1}{2}\mathrm{ln}\left(\frac{L^2M^2}{L^2\eta ^2}\right),$$
(A.7)
where $`L=\frac{1}{2}(l_++l_{}),\eta =\frac{1}{2}(l_+l_{})`$ with $`l_+=\sqrt{\rho ^2+(z+M)^2}`$ and $`l_{}=\sqrt{\rho ^2+(zM)^2}`$ and $`M`$ is the mass of the Schwarzschild solution. Note that $`\psi `$ is just the Newtonian potential due to a rod of length $`2M`$ placed symmetrically about the origin on the z-axis. Both $`\psi _S`$ and $`\gamma _S`$ diverge logarithmically in the limit $`\rho 0`$ (for $`zM`$). In order to recast this solution in the standard Schwarzschild form, one must transform from $`(z,\rho )`$ to the Schwarzschild coordinates $`(r,\theta )`$ by
$$z=(rM)\mathrm{cos}\theta \text{and}\rho ^2=r^2(12M/r)\mathrm{sin}^2\theta $$
(A.8)
This coordinate transformation shows that the horizon, $`r=2M`$, corresponds to the line segment $`\rho =0`$, $`zM`$ in Weyl coordinates. Therefore, the Weyl coordinates cover only the exterior of the horizon.
Now, the key point is that (A.4), the only field equation one has to solve, is linear. Hence we can ‘distort’ the Schwarzschild solution simply by adding to $`\psi _\mathrm{S}`$ any solution $`\psi _D`$ of the flat space Laplace equation which is regular along the $`z`$axis . Thus, we can set
$$\psi =\psi _\mathrm{S}+\psi _D\text{and}\gamma =\gamma _\mathrm{S}+\gamma _D.$$
(A.9)
Substituting these expressions into (A.4) and (A.5) and using the forms of the Schwarzschild functions $`\psi _S`$ and $`\gamma _S`$, one can show that at $`\rho =0`$,
$$\gamma _D_{\rho =0}\widehat{}=2\psi _D_{\rho =0}.$$
(A.10)
This fact plays an important role in analyzing the horizon structure.
In Schwarzschild coordinates, the distorted metric takes the form
$$ds^2=e^{2\psi _D}(12M/r)dt^2+\frac{e^{2(\gamma _D\psi _D)}}{(12M/r)}dr^2+e^{2(\gamma _D\psi _D)}r^2d\theta ^2+r^2\mathrm{sin}^2\theta e^{2\psi _D}d\varphi ^2.$$
(A.11)
As usual, the metric has a coordinate singularity at $`r=2M`$. Let us therefore introduce the Eddington-Finkelstein coordinate $`v`$ as before. The metric can be re-expressed in $`(v,r,\theta ,\varphi )`$ coordinates as
$`ds^2`$ $`=e^{2\psi _D}(12M/r)dv^2+(12M/r)^1e^{2\psi _D}(e^{2(\gamma _D2\psi _D)}1)dr^2`$ (A.12)
$`+2e^{2\psi _D}dvdr+e^{2(\gamma _D\psi _D)}r^2d\theta ^2+e^{2\psi _D}r^2\mathrm{sin}^2\theta d\varphi ^2.`$ (A.13)
Using condition (A.10) it is not difficult to show that the coefficient of $`dr^2`$ in the metric is regular at $`r=2M`$ .
It is immediately obvious from (A.12) that the $`r=2M`$ surface is a Killing horizon of $`\frac{}{v}`$. However, we can not select a preferred normalization for this vector field since the metric is not asymptotically flat. As in the last sub-section, let $`\mathrm{\Delta }`$ be the Killing horizon and choose $`\mathrm{}\frac{}{v}`$. Then, it is straightforward to verify that $`(\mathrm{\Delta },[\mathrm{}])`$ is a complete, non-rotating isolated horizon. Let us calculate the value of surface gravity for $`\mathrm{}\widehat{}=\frac{}{v}`$. We obtain
$$\kappa \widehat{}=(e^{2\psi _D\gamma _D})\frac{1}{2r}\widehat{}=\frac{1}{2r}$$
(A.14)
where we arrived at the last expression by using (A.10). Again, while spin coefficient $`\mathrm{Re}\left[\mu \right]`$, the Weyl component $`\mathrm{\Psi }_2`$ and the scalar curvature $`{}_{}{}^{2}R`$ of the horizon metric all depend on the distortion function $`\psi _D`$, somewhat surprisingly the surface gravity $`\kappa _{(\mathrm{})}`$ does not.
The natural question is whether the above framework can be extended to obtain distorted Reissner-Nordström solutions. This turns out to be non-trivial because the key equation (A.4) now acquires a source term from the electromagnetic field and this field itself depends non-trivially on $`\psi `$ through the Maxwell equations. At first, the coupled system appears to be hopelessly difficult. However, there exists a prescription for defining a new potential $`\stackrel{~}{\psi }`$ in terms of $`\psi `$ and the electromagnetic field such that $`\stackrel{~}{\psi }`$ satisfies the flat space Laplacian (A.4). Using this method, the known distorted black hole solutions were recently generalized to the charged case . The distorted Reissner-Nordström solution is given by the metric
$`ds^2`$ $`=`$ $`(12M/r+Q^2/r^2)e^{2\psi _D}dt^2+(12M/r+Q^2/r^2)^1e^{2(\gamma _D\psi _D)}dr^2`$ (A.16)
$`+e^{2(\gamma _D\psi _D)}r^2d\theta ^2+e^{2\psi _D}r^2\mathrm{sin}^2\theta d\varphi ^2.`$
The forms of $`\psi _D`$ and $`\gamma _D`$ are now substantially more complicated than in the uncharged case. Nonetheless, it is still possible to show that (A.10) continues to hold. As before this equality implies that the apparent singularity at $`r_H^22Mr_H+Q^2=0`$ is only a coordinate singularity. The surface defined by $`r=r_H`$ is a Killing horizon of $`\frac{}{t}`$. There is once again, no natural way to normalize the Killing field, so we only have an equivalence class $`[\mathrm{}^a]`$ of null normals to the Killing horizon. $`(\mathrm{\Delta },[\mathrm{}])`$ is a non-rotating isolated horizon.
The surface gravity of $`\frac{}{t}`$ is given by
$$\kappa =\frac{1}{2r_H}\left(1\frac{Q^2}{r_H^2}\right).$$
(A.17)
Again, the surface gravity is independent of the distortion of the horizon and has the same dependence on the horizon radius $`R_\mathrm{\Delta }`$ (which turns out to be equal to $`r_H`$) and charge $`Q`$ as in Reissner-Nordström spacetime. Considerations of Section V suggest this peculiar behavior of $`\kappa `$ in all these examples is not accidental but can be ‘explained’ from general Hamiltonian considerations which led us to the first law.
## B The Newman-Penrose formalism
### 1 Notation and Conventions
Let us begin with a summary of the Newman-Penrose formalism (see or for a complete account). Apart from the spacetime signature which we take to be $`(,+,+,+)`$, we will follow the conventions used in . Consider a tetrad of null vectors $`n`$, $`\mathrm{}`$, $`m`$ and $`\overline{m}`$ ($`n`$ and $`\mathrm{}`$ are real while $`m`$ is complex) which satisfy
$$\begin{array}{ccccccccc}\hfill n.\mathrm{}& =& 1\hfill & \hfill n.m& =& 0\hfill & \hfill n.\overline{m}& =& 0\hfill \\ \hfill \mathrm{}.m& =& 0\hfill & \hfill \mathrm{}.\overline{m}& =& 0\hfill & \hfill m.\overline{m}& =& 1.\hfill \end{array}$$
(B.1)
The directional derivatives along the basis vectors are denoted by
$$D=\mathrm{}^a_a\mathrm{\Delta }=n^a_a\delta =m^a_a\overline{\delta }=\overline{m}^a_a.$$
(B.2)
The full the information contained in the connection is expressed in terms of twelve complex scalars called the Newman-Penrose spin coefficients defined as follows:
$$\begin{array}{ccccccccc}\hfill \kappa & =& m^a\mathrm{}^b_b\mathrm{}_a\hfill & \hfill ϵ& =& \frac{1}{2}(\overline{m}^a\mathrm{}^b_bm_an^a\mathrm{}^b_b\mathrm{}_a)\hfill & \hfill \pi & =& \overline{m}^a\mathrm{}^b_bn_a\hfill \\ \hfill \sigma & =& m^am^b_b\mathrm{}_a\hfill & \hfill \beta & =& \frac{1}{2}(\overline{m}^am^b_bm_an^am^b_b\mathrm{}_a)\hfill & \hfill \mu & =& \overline{m}^am^b_bn_a\hfill \\ \hfill \rho & =& m^a\overline{m}^b_b\mathrm{}_a\hfill & \hfill \alpha & =& \frac{1}{2}(\overline{m}^a\overline{m}^b_bm_an^a\overline{m}^b_b\mathrm{}_a)\hfill & \hfill \lambda & =& \overline{m}^a\overline{m}^b_bn_a\hfill \\ \hfill \tau & =& m^an^b_b\mathrm{}_a\hfill & \hfill \gamma & =& \frac{1}{2}(\overline{m}^an^b_bm_an^an^b_b\mathrm{}_a)\hfill & \hfill \nu & =& \overline{m}^an^b_bn_a.\hfill \end{array}$$
(B.3)
It is sometimes more useful to express these definitions in terms of covariant derivatives of the basis vectors:
$$\begin{array}{cccccc}\hfill D\mathrm{}& =& (ϵ+\overline{ϵ})\mathrm{}\overline{\kappa }m\kappa \overline{m}\hfill & \hfill Dn& =& (ϵ+\overline{ϵ})n+\pi m+\overline{\pi }\overline{m}\hfill \\ \hfill \mathrm{\Delta }\mathrm{}& =& (\gamma +\overline{\gamma })\mathrm{}\overline{\tau }m\tau \overline{m}\hfill & \hfill \mathrm{\Delta }n& =& (\gamma +\overline{\gamma })n+\nu m+\overline{\nu }\overline{m}\hfill \\ \hfill \delta \mathrm{}& =& (\overline{\alpha }+\beta )\mathrm{}\overline{\rho }m\sigma \overline{m}\hfill & \hfill \delta n& =& (\overline{\alpha }+\beta )n+\mu m+\overline{\lambda }\overline{m}\hfill \\ \hfill Dm& =& \overline{\pi }\mathrm{}\kappa n+(ϵ\overline{ϵ})m\hfill & \hfill \mathrm{\Delta }m& =& \overline{\nu }\mathrm{}\tau n+(\gamma \overline{\gamma })m\hfill \\ \hfill \delta m& =& \overline{\lambda }\mathrm{}\sigma n+(\beta \overline{\alpha })m\hfill & \hfill \overline{\delta }m& =& \overline{\mu }\mathrm{}\rho n+(\alpha \overline{\beta })m.\hfill \end{array}$$
(B.4)
The ten independent components of the Weyl tensor are expressed in terms of five complex scalars $`\mathrm{\Psi }_0`$, $`\mathrm{\Psi }_1`$, $`\mathrm{\Psi }_2`$, $`\mathrm{\Psi }_3`$ and $`\mathrm{\Psi }_4`$. The ten components of the Ricci tensor are defined in terms of four real and three complex scalars $`\mathrm{\Phi }_{00}`$, $`\mathrm{\Phi }_{11}`$, $`\mathrm{\Phi }_{22}`$, $`\mathrm{\Lambda }`$, $`\mathrm{\Phi }_{10}`$, $`\mathrm{\Phi }_{20}`$ and $`\mathrm{\Phi }_{21}`$ . These scalars are defined as follows:
$$\begin{array}{ccccccccc}\hfill \mathrm{\Psi }_0& =& C_{abcd}\mathrm{}^am^b\mathrm{}^cm^d\hfill & \hfill \mathrm{\Phi }_{01}& =& \frac{1}{2}R_{ab}\mathrm{}^am^b\hfill & \hfill \mathrm{\Phi }_{10}& =& \frac{1}{2}R_{ab}\mathrm{}^a\overline{m}^b\hfill \\ \hfill \mathrm{\Psi }_1& =& C_{abcd}\mathrm{}^am^b\mathrm{}^cn^d\hfill & \hfill \mathrm{\Phi }_{02}& =& \frac{1}{2}R_{ab}m^am^b\hfill & \hfill \mathrm{\Phi }_{20}& =& \frac{1}{2}R_{ab}\overline{m}^a\overline{m}^b\hfill \\ \hfill \mathrm{\Psi }_2& =& C_{abcd}\mathrm{}^am^b\overline{m}^cn^d\hfill & \hfill \mathrm{\Phi }_{21}& =& \frac{1}{2}R_{ab}\overline{m}^an^b\hfill & \hfill \mathrm{\Phi }_{12}& =& \frac{1}{2}R_{ab}m^an^b\hfill \\ \hfill \mathrm{\Psi }_3& =& C_{abcd}\mathrm{}^an^b\overline{m}^cn^d\hfill & \hfill \mathrm{\Phi }_{00}& =& \frac{1}{2}R_{ab}\mathrm{}^a\mathrm{}^b\hfill & \hfill \mathrm{\Phi }_{11}& =& \frac{1}{4}R_{ab}(\mathrm{}^an^b+m^a\overline{m}^b)\hfill \\ \hfill \mathrm{\Psi }_4& =& C_{abcd}\overline{m}^an^b\overline{m}^cn^d\hfill & \hfill \mathrm{\Phi }_{22}& =& \frac{1}{2}R_{ab}n^an^b\hfill & \hfill \mathrm{\Lambda }& =& \frac{R}{24}.\hfill \end{array}$$
(B.5)
The six components of the Electromagnetic-Field $`2`$form $`𝐅_{ab}`$ can be defined in terms of three complex scalars:
$$\varphi _0=\mathrm{}^am^b𝐅_{ab}\varphi _1=\frac{1}{2}(\mathrm{}^an^bm^a\overline{m}^b)𝐅_{ab}\varphi _2=n^a\overline{m}^b𝐅_{ab}.$$
(B.6)
The eight real Maxwell equations $`d𝐅=0`$ and $`d{}_{}{}^{}𝐅=0`$ can be written as a set of four complex equations:
$`D\varphi _1\overline{\delta }\varphi _0`$ $`=`$ $`(\pi 2\alpha )\varphi _0+2\rho \varphi _1\kappa \varphi _2`$ (B.7)
$`D\varphi _2\overline{\delta }\varphi _1`$ $`=`$ $`\lambda \varphi _0+2\pi \varphi _1+(\rho 2ϵ)\varphi _2`$ (B.8)
$`\mathrm{\Delta }\varphi _0\delta \varphi _1`$ $`=`$ $`(2\gamma \mu )\varphi _02\tau \varphi _1+\sigma \varphi _2`$ (B.9)
$`\mathrm{\Delta }\varphi _1\delta \varphi _2`$ $`=`$ $`\nu \varphi _02\mu \varphi _1+(2\beta \tau )\varphi _2.`$ (B.10)
### 2 Boundary Conditions
In this section we describe the isolated horizon boundary conditions in the Newman-Penrose formalism. We will restrict ourselves to Einstein-Maxwell theory with zero cosmological constant.
In a null-tetrad adapted to the null hypersurface $`\mathrm{\Delta }`$, take $`\mathrm{}`$ to be a null normal, $`m`$ and $`\overline{m}`$ tangent to $`\mathrm{\Delta }`$ and $`n`$ transverse to $`\mathrm{\Delta }`$. Since $`\mathrm{}`$ is hypersurface orthogonal and null, it is geodesic. This implies that<sup>\***</sup><sup>\***</sup>\***We will denote the NP spin coefficient $`\kappa `$ by $`\kappa _{NP}`$ to distinguish it from the surface gravity $`\kappa _{(\mathrm{})}`$. $`\kappa _{NP}\widehat{}=0`$ and $`\mathrm{Im}\left[\rho \right]\widehat{}=0`$. Thus
$$D\mathrm{}^b:=\mathrm{}^a_a\mathrm{}^b\widehat{}=(ϵ+\overline{ϵ})\mathrm{}^b.$$
(B.11)
The surface gravity is therefore given by $`\kappa _{(\mathrm{})}=ϵ+\overline{ϵ}`$ and the expansion of $`\mathrm{}`$ is $`\theta _{(\mathrm{})}\widehat{}=\mathrm{Re}\left[\rho \right]`$.
For a non-expanding horizon $`\mathrm{\Delta }`$, the conditions on $`\mathrm{}`$ imply $`\rho \widehat{}=0`$ and the Raychaudhuri equation then implies $`\sigma \widehat{}=0`$ and $`\mathrm{\Phi }_{00}=\frac{1}{2}R_{ab}\mathrm{}^a\mathrm{}^b\widehat{}=0`$. Furthermore, from (II.17) (which is a consequence of the energy condition), it follows that $`\varphi _0\widehat{}=0`$. This leads to the following conditions on the Ricci tensor at the horizon
$$\begin{array}{ccccccccc}\hfill \mathrm{\Phi }_{00}& \widehat{}=& 0\hfill & \hfill \mathrm{\Phi }_{01}& \widehat{}=& 0\hfill & \hfill \mathrm{\Phi }_{10}& \widehat{}=& 0\hfill \\ \hfill \mathrm{\Phi }_{02}& \widehat{}=& 0\hfill & \hfill \mathrm{\Phi }_{20}& \widehat{}=& 0\hfill & \hfill \mathrm{\Phi }_{11}& \widehat{}=& 2G\varphi _1\overline{\varphi }_1.\hfill \end{array}$$
(B.12)
The first Maxwell equation (B.7) gives
$$D\varphi _1\widehat{}=0\mathrm{which}\mathrm{implies}\mathrm{D}\mathrm{\Phi }_{11}\widehat{}=0.$$
(B.13)
Also, as shown in (II.8)
$$\mathrm{\Psi }_0\widehat{}=0\mathrm{and}\mathrm{\Psi }_1\widehat{}=0.$$
(B.14)
The intrinsically defined one-form $`\omega _a`$ defined in (II.4) is given by
$$\omega _a=\kappa _{(\mathrm{})}n_a+(\alpha +\overline{\beta })m_a+(\overline{\alpha }+\beta )\overline{m}_a.$$
(B.15)
It is often convenient to choose the null tetrad such that $`\begin{array}{c}dn\hfill \end{array}=0`$ which implies
$$\mu \widehat{}=\overline{\mu }\mathrm{and}\pi \widehat{}=\alpha +\overline{\beta }.$$
(B.16)
In this case we get a foliation of $`\mathrm{\Delta }`$ spanned by $`m`$ and $`\overline{m}`$. Furthermore, by an appropriate spin transformation, we can choose $`ϵ`$ to be real so that $`ϵ\widehat{}=\overline{ϵ}`$ and thus the foliation is Lie dragged along $`\mathrm{}`$:
$$_{\mathrm{}}m^a=(ϵ\overline{ϵ})m^a\widehat{}=0.$$
(B.17)
The one-form $`\omega `$ now becomes
$$\omega _a=\kappa _{(\mathrm{})}n_a+\pi m_a+\overline{\pi }\overline{m}_a.$$
(B.18)
Let us consider a weakly isolated horizon $`(\mathrm{\Delta },[\mathrm{}])`$. The condition $`_{\mathrm{}}\omega =0`$ is equivalent to requiring
$$_{\mathrm{}}\pi \widehat{}=0\mathrm{and}_{\mathrm{}}\kappa _{(\mathrm{})}\widehat{}=0$$
(B.19)
and as we proved in Section II B, these conditions imply that the surface gravity $`\kappa _{(\mathrm{})}`$ is constant on $`\mathrm{\Delta }`$.
As mentioned in Section II B, a weakly isolated horizon with non-zero surface gravity admits a natural foliation. In the Newman-Penrose framework this foliation can be characterized as follows: It is the unique foliation on each leaf of which the pull-back of the 1-form $`\pi m_a+\overline{\pi }\overline{m}_a`$ is divergence-free. This condition was first introduced by Háj́iček in the context of stationary spacetimes.
Finally, since our boundary conditions require that $`_{\mathrm{}}\kappa _{(\mathrm{})}\widehat{}=0`$, in a sense, a part of the zeroth law is simply assumed. As mentioned in Section II B, we could have used a slightly different set of boundary conditions which make no direct requirement on $`\kappa _{(\mathrm{})}`$ and yet lead to the zeroth law (as well as the results of Sections III \- VII).
Let $`(\mathrm{\Delta },[\mathrm{}])`$ be a non-expanding horizon, equipped with an equivalence class $`[\mathrm{}]`$ of null normals to related to each other by constant positive rescalings. As above introduce a null tetrad where $`\mathrm{}`$ is an element of $`[\mathrm{}]`$, $`m`$ and $`\overline{m}`$ are tangent to the foliation, $`n`$ is curl free and $`ϵ`$ is real. In place of Definition 2, let us assume that $`\mathrm{\Delta }`$ admits a foliation by a family $`S_\mathrm{\Delta }`$ of $`2`$spheres transverse to $`[\mathrm{}]`$ such that the NP spin coefficients in an associated null-tetrad satisfy:
$$_{\mathrm{}}\mu \widehat{}=0\mathrm{and}_{\mathrm{}}\pi \widehat{}=0.$$
(B.20)
These conditions now replace the requirement $`_{\mathrm{}}\omega \widehat{}=0`$ used in the definition of a weakly isolated horizon. We can prove the zeroth law from these conditions as follows. First, consider the definition $`2_{[a}_{b]}\xi _c=R_{abc}^{}{}_{}{}^{d}\xi _d`$ of the Riemann tensor. In the NP formalism, these are written as a set of $`18`$ complex equations known as the ‘field equations’. For our purposes, we need only three of these equations
$`D\alpha \overline{\delta }ϵ`$ $`=`$ $`(\rho +\overline{ϵ}2ϵ)\alpha +\beta \overline{\sigma }\overline{\beta }ϵ\kappa \lambda \overline{\kappa }\gamma +(ϵ+\rho )\pi +\mathrm{\Phi }_{10}`$ (B.21)
$`D\beta \delta ϵ`$ $`=`$ $`(\alpha +\pi )\sigma +(\overline{\rho }\overline{ϵ})\beta (\mu +\gamma )\kappa (\overline{\alpha }\overline{\pi })ϵ+\mathrm{\Psi }_1`$ (B.22)
$`D\mu \delta \pi `$ $`=`$ $`(\overline{\rho }ϵ\overline{ϵ})\mu +\sigma \lambda +(\overline{\pi }\overline{\alpha }+\beta )\pi \nu \kappa +\mathrm{\Psi }_2+2\mathrm{\Lambda }.`$ (B.23)
Adding the first equation to the complex conjugate of the second equation and imposing our boundary conditions gives
$$\delta (ϵ+\overline{ϵ})\widehat{}=\delta \kappa _{(\mathrm{})}\widehat{}=0$$
(B.24)
while the third equation reduces to
$$\mathrm{\Psi }_2\widehat{}=(ϵ+\overline{ϵ})\mu .$$
(B.25)
Equation (B.24) tells us that surface gravity is constant on each leaf of the foliation. It now only remains to show that it is also constant along $`\mathrm{}`$. To show this we turn to the Bianchi identity: $`_{[a}R_{bc]de}=0`$. In the NP formalism, this is written as a set of nine complex and two real equations. We shall need only two of these equations
$`D\mathrm{\Psi }_2\overline{\delta }\mathrm{\Psi }_1+\mathrm{\Delta }\mathrm{\Phi }_{00}\overline{\delta }\mathrm{\Phi }_{01}+2D\mathrm{\Lambda }`$ (B.26)
$`=\lambda \mathrm{\Psi }_0+2(\pi \alpha )\mathrm{\Psi }_1+3\rho \mathrm{\Psi }_22\kappa \mathrm{\Psi }_3+\overline{\sigma }\mathrm{\Phi }_{02}`$ (B.27)
$`+(2\gamma +2\overline{\gamma }\overline{\mu })\mathrm{\Phi }_{00}2(\alpha +\overline{\tau })\mathrm{\Phi }_{01}2\tau \mathrm{\Phi }_{10}+2\rho \mathrm{\Phi }_{11}`$ (B.28)
(B.29)
$`D\mathrm{\Phi }_{11}\delta \mathrm{\Phi }_{10}+\mathrm{\Delta }\mathrm{\Phi }_{00}\overline{\delta }\mathrm{\Phi }_{01}+3D\mathrm{\Lambda }`$ (B.30)
$`=(2\gamma +2\overline{\gamma }\mu \overline{\mu })\mathrm{\Phi }_{00}+(\pi 2\alpha 2\overline{\tau })\mathrm{\Phi }_{01}+(\overline{\pi }2\overline{\alpha }2\tau )\mathrm{\Phi }_{10}`$ (B.31)
$`+2(\rho +\overline{\rho })\mathrm{\Phi }_{11}+\overline{\sigma }\mathrm{\Phi }_{02}+\sigma \mathrm{\Phi }_{20}\overline{\kappa }\mathrm{\Phi }_{12}\kappa \mathrm{\Phi }_{21}.`$ (B.32)
Subtracting these equations, imposing our boundary conditions and using $`\mathrm{\Lambda }=0`$, we get $`D\mathrm{\Psi }_2\widehat{}=0`$. Combining this result with (B.25) gives $`D(ϵ+\overline{ϵ})\widehat{}=0`$. This completes the proof of the zeroth law within the alternate definition of weak isolation. Most of the results of this paper were first obtained using that definition. However, since that notion is tied so heavily to the presence of a foliation, its intrinsic meaning is somewhat obscure. Therefore, it was then replaced by Definition 2 used in the main body of the paper.
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# An Eulerian-Lagrangian Approach to the Navier-Stokes equations.
## 1 Introduction
This work presents an Eulerian-Lagrangian approach to the Navier-Stokes equation. An Eulerian-Lagrangian description of the Euler equations has been used in (, ) for local existence results and constraints on blow-up. Eulerian coordinates (fixed Euclidean coordinates) are natural for both analysis and laboratory experiment. Lagrangian variables have a certain theoretical appeal. In this work I present an approach to the Navier-Stokes equations that is phrased in unbiased Eulerian coordinates, yet describes objects that have Lagrangian significance: particle paths, their dispersion and diffusion. The commutator between Lagrangian and Eulerian derivatives plays an important role in the Navier-Stokes equations: it contributes a singular perturbation to the Euler equations, in addition to the Laplacian. The Navier-Stokes equations are shown to be equivalent to the system
$$\mathrm{\Gamma }v=2\nu Cv$$
where $`C`$ are the coefficients of the commutator between Eulerian and Lagrangian derivatives, and $`\mathrm{\Gamma }`$ is the operator of material derivative and viscous diffusion. The physical pressure is not explicitly present in this formulation. The Eulerian velocity $`u`$ is related to $`v`$ in a non-local fashion, and one may recover the physical pressure dynamically from the evolution of the gradient part of $`v`$. When one sets $`\nu =0`$ the commutator coefficients $`C`$ do not enter the equation, and then $`v`$ is a passive rearrangement of its initial value (, , , ). When $`\nu 0`$ the perturbation involves the curvature of the particle paths, and the gradients of $`v`$: a singular perturbation. Fortunately, the coefficients $`C`$ start from zero, and, as long as they remain small $`v`$ does not grow too much.
A different but not unrelated approach (, () is based on a variable $`w`$ that has the same curl as the Eulerian velocity $`u`$. The velocity is recovered then from $`w`$ by projection on divergence-free functions. The evolution equation for $`w`$
$$\mathrm{\Gamma }w+\left(u\right)^{}w=0,$$
conserves local helicity and circulation (when $`\nu =0`$). We will refer to this equation informally as “the cotangent equation” because it is the equation obeyed by the Eulerian gradient of any scalar $`\varphi `$ that solves $`\mathrm{\Gamma }\varphi =0`$. The variable $`w`$ is related to $`v`$:
$$w=(A)^{}v$$
where $`A`$ is the “back-to-labels” maps that corresponds when $`\nu =0`$ to the inverse of the Lagrangian path map. $`A(x,t)`$ is an active vector obeying
$$\mathrm{\Gamma }A=0,$$
$`A(x,0)=x`$. Both $`v`$ and $`A`$ have a Lagrangian meaning when $`\nu =0`$, but the dynamical development of $`w`$ is the product of two processes, the growth of the deformation tensor (given by the evolution of $`A`$) and the rearrangement of a fixed function, given by the evolution of $`v`$. In the presence of viscosity, $`v`$’s evolution is not by rearrangement only. It is therefore useful to study separately the growth of $`A`$ and the shift of $`v`$.
Recently certain model equations (alpha-models) have been proposed (, ) as modifications of the Euler and Navier-Stokes equations. They can be obtained in the context described above simply by smoothing $`u`$ in the cotangent equation. Smoothing means that one replaces the linear zero-order nonlocal operator $`u=𝐏w`$ that relates $`u`$ to $`w`$ in the cotangent equation by a smoothing operator, $`u=K_\alpha 𝐏w`$. When $`\nu =0`$ the models have a Kelvin circulation theorem. They cannot be models for Eulerian averaged equations, the equations describing mean flow in turbulence theory. These Reynolds equations are not conservative in the limit of zero viscosity: the fluctuations introduce additional stresses, the so-called Reynolds stresses, that preclude conservation of circulation along the mean flow. This is a simple yet fundamental objection to the identification of alpha-models with Reynolds equations. The alpha-models might be models of Lagrangian averaged equations, an entirely different concept. (Operationally one may think of Eulerian averaging as long time averages at fixed position, and of Lagrangian averaging as long time averages at fixed initial label. However, a fixed initial label has no obvious physical meaning when one deals with ensembles of flows.)
In this paper we consider the Navier-Stokes equations and obtain rigorous bounds for the particle paths and for the virtual velocity $`v`$. The main bounds concern the Lagrangian displacement, its first and second spatial derivatives, are obtained under general conditions and require no assumptions. Higher derivatives can be bounded also under certain natural quantitative smoothness assumptions.
Some of our bounds can be interpreted as a connection to the Richardson pair dispersion law, one of the empirical laws of fully developed turbulence (), that is consistent with the Kolmogorov two-thirds law (). The pair dispersion law states that the separation $`\delta `$ between fluid particles obeys
$$|\delta |^2ϵt^3$$
where $`ϵ`$ is the rate of dissipation of energy and $`t`$ is time. This is supposed to hold in an inertial range, for times $`t`$ that are neither too small (when the separation is ballistic) nor too large, when viscous and boundary effects are important. The law can be guessed by dimensional analysis by requiring the answer to depend solely on time and $`ϵ`$. Precise laboratory Lagrangian experiments have recently begun to be capable of addressing Lagrangian quantities with preliminary results that seem to be consistent with the Richardson law in some ranges (). If one considers the problem of estimating the pair dispersion mathematically one is faced with the difficulty that the prediction seems to require both non-Lipschitz, Hölder continuous velocities and Lagrangian particle paths. Our approach allows a rigorous formulation and an upper bound
$$|\delta |^23|\delta _0|^2+24E_0t^2+6ϵ_Bt^3,$$
that includes a reference to an initial displacement $`\delta _0`$, an initial kinetic energy $`E_0`$ and a rigorous upper bound on $`ϵ`$. The prefactors are probably not optimal. The conditions under which such a bound can be obtained are quite general, and there are no assumptions. In this paper we chose the case of periodic boundary conditions, and body forces that have a characteristic length scale that remains finite as the size of the periodic box is allowed to diverge. The bound $`ϵ_B`$ does not depend on the size of box. In many physically realistic situations one injects energy at the boundary; in that case one can find $`ϵ_B`$ independently of viscosity (), without any assumptions.
## 2 Velocity and displacement
The Eulerian velocity $`u(x,t)`$ has three components $`u^i,i=1,\mathrm{\hspace{0.17em}2},3`$ and is a function of three Eulerian space coordinates $`x`$ and time $`t`$. We decompose the Eulerian velocity $`u(x,t)`$:
$$u^i(x,t)=\frac{A^m(x,t)}{x_i}v_m(x,t)\frac{n(x,t)}{x_i}.$$
(1)
Repeated indices are summed. There are three objects that appear in this formula. The first one, $`A(x,t)`$, has a Lagrangian interpretation. In the absence of viscosity, $`A`$ is the “back-to-labels” map, the inverse of the particle trajectory map $`ax=X(a,t)`$. The vector
$$\mathrm{}(x,t)=A(x,t)x$$
(2)
will be called the “Eulerian-Lagrangian displacement vector”, or simply “displacement”. $`\mathrm{}`$ joins the current Eulerian position $`x`$ to the original Lagrangian position $`a=A(x,t)`$. $`A(x,t)`$ and $`\mathrm{}(x,t)`$ have dimensions of length, $`A`$ is non-dimensional. A pair of points, $`a=A(x,t)`$, $`b=A(y,t)`$ situated at time $`t=0`$ at distance $`\delta _0=|ab|`$ become separated by $`\delta _t=|xy|`$ at time $`t`$. From the triangle inequality it follows that
$$(\delta _t)^23|\mathrm{}(x,t)|^2+3|\mathrm{}(y,t)|^2+3(\delta _0)^2.$$
(3)
The displacement can be used in this manner to bound pair dispersion.
The second object in (1), $`v(x,t)`$, has dimensions of velocity and, in the absence of viscosity, is just the initial velocity composed with the back-to-labels map (, , , ). We call $`v`$ the “virtual velocity”. Its evolution marks the difference between the Euler and Navier-Stokes equations most clearly. The third object in (1) is a scalar function $`n(x,t)`$ that will be referred to as “the Eulerian-Lagrangian potential”. It plays a mathematical role akin to that played by the physical pressure but has dimensions of length squared per time, like the kinematic viscosity. If $`A(x,t)`$ is known, then there are four functions entering the decomposition of $`u`$, three $`v`$-s and one $`n`$. If the velocity is divergence-free
$$u=0,$$
then there is one relationship between the four unknown functions.
## 3 Eulerian-Lagrangian derivatives and commutators
When one considers the map $`xA(x,t)`$ as a change of variables one can pull back the Lagrangian differentiation with respect to particle position and write it in Eulerian coordinates using the chain rule. Let us call this pull-back of Lagrangian derivatives the Eulerian-Lagrangian derivative,
$$_A=Q^{}_E.$$
(4)
Here
$$Q(x,t)=\left(A(x,t)\right)^1,$$
(5)
and the notation $`Q^{}`$ refers to the transpose of the matrix $`Q`$. The expression of $`_A`$ on components is
$$_A^i=Q_{ji}_j$$
(6)
where we wrote $`_i`$ for differentiation in the $`i`$-th Eulerian Cartesian coordinate direction,
$$_i=_E^i.$$
The Eulerian spatial derivatives can be expressed in terms of the Eulerian-Lagrangian derivatives via
$$_E^i=\left(_iA_m\right)_A^m$$
(7)
The commutation relations
$$[_E^i,_E^k]=0,[_A^i,_A^k]=0$$
hold. The commutators between Eulerian-Lagrangian and Eulerian derivatives do not vanish, in general:
$$[_A^i,_E^k]=C_{m,k;i}_A^m.$$
(8)
The coefficients $`C_{m,k;i}`$ are given by
$$C_{m,k;i}=\{_A^i(_k\mathrm{}_m)\}.$$
(9)
Note that
$$C_{m,k;i}=Q_{ji}_i_kA_m=_A^i(_E^kA_m)=[_A^i,_E^k]A_m.$$
These commutator coefficients play an important role in dynamics.
## 4 The evolution of A
We associate to a given divergence-free velocity $`u(x,t)`$ the operator
$$_t+u\nu \mathrm{\Delta }=\mathrm{\Gamma }_\nu (u,).$$
(10)
We write $`_t`$ for time derivative. We write $`\mathrm{\Gamma }`$ for $`\mathrm{\Gamma }_\nu (u,)`$ when the $`u`$ we use is clear from the context. The coefficient $`\nu >0`$ is the kinematic viscosity of the fluid. When applied to a vector or a matrix, $`\mathrm{\Gamma }`$ acts as a diagonal operator, i.e. on each component separately. The operator $`\mathrm{\Gamma }`$ obeys a maximum principle: If a function $`q`$ solves
$$\mathrm{\Gamma }q=S$$
and the function $`q`$ has homogeneous Dirichlet or periodic boundary conditions, then the sup-norm $`q_{L^{\mathrm{}}(dx)}`$ satisfies
$$q(,t)_{L^{\mathrm{}}(dx)}q(,t_0)_{L^{\mathrm{}}(dx)}+_{t_0}^tS(,s)_{L^{\mathrm{}}(dx)}𝑑s$$
for any $`t_0t`$. The operator $`\mathrm{\Gamma }_\nu (u,)`$ is not a derivation (that means an operator that satisfies the product rule); $`\mathrm{\Gamma }`$ satisfies a product rule that is similar to that of a derivation:
$$\mathrm{\Gamma }(fg)=(\mathrm{\Gamma }f)g+f(\mathrm{\Gamma }g)2\nu (_kf)(_kg).$$
(11)
We require the back-to-labels map $`A`$ to obey
$$\mathrm{\Gamma }A=0.$$
(12)
By (12) we express therefore the advection and diffusion of $`A`$. We will use sometimes the equation obeyed by $`\mathrm{}`$
$$\left(_t+u\nu \mathrm{\Delta }\right)\mathrm{}+u=0$$
(13)
which is obviously equivalent to (12). We will discuss periodic boundary conditions
$$\mathrm{}(x+Le_j,t)=\mathrm{}(x,t),$$
where $`e_j`$ is the unit vector in the $`j`$-th direction. Some of our inequalities will hold also for the physical boundary condition that require $`\mathrm{}(x,t)=0`$ at the boundary.
It is important to note that the initial data for the displacement is zero:
$$\mathrm{}(x,0)=0.$$
(14)
The matrix $`A(x,t)`$ is invertible as long as the evolution is smooth. This is obvious when $`\nu =0`$ because the determinant of this matrix equals $`1`$ for all time, but in the viscous case the statement needs proof. We differentiate (12) in order to obtain the equation obeyed by $`A`$
$$\mathrm{\Gamma }(A)+(A)(u)=0.$$
(15)
The product $`(A)(u)`$ is matrix product in the order indicated. We consider
$$\mathrm{\Gamma }Q=(u)Q+2\nu Q_k(A)_kQ.$$
(16)
It is clear that the solutions of both (15) and (16) are smooth as long as the advecting velocity $`u`$ is sufficiently smooth. It is easy to verify using (11) that the matrix $`Z=(A)QI`$ obeys the equation
$$\mathrm{\Gamma }Z=2\nu Z_k(A)_kQ$$
with initial datum $`Z(x,0)=0`$. Thus, as long as $`u`$ is smooth, $`Z(x,t)=0`$ and it follows that the solution $`Q`$ of (16) is the inverse of $`A`$.
The commutator coefficients $`C_{m,k;i}`$ enter the important commutation relation between the Eulerian-Lagrangian label derivative and $`\mathrm{\Gamma }`$:
$$[\mathrm{\Gamma },_A^i]=2\nu C_{m,k;i}_E^k_A^m$$
(17)
The proof of this formula can be found in Appendix B.
The evolution of the coefficients $`C_{m,k;i}`$ defined in (9) can be computed using (15) and (17):
$$\mathrm{\Gamma }\left(C_{m,k;i}\right)=(_lA_m)_A^i(_k(u_l))$$
$$(_k(u_l))C_{m,l;i}+2\nu C_{j,l;i}_l\left(C_{m,k;j}\right).$$
(18)
The calculation leading to (18) is presented in Appendix B.
## 5 The evolution of v
We require the virtual velocity to obey
$$\mathrm{\Gamma }v=2\nu Cv+Q^{}f.$$
(19)
This equation is, on components
$$\mathrm{\Gamma }_\nu (u,)v_i=2\nu C_{m,k;i}_kv_m+Q_{ji}f_j.$$
(20)
The vector $`f=f(x,t)`$ represents the body forces. The boundary conditions are periodic
$$v(x+Le_j,t)=v(x,t)$$
and the initial data are, for instance
$$v(x,0)=u_0(x).$$
(21)
The reason for requiring the equation (19) is
###### Proposition 1.
Assume that $`u`$ is given by the expression (1) above and that the displacement $`\mathrm{}`$ and the virtual velocity $`v`$ obey the equations (13) and respectively (19). Then the velocity $`u`$ satisfies the Navier-Stokes equation
$$_tu+uu\nu \mathrm{\Delta }u+p=f$$
with pressure $`p`$ determined from the Eulerian-Lagrangian potential by
$$\mathrm{\Gamma }_\nu (u,)n+\frac{|u|^2}{2}+c=p$$
where $`c`$ is a free constant.
Proof. We denote for convenience
$$D_t=_t+u.$$
(22)
We apply $`D_t`$ to the velocity representation (1) and use the commutation relation
$$[D_t,_k]g=(u)^{}g.$$
(23)
We obtain
$$D_t(u^i)=\left(_i(D_tA^m)\right)v_m+(_iA^m)D_tv_m_i\left(\frac{|u|^2}{2}+D_tn\right).$$
We substitute the equations for $`A`$ (13) and for $`v`$ (19):
$$D_t(u^i)=_i\left(\frac{|u|^2}{2}+D_tn\right)+\left(_i(\nu \mathrm{\Delta }A^m)\right)v_m+$$
$$(_iA^m)\left\{\nu \mathrm{\Delta }v_m+Q_{mj}^{}\left(2\nu _k(\mathrm{})_{jl}^{}_kv_l+f_j\right)\right\}.$$
Now we use the facts that
$$(_iA^m)Q_{mj}^{}=\delta _{ij}$$
(Kronecker’s delta), and
$$_k(\mathrm{})_{il}^{}=_k(A)_{il}^{}=_k(_iA^l)$$
to deduce
$$D_t(u^i)=_i\left(\frac{|u|^2}{2}+D_tn\right)+f_i$$
$$+\nu (\mathrm{\Delta }_iA^m)v_m+\nu (_iA^m)\mathrm{\Delta }v_m+2\nu _k(_iA^l)_kv_l$$
and so, changing the dummy summation index $`l`$ to $`m`$ in the last expression
$$D_t(u^i)=_i\left(\frac{|u|^2}{2}+D_tn\right)+\nu \mathrm{\Delta }((_iA^m)v_m)+f_i.$$
Using (1) we obtained
$$D_t(u^i)=\nu \mathrm{\Delta }u_i_i\left(\frac{|u|^2}{2}\nu \mathrm{\Delta }n+D_tn\right)+f_i$$
and that concludes the proof.
Observation The incompressibility of velocity has not yet been used. This is why no restriction on the potential $`n(x,t)`$ was needed. The incompressibility
$$u=0$$
(24)
can be imposed in two ways. The first approach is static: one considers the ansatz (1) and one requires that $`n`$ maintains the incompressibility at each instance of time. This results in the equation
$$\mathrm{\Delta }n=(A)^{}v).$$
(25)
In this way $`n`$ is computed from $`A`$ in a time independent manner and the basic formula (1) can be understood as
$$u=𝐏\left((A)^{}v\right)$$
(26)
where $`𝐏`$ is the Leray-Hodge projector on divergence-free functions. The second approach is dynamic: one computes the physical Navier-Stokes pressure
$$p=R_iR_j(u^iu^j)+c$$
(27)
where $`c`$ is a free constant and $`R_i=(\mathrm{\Delta })^{\frac{1}{2}}_i`$ is the Riesz transform for periodic boundary conditions. The formula for $`p`$ follows by taking the divergence of the Navier-Stokes equation and using (24). Substituting (27) in the expression for the pressure in Proposition 1 one obtains the evolution equation
$$\mathrm{\Gamma }n=R_iR_j(u^iu^j)\frac{|u|^2}{2}+c$$
(28)
for $`n`$. Incompressibility can be enforced either by solving at each time the static equation (25) or by evolving $`n`$ according to (28).
###### Proposition 2.
Let $`u`$ be given by (1) and assume that the displacement solves (13) and that the virtual velocity solves (19). Assume in addition that the potential obeys (25) (respectively (28)). Then $`u`$ obeys the incompressible Navier-Stokes equations,
$$_tu+uu\nu \mathrm{\Delta }u+p=f,u=0,$$
the pressure $`p`$ satisfies (27) and the potential obeys also (28) (respectively (25)).
The same results hold for the case of the whole $`𝐑^3`$ with boundary conditions requiring $`u`$ and $`\mathrm{}`$ to vanish at infinity. In the presence of boundaries, if the boundary conditions for $`u`$ are homogeneous Dirichlet ($`u=0`$) then the boundary conditions for $`v`$ are Dirichlet, but not homogeneous. In that case one needs to solve either one of the equations (25),(28) for $`n`$ (with Dirichlet or other physical boundary condition) and the $`v`$ equation (19) with
$$v=_An$$
at the boundary.
###### Proposition 3.
Let $`u`$ be an arbitrary spatially periodic smooth function and assume that a displacement $`\mathrm{}`$ solves the equation (13) and a virtual velocity $`v`$ obeys the equation (19) with periodic boundary conditions and with $`C`$ computed using $`A=x+\mathrm{}`$. Then $`w`$ defined by
$$w_i=(_iA^m)v_m$$
(29)
obeys the cotangent equation
$$\mathrm{\Gamma }w+(u)^{}w=f.$$
(30)
Proof. The proof is a straightforward calculation. One uses (11) to write
$$\mathrm{\Gamma }w_i=(_iA^m)\mathrm{\Gamma }v_m+v_m\mathrm{\Gamma }(_iA^m)2\nu (_k_iA^m)_kv_m.$$
The equation (20) is used for the first term and the equation (15) for the second term. One obtains
$$\mathrm{\Gamma }w_i=f_i(_iu_j)w_j+2\nu \left\{(_iA^m)C_{r,q;m}_qv_r(_k_iA^m)_kv_m\right\}.$$
The proof ends by showing that the term in braces vanishes because of the identity
$$(_iA^m)C_{r,q;m}=_q_iA^r.$$
An approach to the Euler equations based entirely on a variable $`w`$ (, () is well-known. The function $`w`$ has the same curl as $`u`$, $`\omega =\times u=\times w`$. In the case of zero viscosity and no forcing, the local helicity $`w\omega `$ is conserved $`D_t(w\omega )=0`$; this is easily checked using the fact that the vorticity obeys the “tangent” equation $`D_t\omega =(u)\omega `$ and the inviscid, unforced form of (30). The same proof verifies the Kelvin circulation theorem
$$\frac{d}{dt}_{\gamma (t)}w𝑑X=0$$
on loops $`\gamma (t)`$ advected by the flow of $`u`$. Although obviously related, the two variables $`v`$ and $`w`$ have very different analytical merits. While the growth of $`w`$ is difficult to control, in the inviscid case $`v`$ does not grow at all, and in the viscous case its growth is determined by the magnitude of $`C`$ which starts from zero. This is why we emphasize $`v`$ as the primary variable and consider $`w`$ a derived variable.
## 6 Gauge Invariance
Consider a scalar function $`\varphi `$. If one transforms $`v\stackrel{~}{v}=v+_A\varphi `$ and $`n\stackrel{~}{n}=n+\varphi `$ then $`u`$ remains unchanged in (1): $`uu`$. The requirement that $`_Eu=0`$ does not specify this arbitrary $`\varphi `$.
Assume now that the scalar $`\varphi `$ is advected passively by $`u`$ and diffuses with diffusivity $`\nu `$:
$$\mathrm{\Gamma }\varphi =0.$$
Then, in view of (17), if $`v`$ solves (19) then
$$\stackrel{~}{v}=v+_A\varphi $$
also solves (19). If $`n`$ solves (28) then
$$\stackrel{~}{n}=n+\varphi $$
also solves (28). If $`w`$ solves the equation (30) then
$$\stackrel{~}{w}=w+_E\varphi $$
also solves (30). The vector fields obtained by taking the Eulerian gradient of passive scalars are homogeneous solutions of (30). The vector fields obtained by taking the Eulerian-Lagrangian gradient of passive scalars, $`_A\varphi `$ are homogeneous solutions of (19). This can be used to show that if one chooses an initial datum for $`v`$ that differs from $`u_0`$ by the gradient of an arbitrary function $`\varphi _0`$ there is no change in the evolution of $`u`$.
###### Proposition 4.
Let each of two functions $`v_j`$, $`j=1,2`$ solve the system
$$\mathrm{\Gamma }(u_j,)v_j=2\nu C_jv_j+Q_j^{}f$$
with periodic boundary conditions, coupled with
$$\mathrm{\Gamma }(u_j,)A_j=0$$
with periodic boundary conditions for $`\mathrm{}_j=A_jx`$. Assume that the initial data for $`A_j`$ are the same, $`\mathrm{}_j(x,0)=0`$. Assume that each velocity is determined from its corresponding virtual velocity by the rule
$$u_j=𝐏\left((A_j)^{}v_j\right).$$
Assume, moreover, that at time $`t=0`$ the virtual velocities differ by a gradient
$$𝐏v_1=𝐏v_2=u_0.$$
Then, as long as one of the solutions $`v_j`$ is smooth one has
$$u_1(x,t)=u_2(x,t),A_1(x,t)=A_2(x,t)$$
The same kind of result can be proved for (30) using the Eulerian gauge invariance.
## 7 K-bounds
We are going to describe here bounds that are based solely on the kinetic energy balance in the Navier-Stokes equation ( () and references therein). These are very important, as they are the only unconditional bounds that are known for arbitrary time intervals. We call them kinetic energy bounds or in short, K-bounds. We start with the most important, the energy balance itself. From the Navier-Stokes equation one obtains the bound
$$|u(x,t)|^2𝑑x+\nu _{t_0}^t|u(x,s)|^2𝑑x𝑑sK_0$$
(31)
with
$$K_0=\mathrm{min}\{k_0;k_1\}$$
(32)
where
$$k_0=2|u(x,t_0)|^2𝑑x+3(tt_0)_{t_0}^t|f(x,s)|^2𝑑x𝑑s$$
(33)
and
$$k_1=|u(x,t_0)|^2𝑑x+\frac{1}{\nu }_{t_0}^t|\mathrm{\Delta }^{\frac{1}{2}}f(x,s)|^2𝑑x𝑑s$$
(34)
Note that we have not normalized the volume of the domain. The prefactors are not optimal. The energy balance holds for all solutions of the Navier-Stokes equations. We took an arbitrary starting time $`t_0`$. The bound $`K_0`$ is a nondecreasing function of $`tt_0`$. We will use this fact tacitly below. In order to give a physical interpretation to this general bound it is useful to denote by
$$ϵ(s)=\nu L^3|u(x,s)|^2𝑑x$$
the volume average of the instantaneous energy dissipation rate, by
$$E(t)=\frac{1}{2L^3}|u(x,t)|^2𝑑x$$
the volume average of the kinetic energy; for any time dependent function $`g(s)`$, we write
$$g()_t=\frac{1}{tt_0}_{t_0}^tg(s)𝑑s$$
for the time average. We also write
$$F^2=L^3|f(x,)|^2𝑑x_t,$$
$$G^2=L^3|\mathrm{\Delta }^{\frac{1}{2}}f(x,)|^2𝑑x_t$$
and define the forcing length scale by
$$L_f^2=\frac{G^2}{F^2}.$$
Then (31) implies
$$2E(t)+(tt_0)ϵ()_t4E(t_0)+(tt_0)F^2\mathrm{min}\{\frac{L_f^2}{\nu };3(tt_0)\}.$$
(35)
After a long enough time
$$tt_0\frac{L_f^2}{3\nu },$$
the kinetic energy grows at most linearly in time
$$E(t)2E(t_0)+F^2\left(\frac{(tt_0)L_f^2}{\nu }\right).$$
The long time for the average dissipation rate is bounded
$$lim\underset{t\mathrm{}}{sup}ϵ()_t\frac{F^2L_f^2}{\nu }=ϵ_B.$$
(36)
These bounds are uniform in the size $`L`$ of the period which we assume to be much larger than $`L_f`$. If the size of the period is allowed to enter the calculations then the kinetic energy is bounded by
$$E(t)L^2\frac{L_f^2F_{}^2}{\nu ^2}+\left(E(t_0)L^2\frac{L_f^2F_{}^2}{\nu ^2}\right)e^{\frac{\nu (tt_0)}{L^2}}$$
where
$$L_f^2F_{}^2=\underset{t}{sup}L^3|(\mathrm{\Delta })^{\frac{1}{2}}f(x,t)|^2𝑑x.$$
This means that for much longer times
$$tt_0\frac{L^2}{\nu }$$
the kinetic energy saturates to a value that depends on the large scale. But the bound (35) that is independent of $`L`$ is always valid; it can be written in terms of
$$B=4E(t_0)+(tt_0)ϵ_B$$
(37)
as
$$E(t)+(tt_0)ϵ_tB.$$
(38)
A useful K-bound is
$$_{t_0}^tu(,s)_{L^{\mathrm{}}(dx)}𝑑sK_{\mathrm{}}$$
(39)
The constant $`K_{\mathrm{}}`$ has dimensions of length and depends on the initial kinetic energy, viscosity, body forces and time. The bound follows by interpolation from () and is derived in Appendix A together with the formula
$$K_{\mathrm{}}=C\{\frac{K_0}{\nu ^2}+\sqrt{\nu (tt_0)}+\frac{tt_0}{\nu ^2}_{t_0}^tf(,s)_{L^2}^2ds\}.$$
(40)
The displacement $`\mathrm{}`$ satisfies certain K-bounds that follow from the bounds above and (13). We mention here
$$\mathrm{}(,t)_{L^{\mathrm{}}(dx)}_{t_0}^tu(,s)_{L^{\mathrm{}}(dx)}𝑑sK_{\mathrm{}},$$
(41)
The inequality (41) follows from (13) by multiplying with $`\mathrm{}|\mathrm{}|^{2(m1)}`$, integrating,
$$\frac{1}{2m}\frac{d}{dt}|\mathrm{}(x,t)|^{2m}𝑑x+\nu |\mathrm{}(x,t)|^2|\mathrm{}(x,t)|^{2(m1)}𝑑x+$$
$$+\nu \frac{m1}{2}||\mathrm{}(x,t)|^2|^2|\mathrm{}(x,t)|^{2(m2)}𝑑x+$$
$$+u(x,t)\mathrm{}(x,t)|\mathrm{}(x,t)|^{2(m1)}𝑑x0,$$
(42)
and then ignoring the viscous terms, using Hölder’s inequality in the last term, multiplying by $`m`$, taking the $`m`$-th root, integrating in time and then letting $`m\mathrm{}`$.
The case $`m=1`$ gives
$$\frac{d}{2dt}|\mathrm{}(x,t)|^2𝑑x+\nu |\mathrm{}(x,t)|^2𝑑x\sqrt{K_0}\sqrt{|\mathrm{}(x,t)|^2𝑑x}$$
and consequently, we obtain by integration from $`t_0=0`$
$$\sqrt{|\mathrm{}(x,t)|^2𝑑x}t\sqrt{K_0},$$
(43)
and then, using (43) we deduce the inequality
$$_0^t|\mathrm{}(x,s)|^2𝑑x𝑑s\frac{K_0t^2}{2\nu }.$$
(44)
Now we multiply (13) by $`\mathrm{\Delta }\mathrm{}`$, integrate by parts, use Schwartz’s inequality to write
$$\frac{d}{dt}|\mathrm{}(x,t)|^2𝑑x+\nu |\mathrm{\Delta }\mathrm{}(x,t)|^2𝑑x\sqrt{|u(x,t)|^2𝑑x}\sqrt{|\mathrm{}(x,t)|^2𝑑x}$$
$$2\text{T}race\left\{(\mathrm{}(x,t))(u(x,t))(\mathrm{}(x,t))^{}\right\}𝑑x$$
and then use the elementary inequality
$$\left(|\mathrm{}(x,t)|^4𝑑x\right)^{\frac{1}{2}}C\mathrm{}(,t)_L^{\mathrm{}}\left(|\mathrm{\Delta }\mathrm{}(x,t)|^2𝑑x\right)^{\frac{1}{2}},$$
in conjunction with the Hölder inequality and (41) to deduce
$$\frac{d}{dt}|\mathrm{}(x,t)|^2𝑑x+\nu |\mathrm{\Delta }\mathrm{}(x,t)|^2𝑑x$$
$$\sqrt{|\mathrm{}(x,t)|^2𝑑x}\sqrt{|u(x,t)|^2𝑑x}+C\frac{K_{\mathrm{}}^2}{\nu }|u(x,t)|^2𝑑x.$$
We obtain, after integration and use of (31, 44)
$$|\mathrm{}(x,t)|^2𝑑x+\nu _0^t|\mathrm{\Delta }\mathrm{}(x,s)|^2𝑑x𝑑sC\left(\frac{K_0t}{\nu }+\frac{K_{\mathrm{}}^2K_0}{\nu ^2}\right).$$
(45)
Recalling the bound (37, 38) on kinetic energy we have:
###### Theorem 1.
Assume that the vector valued function $`\mathrm{}`$ obeys (13) and assume that the velocity $`u(x,t)`$ is a solution of the Navier-Stokes equations (or, more generally, that it is a divergence-free periodic function that satisfies the bounds (31) and (39)). Then $`\mathrm{}`$ satisfies the inequality (41) together with
$$\frac{1}{L^3}|\mathrm{}(x,t)|^2𝑑x(4E_0+tϵ_B)t^2,$$
(46)
$$\frac{1}{L^3t}_0^t|\mathrm{}(x,s)|^2𝑑x𝑑s\frac{Bt}{2\nu },$$
(47)
and
$$|\mathrm{}(x,t)|^2\frac{dx}{L^3}+\nu _0^t|\mathrm{\Delta }\mathrm{}(x,s)|^2\frac{dx}{L^3}𝑑sC\left(\frac{Bt}{\nu }+\frac{K_{\mathrm{}}^2B}{\nu ^2}\right).$$
(48)
In these inequalities
$$E_0=\frac{1}{2L^3}|u(x,0)|^2𝑑x,$$
$$B=4E_0+tϵ_B$$
and $`ϵ_B`$ is given in (36).
Let us consider the pair dispersion
$$\delta _t^2=L^6_{\{(x,y);|A(x,t)A(y,t)|\delta _0\}}|xy|^2𝑑x𝑑y.$$
(49)
Using the triangle inequality (3) in (46) we obtain
###### Theorem 2.
Consider periodic solutions of the Navier-Stokes equation with large period $`L`$, and assume that the body forces have $`L_f`$ finite. Then the pair dispersion obeys
$$\delta _t^23\delta _0^2+24tE_0t^2+6ϵ_Bt^3.$$
(50)
Comment Use of the ODE $`\frac{dX}{dt}=u(X,t)`$ requires information about the gradient $`A`$ and produces worse bounds.
## 8 $`ϵ`$-bounds
This section is devoted to bounds on higher order derivatives of $`\mathrm{}`$. These bounds require assumptions. We are going to apply the Laplacian to (13), multiply by $`\mathrm{\Delta }\mathrm{}`$ and integrate. We obtain
$$\frac{1}{2}\frac{d}{dt}|\mathrm{\Delta }\mathrm{}(x,t)|^2𝑑x+\nu |\mathrm{\Delta }\mathrm{}(x,t)|^2𝑑x=$$
$$_ku(x,t)_k\mathrm{\Delta }\mathrm{}(x,t)dx+I$$
(51)
where
$$I=_k(u(x,t)\mathrm{}(x,t))_k\mathrm{\Delta }\mathrm{}(x,t)dx.$$
Now
$$I=(_ku)\mathrm{}(x,t)_k\mathrm{\Delta }\mathrm{}(x,t)dx+II$$
where
$$II=u(x,t)(_k\mathrm{}(x,t))\mathrm{\Delta }_k\mathrm{}(x,t)dx$$
and, integrating by parts
$$II=_lu(x,t)(_k\mathrm{}(x,t))_l_k\mathrm{}(x,t)dx$$
and then again
$$II=_lu(x,t)_l_k\mathrm{}(x,t)(_k\mathrm{}(x,t))dx$$
and so
$$I=_lu_i(x,t)_k\mathrm{}_j(x,t)\left\{_i_k+\delta _{ik}\mathrm{\Delta }\right\}_l\mathrm{}_j(x,t)dx$$
Putting things together we get
$$|I|C\mathrm{}(,t)_L^{\mathrm{}}u(,t)_{L^2}\mathrm{\Delta }\mathrm{}(,t)_{L^2}$$
Thus
$$\frac{1}{2}\frac{d}{dt}|\mathrm{\Delta }\mathrm{}(x,t)|^2𝑑x+\nu |\mathrm{\Delta }\mathrm{}(x,t)|^2𝑑x$$
$$\frac{C}{\nu }|u(x,t)|^2𝑑x+C\mathrm{}(,t)_L^{\mathrm{}}u(,t)_{L^2}\mathrm{\Delta }\mathrm{}(,t)_{L^2}$$
(52)
Now we use an interpolation inequality that is valid for periodic functions with zero mean and implies that
$$\mathrm{}_L^{\mathrm{}}c\mathrm{\Delta }\mathrm{}_{L^2}^{\frac{1}{2}}\mathrm{\Delta }\mathrm{}_{L^2}^{\frac{1}{2}}$$
Using this inequality we obtain
$$\frac{d}{dt}|\mathrm{\Delta }\mathrm{}(x,t)|^2𝑑x+\nu |\mathrm{\Delta }\mathrm{}(x,t)|^2𝑑x$$
$$\frac{C}{\nu }|u(x,t)|^2𝑑x+C\nu ^3u(,t)_{L^2}^4\mathrm{\Delta }\mathrm{}(,t)_{L^2}^2$$
(53)
Therefore we deduce
$$|\mathrm{\Delta }\mathrm{}(x,t)|^2\frac{dx}{L^3}c\frac{B}{\nu ^2}\text{exp}\left\{\frac{cL^6}{\nu ^5}_0^tϵ^2(s)𝑑s\right\}$$
(54)
where
$$ϵ(s)=\nu L^3|u(x,s)|^2𝑑x$$
(55)
is the instantaneous energy dissipation.
###### Proposition 5.
If $`\mathrm{}`$ solves (13) with periodic boundary conditions on a time interval $`t[0,T]`$ and if the integral
$$_0^Tϵ^2(s)𝑑s$$
is finite, then
$$|\mathrm{\Delta }\mathrm{}(x,t)|^2\frac{dx}{L^3}+\nu _0^t|\mathrm{\Delta }\mathrm{}(x,s)|^2\frac{dx}{L^3}c\frac{B}{\nu ^2}\text{exp}\left\{\frac{cL^6}{\nu ^5}_0^tϵ^2(s)𝑑s\right\}$$
holds for all $`0tT`$.
## 9 Bounds for the virtual velocity
We prove here the assertion that $`v`$ does not grow too much as long as the $`L^3`$ norm of $`C`$ is not too large. We recall that $`v`$ solves (20)
$$\mathrm{\Gamma }v_i=2\nu C_{m,k;i}_kv_m+Q_{ji}f_j.$$
(56)
We multiply by $`v_i|v|^{2(m1)}`$ and integrate:
$$\frac{1}{2m}\frac{d}{dt}|v(x,t)|^{2m}𝑑x+\nu |v(x,t)|^2|v(x,t)|^{2(m1)}𝑑x+$$
$$+\nu \frac{m1}{2}||v(x,t)|^2|^2|v(x,t)|^{2(m2)}𝑑x=$$
$$=2\nu C_{m,k;i}(x,t)(_kv_m(x,t))v_i(x,t)|v(x,t)|^{2(m1)}𝑑x+$$
$$+Q_{ji}(x,t)f_j(x,t)v_i(x,t)|v(x,t)|^{2(m1)}𝑑x.$$
(57)
We bound
$$2\nu \left|C_{m,k;i}(x,t)(_kv_m(x,t))v_i(x,t)|v(x,t)|^{2(m1)}𝑑x\right|$$
$$\nu |v(x,t)|^2|v(x,t)|^{2(m1)}𝑑x+\nu |C(x,t)|^2|v(x,t)|^{2m}𝑑x$$
where
$$|C(x,t)|^2=\underset{m,k,i}{}|C_{m,k;i}(x,t)|^2,$$
(58)
and we bound
$$\left|Q_{ji}(x,t)f_j(x,t)v_i(x,t)|v(x,t)|^{2(m1)}𝑑x\right|$$
$$\left\{|g(x,t)|^{2m}𝑑x\right\}^{\frac{1}{2m}}\left\{|v(x,t)|^{2m}𝑑x\right\}^{\frac{2m1}{2m}}$$
where
$$g_i(x,t)=Q_{ji}(x,t)f_j(x,t).$$
(59)
The inequality obtained is
$$\frac{d}{dt}|v(x,t)|^{2m}𝑑x+\nu m(m1)||v(x,t)|^2|^2|v(x,t)|^{2(m2)}𝑑x$$
$$2m\nu |C(x,t)|^2|v(x,t)|^{2m}𝑑x+$$
$$+\mathrm{\hspace{0.17em}2}m\left\{|g(x,t)|^{2m}𝑑x\right\}^{\frac{1}{2m}}\left\{|v(x,t)|^{2m}𝑑x\right\}^{\frac{2m1}{2m}}$$
(60)
Let us consider for any $`m1`$ the quantity
$$q(x,t)=|v(x,t)|^m.$$
The inequality (60) implies
$$\frac{d}{dt}(q(x,t))^2+4\nu (1\frac{1}{m})|q(x,t)|^22m\nu |C(x,t)|^2(q(x,t))^2𝑑x+$$
$$+\mathrm{\hspace{0.17em}2}m\left\{|g(x,t)|^{2m}𝑑x\right\}^{\frac{1}{2m}}\left\{(q(x,t))^2𝑑x\right\}^{\frac{2m1}{2m}}$$
Using the well-known Morrey-Sobolev inequality
$$\left\{(q(x))^6𝑑x\right\}^{\frac{1}{3}}C_0\left\{|q(x)|^2𝑑x+L^2(q(x))^2𝑑x\right\}$$
and Hölder’s inequality we deduce
$$\frac{d}{dt}(q(x,t))^2+4\nu (1\frac{1}{m})|q(x,t)|^2$$
$$2m\nu C_0\left\{|C(x,t)|^3𝑑x\right\}^{\frac{2}{3}}\left\{|q(x,t)|^2𝑑x+L^2(q(x,t))^2𝑑x\right\}$$
$$+\mathrm{\hspace{0.17em}2}m\left\{|g(x,t)|^{2m}𝑑x\right\}^{\frac{1}{2m}}\left\{(q(x,t))^2𝑑x\right\}^{\frac{2m1}{2m}}$$
Therefore, if on the time interval $`t[0,\tau ]`$ $`C(x,t)`$ obeys the smallness condition
$$\left\{|C(x,t)|^3𝑑x\right\}^{\frac{1}{3}}\sqrt{\frac{2(m1)}{C_0m^2}}$$
(61)
then we have the inequality
$$\frac{d}{dt}v(,t)_{L^{2m}}\frac{\nu (m1)}{2m^2L^2}v(,t)_{L^{2m}}+g(,t)_{L^{2m}}$$
for $`t[0,\tau ]`$ and consequently
$$v(,t)_{L^{2m}}v_0_{L^{2m}}e^{\frac{\nu (m1)t}{2m^2L^2}}+_0^tg(,s)_{L^{2m}}$$
(62)
holds on the same time interval.
## 10 Appendix A
In this appendix we prove the inequality (39) and derive the explicit expression for $`K_{\mathrm{}}`$. The calculation is based on (). All constants $`C`$ are non-dimensional and may change from line to line. Solutions $`u`$ of the Navier-Stokes equations obey the differential inequality
$$\frac{d}{ds}|u(x,s)|^2𝑑x+\nu |\mathrm{\Delta }u(x,s)|^2𝑑x$$
$$\frac{C}{\nu ^3}\left(|u(x,s)|^2𝑑x\right)^3+\frac{C}{\nu }|f(x,s)|^2𝑑x.$$
The idea of () was to divide by an appropriate quantity to make use of the balance (31). The quantity is
$$(G(s))^2=\left(\gamma ^2+|u(x,s)|^2𝑑x\right)^2$$
where $`\gamma `$ is a positive constant that does not depend on $`s`$ and will be specified later. Dividing by $`(G(s))^2`$, integrating in time from $`t_0`$ to $`t`$ and using (31) one obtains
$$_{t_0}^t\mathrm{\Delta }u(,s)_{L^2}^2(G(s))^2𝑑s$$
$$C\left(\frac{K_0}{\nu ^5}+\frac{1}{\nu \gamma ^2}+\frac{1}{\nu ^2\gamma ^4}_{t_0}^tf(,s)^2𝑑s\right).$$
The three dimensional Sobolev embedding-interpolation inequality for periodic mean-zero functions
$$u_L^{\mathrm{}}Cu_{L^2}^{\frac{1}{2}}\mathrm{\Delta }u_{L^2}^{\frac{1}{2}}$$
is elementary. From it we deduce
$$u(,s)_L^{\mathrm{}}Cu(,s)_{L^2}^{\frac{1}{2}}(G(s))^{\frac{1}{2}}\left[\mathrm{\Delta }u(,s)_{L^2}G(s)^1\right]^{\frac{1}{2}}$$
Integrating in time, using the Hölder inequality, the inequality (31) and the inequalities above we deduce
$$_{t_0}^tu(,s)_L^{\mathrm{}}Cr$$
where the length $`r=r(\gamma ,t,\nu ,K_0)`$ is given in terms of six length scales
$$\frac{K_0}{\nu ^2}=r_0,\frac{\nu ^2}{\gamma ^2}=r_1,\frac{(tt_0)\gamma ^2}{\nu }=r_2,$$
$$(\gamma (tt_0))^{\frac{2}{3}}=r_3,\frac{tt_0}{\nu ^2}_{t_0}^tf(,s)_{L^2}^2𝑑s=r_4$$
and
$$r_5=\sqrt{\nu (tt_0)}.$$
The expression for $`r`$ is
$$r=r_0+(r_0)^{\frac{3}{4}}(r_1)^{\frac{1}{4}}+(r_0)^{\frac{1}{2}}(r_2)^{\frac{1}{2}}+(r_0)^{\frac{1}{4}}(r_3)^{\frac{3}{4}}+$$
$$(r_0)^{\frac{1}{4}}(r_4)^{\frac{1}{4}}(r_5)^{\frac{1}{2}}+(r_0)^{\frac{3}{4}}(r_4)^{\frac{1}{4}}\left(\frac{r_1}{r_2}\right)^{\frac{1}{4}}$$
The choice
$$\gamma ^4=\frac{\nu ^3}{tt_0}$$
entrains
$$r_1=r_2=r_3=r_5$$
reducing thus the number of length scales to three, the energy viscous length scale $`r_0`$, the diffusive length scale $`r_5`$ and the force length scale $`r_4`$. The bound becomes
$$K_{\mathrm{}}=C(r_0+r_4+r_5)$$
i.e. (40).
## 11 Appendix B
We prove her the commutation relation (17). We take an arbitrary function $`g`$ and compute $`[\mathrm{\Gamma },L_ig])`$ where $`\mathrm{\Gamma }=\mathrm{\Gamma }_\nu (u,)`$ and $`L_i=_A^i`$. We use first (11):
$$[\mathrm{\Gamma },L_ig]=\mathrm{\Gamma }\left(Q_{ji}_jg\right)Q_{ji}_j\mathrm{\Gamma }g=$$
$$\mathrm{\Gamma }(Q_{ji})_jg+Q_{ji}\mathrm{\Gamma }_jg2\nu _k(Q_{ji})_k_jgQ_{ji}_j\mathrm{\Gamma }g=$$
(commuting in the last term $`_j`$ and $`\mathrm{\Gamma }`$)
$$\mathrm{\Gamma }(Q_{ji})_jg2\nu _k(Q_{ji})_k_jgQ_{ji}_j(u_k)_kg=$$
(changing names of dummy indices in the last term)
$$\left(\mathrm{\Gamma }(Q_{ji})Q_{ki}_k(u_j)\right)_jg2\nu _k(Q_{ji})_k_jg=$$
(using (16))
$$2\nu Q_{jp}(_l_kA_p)(_kQ_{li})_jg2\nu _k(Q_{ji})_k_jg=$$
(using the definition (6) of $`_A`$
$$2\nu (_l_kA_p)(_kQ_{li})(L_pg)2\nu _k(Q_{ji})_k_jg=$$
(renaming dummy indices in the last expression)
$$2\nu (_k(Q_{li}))\left\{(_l_kA_p)(L_pg)_l_kg\right\}=$$
(using (7) in the last expression)
$$2\nu (_k(Q_{li}))\left\{(_l_kA_p)(L_pg)_k\left(_l(A_p)(L_pg)\right)\right\}=$$
(carrying out the differentiation in the last term and cancelling)
$$2\nu (_k(Q_{li}))_l(A_p)_k(L_p(g))=$$
(using the differential consequence of the fact that $`Q`$ and $`A`$ are inverses of each other)
$$2\nu Q_{li}(_k_l(A_p))_k(L_pg)=$$
(using the definition (6) of $`_A`$)
$$2\nu L_i(_k(A_p))_k(L_pg)=$$
(using the definition (9) of $`C_{m,k;i}`$)
$$2\nu C_{p,k;i}_k(L_pg),$$
and that concludes the proof. We proceed now to prove (18). We start with (15)
$$\mathrm{\Gamma }(_kA_m)=(_ku_j)(_jA_m)$$
and apply $`L_i`$:
$$L_i(\mathrm{\Gamma }(_kA_m))=L_i\left\{(_ku_j)(_jA_m)\right\}.$$
Using the commutation relation (17) and the definition (9) we get
$$\mathrm{\Gamma }(C_{m,k;i})=L_i\left\{(_ku_j)(_jA_m)\right\}+2\nu C_{p,l;i}_lL_p(_kA_m).$$
Using the fact that $`L_i`$ is a derivation in the first term and the definition (9) in the last term we conclude that
$$\mathrm{\Gamma }(C_{m,k;i})=(_ku_j)C_{m,j;i}(_jA_m)(L_i(_ku_j))+2\nu C_{p,l;i}(_lC_{m,k;p})$$
which is (18). We compute now the formal adjoint of $`_A^i`$
$$(_A^i)^{}g=_j(Q_{ji}g)=_A^i(g)(_j(Q_{ji}))g$$
(with (7))
$$(_A^i)^{}g=_A^i(g)\left\{(_jA_p)L_p(Q_{ji})\right\}g=$$
(using the fact that $`Q`$ is the inverse of $`A`$)
$$(_A^i)^{}g=_A^i(g)+Q_{ji}C_{p,j;p}g.$$
Acknowledgments. Part of this work was done at the Institute for Theoretical Physics in Santa Barbara, whose hospitality is gratefully acknowledged. This research is supported in part by NSF- DMS9802611.
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# Vacuum Polarization Using Quantum Mechanical Path Integrals
## 1 Introduction
This paper was initiated by several recent articles on so-called string-inspired calculations in scalar and spinor quantum electrodynamics (QED). World-line techniques and standard quantum mechanical path integrals were used to calculate effective actions and loop processes without employing any loop-momentum integrals. Our intention here is to demonstrate in a didactic and rather explicit manner how these methods can be put to work in a simplified model which avoids charge, spin or polarization degrees of freedom, but brings out all the necessary physics that one would encounter in more realistic models.
String-inspired methods in quantum field theory were first used in the works of Bern and Kosower <sup>1</sup>. These authors and Strassler <sup>2</sup> then realized that some of the well-known vacuum processes in QED and QCD can be computed rather easily with the aid of one-dimensional path integrals for relativistic point particles. Similar techniques and results can also be found in the monograph by Polyakov <sup>3</sup>. String-inspired methods, particularly in QED, were then extensively studied in a series of papers by Schmidt, Schubert and Reuter; cf., e.g., Ref. 4, where the state of the art is reviewed extensively. There are also contributions by McKeon <sup>5</sup> and various co-authors who have proved that world-line methods are extremely useful.
## 2 Vacuum Polarization in a Model Field Theory <br>$`^{}=\frac{g}{2}\psi ^2\varphi `$
To have a comparatively simple model let us consider an interaction Lagrangian
$$^{}=\frac{g}{2}\psi ^2(x)\varphi (x),$$
(1)
where $`g`$ is the coupling constant. In the following we will be mainly interested in the one-loop vacuum graph (Fig. 1).
In QED the particle circulating in the loop would be the electron which is tied to an arbitrary number of off-shell photons. As is well known (see, e.g., Ref. 6), loop graphs belong to a subclass of Feynman diagrams called one-particle-irreducible diagrams. Their associated one-particle-irreducible amplitudes $`\mathrm{\Gamma }_N(x_1,x_2,\mathrm{},x_N)`$ can be obtained with the aid of a generating functional:
$$\mathrm{\Gamma }[\varphi ]=\underset{N=0}{\overset{\mathrm{}}{}}\frac{1}{N!}d^4x_1d^4x_2\mathrm{}d^4x_N\mathrm{\Gamma }_N(x_1,\mathrm{},x_N)\varphi (x_1)\mathrm{}\varphi (x_N).$$
(2)
Our interest lies with $`N=2`$. But for the time being we will let $`N=1,2,\mathrm{},\mathrm{}`$. Now recall from potential theory that a particle of mass $`m`$ travelling to all orders in an external field $`\varphi (x)`$ is given by
(3)
where $`\mathrm{\Delta }_+`$ is the Green’s function of the freely propagating $`\psi `$-particle ($`\mathrm{\Delta }_+\mathrm{\Delta }_+[\varphi =0]`$) which satisfies the Green’s function equation $`(^2+m^2)\mathrm{\Delta }_+(xy)=\delta ^4(xy)`$, or in momentum space, $`(p^2+m^2)\mathrm{\Delta }_+(p)=1`$. Our metric signature is $`(,+,+,+)`$. Summing up the terms in the geometric (Born-)series (3) we obtain
$`\mathrm{\Delta }_+[\varphi ]`$ $`=`$ $`\mathrm{\Delta }_+\left(1+g\varphi \mathrm{\Delta }_++g\varphi \mathrm{\Delta }_+g\varphi \mathrm{\Delta }_++\mathrm{}\right)`$ (4)
$`=`$ $`\mathrm{\Delta }_+\left(1g\varphi \mathrm{\Delta }_+\right)^1.`$
It is rather trivial to rewrite Eq. (4) in the form
$$(p^2+m^2g\varphi )\mathrm{\Delta }_+[\varphi ]=1,$$
(5)
or in $`x`$-representation
$$\left(^2+m^2g\varphi (x)\right)\mathrm{\Delta }_+(x,y|\varphi )=\delta ^4(xy).$$
(6)
Given these simple facts we can give an analytical expression for our graph Fig. 1, namely,
$$\text{i}\mathrm{\Gamma }_N(x_1,\mathrm{},x_N)=(N1)!\frac{1}{2}\left(g\mathrm{\Delta }_+(x_1,x_2)\right)\left(g\mathrm{\Delta }_+(x_2,x_3)\right)\mathrm{}\left(g\mathrm{\Delta }_+(x_N,x_1)\right).$$
(7)
The individual factors have the following origin:
* The factor $`(N1)!`$ takes into account that, after fixing one of the $`\varphi `$-lines, a total of $`(N1)!`$ topological inequivalent graphs can be created by permutation of the remaining $`\varphi `$-lines.
* One factor $`g`$ is assigned to each vertex.
* A free propagator $`\mathrm{\Delta }_+`$ is assigned to each $`\psi `$-line.
* The factor $`\frac{1}{2}`$ is indicative of a neutral scalar field theory.
$`\mathrm{\Gamma }_N`$ does not, by definition, contain external propagators. Substituting Eq. (7) into (2) we obtain $`\mathrm{\Gamma }[\varphi ]`$ in one-loop approximation
$`\text{i}\mathrm{\Gamma }[\varphi ]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{N=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(N1)!}{N!}}{\displaystyle d^4x_1\mathrm{}d^4x_N(g)^N\mathrm{\Delta }_+(x_1,x_2)\mathrm{}\mathrm{\Delta }_+(x_N,x_1)\varphi (x_1)\mathrm{}\varphi (x_N)}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{N=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{1}{N}}{\displaystyle d^4x_1\mathrm{}d^4x_N\left(g\varphi (x_1)\mathrm{\Delta }_+(x_1,x_2)\right)\mathrm{}\left(g\varphi (x_N)\mathrm{\Delta }_+(x_N,x_1)\right)}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^4x_1\underset{N=1}{\overset{\mathrm{}}{}}\frac{1}{N}d^4x_2\mathrm{}d^4x_Nx_1|g\varphi \mathrm{\Delta }_+|x_2\mathrm{}x_N|g\varphi \mathrm{\Delta }_+|x_1}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^4x_1\underset{N=1}{\overset{\mathrm{}}{}}\frac{1}{N}x_1|(g\varphi \mathrm{\Delta }_+)^N|x_1}={\displaystyle \frac{1}{2}}{\displaystyle d^4xx|\underset{N=1}{\overset{\mathrm{}}{}}\frac{1}{N}(g\varphi \mathrm{\Delta }_+)^N|x}`$
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^4xx|\mathrm{ln}(1g\varphi \mathrm{\Delta }_+)|x}={\displaystyle \frac{1}{2}}\text{Tr}_x\mathrm{ln}(1g\varphi \mathrm{\Delta }_+).`$
Here we used the series $`\mathrm{ln}(1x)=_{n=1}^{\mathrm{}}\frac{x^n}{n}`$, $`x(1,+1)`$. We also made use of the completeness relation $`d^4y|yy|=1`$ and wrote $`x|g\varphi \mathrm{\Delta }_+|y=g\varphi (x)x|\mathrm{\Delta }_+|y=g\varphi (x)\mathrm{\Delta }_+(x,y)`$. So our final result reads
$$\text{i}\mathrm{\Gamma }[\varphi ]=\frac{1}{2}\text{Tr}\mathrm{ln}(1g\varphi \mathrm{\Delta }_+)^1,$$
(8)
or with the aid of (4):
$$\text{i}\mathrm{\Gamma }[\varphi ]=\frac{1}{2}\text{Tr}\mathrm{ln}\left[\frac{\mathrm{\Delta }_+[\varphi ]}{\mathrm{\Delta }_+[0]}\right]$$
(9)
and since
$$\mathrm{\Delta }_+[0]\mathrm{\Delta }_+=\frac{1}{p^2+m^2\text{i}ϵ}\text{and}\mathrm{\Delta }_+[\varphi ]=\frac{1}{p^2+m^2g\varphi \text{i}ϵ},$$
we have
$$\text{i}\mathrm{\Gamma }[\varphi ]=\frac{1}{2}\text{Tr}\mathrm{ln}\frac{p^2+m^2g\varphi \text{i}ϵ}{p^2+m^2\text{i}ϵ}.$$
(10)
Here we employ the formula
$$\mathrm{ln}\frac{a}{b}=\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{s}\text{e}^{\text{i}s(b\text{i}ϵ)}\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{s}\text{e}^{\text{i}s(a\text{i}ϵ)}$$
(11)
and obtain
$$\text{i}\mathrm{\Gamma }[\varphi ]=\frac{1}{2}\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{s}\text{Tr}\text{e}^{\text{i}s(p^2g\varphi +m^2\text{i}ϵ)}+\frac{1}{2}\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{s}\text{Tr}\text{e}^{\text{i}s(p^2+m^2\text{i}ϵ)}.$$
(12)
Since the last term is $`\varphi `$-independent it is usually dropped.
Now we turn to the computation of
$`\text{i}\mathrm{\Gamma }[\varphi ]`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{s}}\text{Tr}_x\text{e}^{\text{i}s(p^2g\varphi +m^2)},m^2m^2\text{i}ϵ,`$ (13)
$`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{s}}\text{e}^{\text{i}sm^2}{\displaystyle d^4xx|\text{e}^{\text{i}(p^2g\varphi )s}|x}.`$
Introducing the “Hamiltonian” $`H=p^2g\varphi `$, $`p_\mu =\frac{1}{\text{i}}_\mu `$, we need to calculate the trace of the quantum mechanical transition amplitude
$$x,s|y,0=x|\text{e}^{\text{i}Hs}|y.$$
(14)
Instead of working with the Hamiltonian we now switch over to the Lagrangian description of our system so that we can make use of Feynman’s path integral representation of Eq. (14).
Since $`L=p\dot{x}H`$, $`\dot{x}=\frac{H}{p}=2p`$: $`p=\frac{\dot{x}}{2}`$, we find
$$L=\frac{\dot{x}^2}{2}\left(\frac{\dot{x}^2}{4}g\varphi \right)=\frac{\dot{x}^2}{4}+g\varphi .$$
(15)
Now, Feynman’s path integral representation of the transition amplitude is given by
$$x,s|y=𝒩\underset{x(0)=y,x(s)=x}{}𝒟x(\tau )\text{e}^{\text{i}S_{\text{cl}}[x(\tau )]},$$
(16)
where
$$S_{\text{cl}}[x(\tau )]=\underset{0}{\overset{s}{}}𝑑\tau L(x(\tau ),\dot{x}(\tau )),L=\frac{\dot{x}^2}{4}+g\varphi .$$
(17)
This is all we need from Feynman’s book <sup>7</sup> or any other monograph on Feynman path integrals in single-particle quantum mechanics<sup>8</sup>. Already at this stage we want to emphasize that nowhere in the sequel do we have to compute a loop-momentum integral as is usually required in other field theoretic approaches.
The normalization factor $`𝒩`$ in (16) is determined from the free theory, $`g=0`$, i.e., $`H_0=p^2`$. In this case we have
$$x|\text{e}^{\text{i}sp^2}|y=𝒩\underset{x(0)=y}{\overset{x(s)=x}{}}𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \frac{\dot{x}^2}{4}}.$$
(18)
Taking the trace on the left-hand side gives us $`\left((dx)d^4x\right)`$
$`{\displaystyle (dx)x|\text{e}^{\text{i}sp^2}|x}={\displaystyle (dx)x|\text{e}^{\text{i}sp^2}\mathrm{𝟙}|x},\mathrm{𝟙}={\displaystyle (dp)|pp|}`$ (19)
$`=`$ $`{\displaystyle (dx)(dp)\text{e}^{\text{i}sp^2}x|pp|x},x|p={\displaystyle \frac{\text{e}^{\text{i}px}}{(2\pi )^2}},{\displaystyle (dx)}=V_4`$
$`=`$ $`V_4{\displaystyle \frac{(dp)}{(2\pi )^4}\text{e}^{\text{i}sp^2}}=V_4\left({\displaystyle \frac{\text{i}}{(4\pi )^2}}\right){\displaystyle \frac{1}{s^2}}.`$
Here the four-dimensional integral over the (3+1)-dimensional momentum space is computed as
$`{\displaystyle \frac{(dp)}{(2\pi )^4}\text{e}^{\text{i}sp^2}}`$ $`=`$ $`\left({\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp_1}{2\pi }}\text{e}^{\text{i}sp_1^2}\right)^3\left({\displaystyle \underset{\mathrm{}}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dp_0}{2\pi }}\text{e}^{\text{i}sp_0^2}\right)=\left({\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{\pi }{\text{i}s}}\right)^{1/2}\right)^3{\displaystyle \frac{1}{2\pi }}\left({\displaystyle \frac{\pi \text{i}}{s}}\right)^{1/2}`$
$`=`$ $`{\displaystyle \frac{1}{(4\pi )^2}}{\displaystyle \frac{1}{\text{i}s^2}}.`$
The trace on the right-hand side of Eq. (18) is written as
$$𝒩(dx)\underset{x(0)=x(s)}{}𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \frac{\dot{x}^2}{4}}𝒩\underset{x(0)=x(s),\text{arbitrary}}{}𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \frac{\dot{x}^2}{4}},$$
so that we end up with the useful relation
$$𝒩\underset{x(0)=x(s),\text{arbitr.}}{}𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \frac{\dot{x}^2}{4}}=V\frac{1}{\text{i}(4\pi )^2}\frac{1}{s^2}.$$
(20)
Here then is our path integral representation for the one-loop process with an arbitrary number of external off-shell $`\varphi `$-particle lines ( s. Eq. (13)):
$$\text{i}\mathrm{\Gamma }[\varphi ]=\frac{1}{2}\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{s}\text{e}^{\text{i}m^2s}𝒩\underset{x(0)=x(s),\text{arbitr.}}{}𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \left[\frac{\dot{x}^2}{4}+g\varphi \right]}.$$
(21)
Using a perturbative expansion of the right-hand side of Eq. (21) we can write
$`\mathrm{\Gamma }[\varphi ]`$ $`=`$ $`{\displaystyle \frac{\text{i}}{2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{s}}\text{e}^{\text{i}m^2s}𝒩{\displaystyle \underset{x(0)=x(s),}{}}𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \frac{\dot{x}^2}{4}}{\displaystyle \underset{N=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\text{i}g)^N}{N!}}\left({\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau \varphi (x(\tau ))\right)^N`$ (22)
$`=`$ $`{\displaystyle \frac{\text{i}}{2}}{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{s}}\text{e}^{\text{i}m^2s}𝒩{\displaystyle \underset{x(0)=x(s),}{}}𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \frac{\dot{x}^2}{4}}{\displaystyle \underset{N=1}{\overset{\mathrm{}}{}}}{\displaystyle \frac{(\text{i}g)^N}{N!}}{\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _i\varphi (x(\tau _i))`$
$`=:`$ $`{\displaystyle \underset{N=1}{\overset{\mathrm{}}{}}}\mathrm{\Gamma }_N[\varphi ],`$
with
$$\mathrm{\Gamma }_N[\varphi ]=\frac{\text{i}}{2}\frac{(\text{i}g)^N}{N!}\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{s}\text{e}^{\text{i}m^2s}𝒩\underset{x(0)=x(s),}{}𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \frac{\dot{x}^2}{4}}\underset{i=1}{\overset{N}{}}\underset{0}{\overset{s}{}}𝑑\tau _i\varphi (x(\tau _i)).$$
(23)
At this point we specialize the $`\varphi `$-field to a sum of plane waves,
$$\varphi (x)=\underset{i=1}{\overset{N}{}}\text{e}^{\text{i}k_ix}.$$
(24)
For $`N=2`$ the $`\varphi `$-term on the right-hand side of Eq. (23) would read
$`{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _1𝑑\tau _2\varphi (x(\tau _1))\varphi (x(\tau _2))={\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _1𝑑\tau _2\left(\text{e}^{\text{i}k_1x(\tau _1)}+\text{e}^{\text{i}k_2x(\tau _1)}\right)\left(\text{e}^{\text{i}k_1x(\tau _2)}+\text{e}^{\text{i}k_2x(\tau _2)}\right)`$
$`={\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _1𝑑\tau _2\left(\mathrm{}+\text{e}^{\text{i}k_1x(\tau _1)}\text{e}^{\text{i}k_2x(\tau _2)}+\text{e}^{\text{i}k_2x(\tau _1)}\text{e}^{\text{i}k_1x(\tau _2)}+\mathrm{}\right)`$
$`=2!{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _1𝑑\tau _2\text{e}^{\text{i}k_1x(\tau _1)}\text{e}^{\text{i}k_2x(\tau _2)}+\mathrm{}.`$ (25)
Here we kept only mixed terms in $`k_1`$ and $`k_2`$, i.e., each $`\varphi `$-mode occurs only once.
Generalizing to $`N`$ we would find $`N!`$ instead of $`2!`$ in Eq. (25). This factor $`N!`$ then cancels the $`N!`$ that stands in the denominator of Eq. (23). So far we have
$$\mathrm{\Gamma }_N[k_1,k_2,\mathrm{},k_N]=(\text{i})\frac{1}{2}(\text{i}g)^N\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{s}𝒩𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \left[\frac{\dot{x}^2}{4}m^2\right]}\left(\underset{i=1}{\overset{N}{}}\underset{0}{\overset{s}{}}𝑑\tau _i\text{e}^{\text{i}k_ix(\tau _i)}\right).$$
(26)
Introducing the “current” for $`x`$,
$$j(\tau )=\text{i}\underset{j=1}{\overset{N}{}}k_j\delta (\tau \tau _j),$$
(27)
we can rewrite Eq. (26) in the form
$`\mathrm{\Gamma }_N[k_1,k_2,\mathrm{},k_N]`$ $`=`$ $`(\text{i}){\displaystyle \frac{1}{2}}(\text{i}g)^N{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{s}}\text{e}^{\text{i}m^2s}\left({\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _i\right)𝒩{\displaystyle 𝒟x\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \frac{\dot{x}^2}{4}}\text{e}^{\underset{0}{\overset{s}{}}𝑑\tau j(\tau )x(\tau )}}`$
$`=`$ $`(\text{i}){\displaystyle \frac{1}{2}}(\text{i}g)^N{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{s}}\text{e}^{\text{i}m^2s}\left({\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _i\right)𝒩{\displaystyle 𝒟x\text{e}^{\frac{\text{i}}{4}\underset{0}{\overset{s}{}}𝑑\tau x\frac{d^2}{d\tau ^2}x}\text{e}^{\underset{0}{\overset{s}{}}𝑑\tau jx}}.`$
The operator $`\frac{d^2}{d\tau ^2}`$, acting on $`x(\tau )`$ with periodical boundary condition $`x(s)=x(0)`$, has zero-modes $`x_0`$:
$$\frac{d^2}{d\tau ^2}x(\tau )=\lambda _nx(\tau ),x(s)=x(0),$$
zero mode $`xx_0=`$const.: $`\frac{d^2}{d\tau ^2}x_0=0`$, $`\lambda _0=0`$.
These zero modes will be separated from their orthogonal non-zero mode partners $`\xi (\tau )`$ by writing $`x(\tau )=x_0+\xi `$ with $`_0^s𝑑\tau \xi (\tau )=0`$, i.e., $`_0^s𝑑\tau x(\tau )=x_0`$ and
$$𝒟x=d^4x_0𝒟\xi .$$
(29)
Since $`x(\tau )`$ (and therefore $`\xi (\tau )`$) is periodic we can write
$$x(\tau )=x_0+\underset{n0}{}c_n\text{e}^{\frac{2\pi \text{i}n}{s}\tau }.$$
(30)
The zero-mode contribution in the source term of Eq. (LABEL:28) yields
$$\text{e}^{\underset{0}{\overset{s}{}}jx𝑑\tau }=\text{e}^{x_0\text{i}\underset{0}{\overset{s}{}}𝑑\tau _{j=1}^Nk_j\delta (\tau \tau _j)}=\text{e}^{\text{i}x_0_{j=1}^Nk_j}$$
and therefore
$$d^4x_0\text{e}^{\text{i}x_0_{j=1}^Nk_j}=(2\pi )^4\delta ^4(k_1+k_2+\mathrm{}+k_N).$$
(31)
The result of combining all this information with Eq. (LABEL:28) produces the result
$`\mathrm{\Gamma }_N[k_1,\mathrm{},k_N]`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\text{i}g)^N(2\pi )^4\delta \left({\displaystyle \underset{j=1}{\overset{N}{}}}k_j\right){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{(4\pi )^2s^3}}\text{e}^{\text{i}m^2s}\left({\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _i\right)`$ (32)
$`\times {\displaystyle \frac{𝒟\xi \text{e}^{\frac{\text{i}}{4}\underset{0}{\overset{s}{}}𝑑\tau \xi \frac{d^2}{d\tau ^2}\xi }\text{e}^{\underset{0}{\overset{s}{}}𝑑\tau j\xi }}{𝒟\xi \text{e}^{\frac{\text{i}}{4}\underset{0}{\overset{s}{}}𝑑\tau \xi \frac{d^2}{d\tau ^2}\xi }}}.`$
Use has also been made of formula (20):
$$𝒩d^4x_0\underset{\xi (s)=\xi (0)}{}𝒟\xi \text{e}^{\frac{\text{i}}{4}\underset{0}{\overset{s}{}}𝑑\tau [\dot{\xi }+\dot{x}_0]^2}=\frac{\text{i}}{(4\pi )^2}\frac{1}{s^2}V,$$
where $`\dot{x}_0=0`$ and the four-volume $`d^4x_0=V`$ is cancelled on both sides.
In Eq. (32) we meet the path integral
$$\underset{\xi (0)=\xi (s)}{}𝒟\xi \text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \left(\frac{\dot{\xi }^2}{4}\text{i}j\xi \right)}=\underset{\xi (0)=\xi (s)}{}𝒟\xi \text{e}^{\text{i}S[\xi ]},$$
where
$$S[\xi ]=\underset{0}{\overset{s}{}}𝑑\tau L(\xi ,\dot{\xi }),\text{and}L(\xi ,\dot{\xi })=\frac{\dot{\xi }^2}{4}\text{i}j\xi .$$
(33)
Using Eq. (33) we find from $`\frac{d}{dt}\frac{L}{\dot{\xi }}\frac{L}{\xi }=0`$ the equation of motion,
$$\frac{1}{2}\frac{d^2}{d\tau ^2}\xi (\tau )=\text{i}j(\tau ),$$
(34)
which can be solved with the ansatz
$$\xi (\tau )=\text{i}\underset{0}{\overset{s}{}}𝑑\tau ^{}G(\tau ,\tau ^{})j(\tau ^{}),$$
where the Green’s function equation is given by (s. later)
$$\frac{1}{2}\frac{d^2}{d\tau ^2}G(\tau ,\tau ^{})=\delta (\tau \tau ^{})\frac{1}{s}.$$
(35)
Indeed we can write
$`{\displaystyle \frac{1}{2}}{\displaystyle \frac{d^2}{d\tau ^2}}\xi (\tau )`$ $`=`$ $`\text{i}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau ^{}\left[\delta (\tau \tau ^{}){\displaystyle \frac{1}{s}}\right]j(\tau ^{})=\text{i}j(\tau )+{\displaystyle \frac{\text{i}}{s}}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau ^{}j(\tau ^{})`$
$`=`$ $`\text{i}j(\tau )+{\displaystyle \frac{\text{i}}{s}}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau ^{}\text{i}{\displaystyle \underset{j=1}{\overset{N}{}}}k_j\delta (\tau ^{}\tau _j)=\text{i}j(\tau ){\displaystyle \frac{1}{s}}{\displaystyle \underset{j=1}{\overset{N}{}}}k_j`$
$`=`$ $`\text{i}j(\tau )`$
following from the $`\delta `$-function in Eq. (32): $`_{j=1}^Nk_j=0`$.
We now must compute the path integrals occurring in Eq. (32) which are of Gaussian type. So in an intermediate step we just consider a single Gaussian integral:
$`{\displaystyle \frac{dx}{\sqrt{\text{i}\pi }}\text{e}^{\text{i}(mx^2+vx)}}`$ $`=`$ $`{\displaystyle \frac{dx}{\sqrt{\text{i}\pi }}\text{e}^{\text{i}m\left(x+\frac{v}{2m}\right)^2}\text{e}^{\text{i}\frac{v^2}{4m}}}={\displaystyle \frac{dx}{\sqrt{\text{i}\pi }}\text{e}^{\text{i}mx^2}\text{e}^{\text{i}\frac{v^2}{4m}}}`$
$`=`$ $`{\displaystyle \frac{ds\text{e}^{s^2}}{\sqrt{\pi }}\frac{\text{e}^{\text{i}\frac{v^2}{4m}}}{\sqrt{m}}}=1{\displaystyle \frac{\text{e}^{\text{i}\frac{v^2}{4m}}}{\sqrt{m}}},x={\displaystyle \frac{\sqrt{\text{i}}}{\sqrt{m}}}.`$
The result for a product of coupled Gaussian integrals is therefore the generalization
$$\frac{dx_1}{\sqrt{\text{i}\pi }}\mathrm{}\frac{dx_n}{\sqrt{\text{i}\pi }}\text{e}^{\text{i}\left(_{l,m}x_lM_{lm}x_m+_lx_lv_l\right)}=\frac{\text{e}^{\text{i}_{l,m}v_l\frac{(M^1)_{lm}}{4}v_m}}{\sqrt{detM}},$$
(36)
which can be proved by making a rotation on the $`x`$’s and $`v`$’s which diagonalizes $`M`$ and so reduces to the one-variable case where now $`\sqrt{m}`$ is replaced by $`\sqrt{detM}=\sqrt{_{j=1}^nM_j}`$, where the $`M_j`$ are the eigenvalues of $`M`$. So when we first compute the discrete version of the path integrals in Eq. (32) and then go to the continuous limit, we evidently obtain (substituting $`M=\frac{1}{4}\frac{d^2}{d\tau ^2}`$, $`v=\frac{1}{\text{i}}j`$ in Eq. (36))
$`{\displaystyle \frac{𝒟\xi \text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \xi \left(\frac{1}{4}\frac{d^2}{d\tau ^2}\right)\xi }\text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \frac{1}{\text{i}}j\xi }}{𝒟\xi \text{e}^{\text{i}\underset{0}{\overset{s}{}}𝑑\tau \xi \left(\frac{1}{4}\frac{d^2}{d\tau ^2}\right)\xi }}}\left(\widehat{=}{\displaystyle \frac{\text{e}^{\text{i}_{i,j}v_i\left(\frac{1}{4M}\right)_{ij}v_j}}{\sqrt{detM}}}\sqrt{detM}\right)`$
$`\widehat{=}\text{e}^{\frac{\text{i}}{2}\underset{0}{\overset{s}{}}𝑑\tau \underset{0}{\overset{s}{}}𝑑\tau ^{}_{i,j}j_i\mathrm{\hspace{0.17em}2}\left(\frac{d^2}{d\tau ^2}\right)_{ij}^1j_j}=\text{e}^{\frac{\text{i}}{2}\underset{0}{\overset{s}{}}𝑑\tau \underset{0}{\overset{s}{}}𝑑\tau ^{}j^\mu (\tau )G_{\mu \nu }(\tau ,\tau ^{})j^\nu (\tau ^{})},`$
where
$$G_{\mu \nu }(\tau ,\tau ^{})=\eta _{\mu \nu }G(\tau ,\tau ^{}),\text{and}G(\tau ,\tau ^{})=\tau |2\left(\frac{d^2}{d\tau ^2}\right)^1|\tau ^{}.$$
(37)
Now we can write Eq. (32) in the form
$`\mathrm{\Gamma }_N[k_1,\mathrm{},k_N]`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\text{i}g)^N(2\pi )^4\delta \left(k_1+\mathrm{}+k_N\right)`$
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{(4\pi )^2s^3}}\text{e}^{\text{i}m^2s}\left({\displaystyle \underset{i=1}{\overset{N}{}}}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _i\right)\text{e}^{\frac{\text{i}}{2}\underset{0}{\overset{s}{}}𝑑\tau \underset{0}{\overset{s}{}}𝑑\tau ^{}j^\mu (\tau )G_{\mu \nu }(\tau ,\tau ^{})j^\nu (\tau ^{})}.`$
Note the structure expressed in Eq. (2), where loop-particle, mass $`m`$, and off-shell $`\varphi `$-particles are factorized in such a way that the scalar particle circulating in the loop becomes multiplied by the exponential term which is solely due to the $`\varphi `$-particles tied to the loop. With $`j(\tau )`$ given in Eq. (27) and
$`{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau {\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau ^{}jGj`$ $`=`$ $`{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau {\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau ^{}{\displaystyle \underset{i,j=1}{\overset{N}{}}}k_i\delta (\tau \tau _i)G(\tau ,\tau ^{})k_j\delta (\tau ^{}\tau _j)`$
$`=`$ $`{\displaystyle \underset{i,j=1}{\overset{N}{}}}k_ik_jG(\tau _i,\tau _j),`$
we finally obtain
$$\mathrm{\Gamma }_N[k_1,\mathrm{},k_N]=\frac{1}{2}(\text{i}g)^N(2\pi )^4\delta \left(k_1+\mathrm{}+k_N\right)\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{(4\pi )^2s^3}\text{e}^{\text{i}m^2s}\underset{i=1}{\overset{N}{}}\underset{0}{\overset{s}{}}𝑑\tau _i\text{e}^{\frac{\text{i}}{2}_{i,j}^Nk_ik_jG(\tau _i,\tau _j)}.$$
(39)
As is shown below, $`G(\tau ,\tau )=0=\dot{G}(\tau ,\tau )`$, so that there are no terms with $`k_i^2`$ present, i.e., without the use of on-shell conditions.
Now we must devote a few lines to the Green’s function of the problem. First note that the spectrum and the eigenmodes of the operator $`\frac{}{\tau }`$ are given by
$$\text{Spectrum}(_\tau )=\text{i}\frac{2\pi }{s}n,\tau |f_n=f_n(\tau )=\frac{1}{\sqrt{s}}\text{e}^{\text{i}\left(\frac{2\pi }{s}\right)n\tau },n,f_n(\tau +s)=f_n(\tau ).$$
One can check orthogonality,
$$f_n|f_m=\underset{0}{\overset{s}{}}𝑑\tau f_n^{}(\tau )f_m(\tau )=\delta _{nm},$$
and completeness,
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}f_n(\tau _2)f_n^{}(\tau _1)=\frac{1}{s}\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\text{e}^{\frac{2\pi \text{i}}{s}(\tau _2\tau _1)n}=\underset{m=\mathrm{}}{\overset{\mathrm{}}{}}\delta (\tau _2\tau _1ms),$$
by Poisson’s formula. Note that for $`0\tau _1,\tau _2s`$ we obtain
$$\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}f_n(\tau _2)f_n^{}(\tau _1)=\delta (\tau _2\tau _1)\text{in}\underset{0}{\overset{s}{}}𝑑\tau \mathrm{}.$$
We are interested in the spectrum of $`_\tau ^2`$ which is given by Spectrum$`(_\tau ^2)=\left(\text{i}\frac{2\pi }{s}n\right)^2=\frac{4\pi ^2}{s^2}n^2`$.
Earlier we defined the Green’s function
$$G(\tau _2,\tau _1)2\tau _2|(_\tau ^2)^1|\tau _1,$$
(40)
which we write as
$`2\tau _2|{\displaystyle \frac{1}{_\tau ^2}}|\tau _1`$ $`=`$ $`2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}\tau _2|f_nf_n|{\displaystyle \frac{1}{_\tau ^2}}|f_mf_m|\tau _1`$ (41)
$`=`$ $`2{\displaystyle \underset{n=\mathrm{}}{\overset{\mathrm{}}{}}}f_n(\tau _2){\displaystyle \frac{1}{\frac{4\pi ^2}{s^2}n^2}}f_n^{}(\tau _1)=2s{\displaystyle \underset{n=\mathrm{},n0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\text{e}^{2\pi \text{i}n\frac{\tau _2\tau _1}{s}}}{(2\pi \text{i}n)^2}}`$
$`=`$ $`|\tau _2\tau _1|{\displaystyle \frac{(\tau _2\tau _1)^2}{s}}{\displaystyle \frac{s}{6}}=G(\tau _2,\tau _1).`$ (42)
It can be easily seen that the constant $`\frac{s}{6}`$ drops out in scattering amplitude calculations, so that we can omit it at the beginning.
Now we are able to prove the Green’s function equation mentioned earlier in Eq. (35):
$`{\displaystyle \frac{1}{2}}_\tau ^2G(\tau ,\tau ^{})`$ $`\stackrel{(\text{40}),(\text{41})}{=}`$ $`{\displaystyle \frac{1}{s}}{\displaystyle \underset{n=\mathrm{},n0}{\overset{\mathrm{}}{}}}\text{e}^{\frac{2\pi \text{i}}{s}n(\tau \tau ^{})}={\displaystyle \frac{1}{s}}{\displaystyle \underset{n}{}}\text{e}^{\frac{2\pi \text{i}}{s}n(\tau \tau ^{})}{\displaystyle \frac{1}{s}}={\displaystyle \underset{m}{}}\delta (\tau \tau ^{}ms){\displaystyle \frac{1}{s}}`$
$`\stackrel{m=0}{=}`$ $`\delta (\tau \tau ^{}){\displaystyle \frac{1}{s}}\text{for}0\tau ,\tau ^{}s.`$
($`\frac{1}{s}`$ comes from the zero-mode, $`n=0`$, which is subtracted.) With the above definition of the Green’s function Eq. (42), we can easily prove the following relations:
$`G(\tau ,\tau ^{})`$ $`=`$ $`G(\tau \tau ^{})=G(\tau ^{},\tau ),\text{symmetry}`$
$`_\tau G(\tau ,\tau ^{})`$ $``$ $`\dot{G}(\tau ,\tau ^{})=\text{sign}(\tau \tau ^{}){\displaystyle \frac{2(\tau \tau ^{})}{s}}`$ (43)
$`\dot{G}(\tau ,\tau ^{})`$ $`=`$ $`\dot{G}(\tau ^{},\tau ),\dot{}={\displaystyle \frac{d}{d\tau }}.`$
Since we wanted to calculate the vacuum polarization diagram with two $`\varphi `$-lines we now specialize to $`N=2`$ in Eq. (39):
$`\mathrm{\Gamma }_2[k_1,k_2]`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^2(2\pi )^4\delta ^4(k_1+k_2){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{(4\pi )^2s^3}}\text{e}^{\text{i}m^2s}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _1{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _2\text{e}^{\frac{\text{i}}{2}\left(k_1k_2G(\tau _1,\tau _2)+k_2k_1G(\tau _2,\tau _1)\right)}`$ (44)
$`=`$ $`{\displaystyle \frac{1}{2}}g^2(2\pi )^4\delta ^4(k_1+k_2){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{(4\pi )^2s^3}}\text{e}^{\text{i}m^2s}{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _1{\displaystyle \underset{0}{\overset{s}{}}}𝑑\tau _2\text{e}^{\text{i}k_1k_2G(\tau _1,\tau _2)},`$
where we used the symmetry of $`G`$: $`G(\tau _1,\tau _2)=G(\tau _2,\tau _1)`$. Since $`G(\tau _1,\tau _2)`$ is periodic and is only dependent on the difference $`(\tau _1\tau _2)`$, the integration over $`\tau _2`$ is trivial; after the $`\tau _1`$-integration there is no dependence on $`\tau _2`$ left and hence the integration over $`\tau _2`$ gives just $`s`$. We can then choose $`\tau _2=0`$:
$$\mathrm{\Gamma }_2[k_1,k_2]=\frac{1}{2}g^2(2\pi )^4\delta ^4(k_1+k_2)\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{(4\pi )^2s^3}\text{e}^{\text{i}m^2s}\underset{0}{\overset{s}{}}𝑑\tau _1\text{e}^{\text{i}k_1k_2G(\tau _1)}.$$
(45)
Now we use
$`G(\tau _1)`$ $``$ $`G(\tau _1,0)={\displaystyle \frac{\tau _1^2}{s}}+\tau _1`$
$`v`$ $`=`$ $`{\displaystyle \frac{2\tau _1}{s}}1,d\tau _1={\displaystyle \frac{s}{2}}dv,\tau _1=(v+1){\displaystyle \frac{s}{2}}`$
$`0\tau _1s,1v1,G={\displaystyle \frac{s}{4}}(1v^2)`$
and obtain
$`\mathrm{\Gamma }_2[k_1,k_2]`$ $`=`$ $`{\displaystyle \frac{1}{2}}g^2(2\pi )^4\delta ^4(k_1+k_2){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{2(4\pi )^2s}}\text{e}^{\text{i}m^2s}{\displaystyle \underset{1}{\overset{1}{}}}𝑑v\text{e}^{\text{i}k_1^2\frac{s}{4}(1v^2)}`$
$`=`$ $`(2\pi )^4\delta ^4(k_1+k_2)\mathrm{\Pi }(k^2),`$
with
$$\mathrm{\Pi }(k^2)=\frac{g^2}{2}\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{(4\pi )^2}\frac{1}{s}\text{e}^{\text{i}m^2s}\frac{1}{2}\underset{1}{\overset{1}{}}𝑑v\text{e}^{\text{i}k^2\frac{s}{4}(1v^2)}.$$
(47)
An integration by parts on the variable $`v`$ then produces
$$\mathrm{\Pi }(k^2)=\frac{g^2}{2(4\pi )^2}\underset{0}{\overset{\mathrm{}}{}}\frac{ds}{s}\text{e}^{\text{i}m^2s}+\frac{g^2}{2(4\pi )^2}\frac{k^2}{2}\underset{0}{\overset{1}{}}𝑑vv^2\frac{1}{\left[m^2+\frac{k^2}{4}(1v^2)\right]}.$$
(48)
Now let us assume that the $`\varphi `$-particle is massless; then we can regularize $`\mathrm{\Pi }(k^2)`$ at $`k^2=0`$ and construct $`\mathrm{\Pi }^\text{R}(k^2)=\mathrm{\Pi }(k^2)\mathrm{\Pi }(0)`$, which is equivalent to writing Eq. (48) in the simple form (we drop R again):
$$\mathrm{\Pi }(k^2)=\frac{g^2}{2(4\pi )^2}\frac{k^2}{2}\underset{0}{\overset{1}{}}𝑑vv^2\frac{1}{\left[m^2+\frac{k^2}{4}(1v^2)\right]}.$$
(49)
The same procedure can be carried through in Eq. (47). There we would need
$`{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{s}}{\displaystyle \underset{0}{\overset{1}{}}}𝑑v\text{e}^{\text{i}m^2s}\left(\text{e}^{\text{i}s\frac{k^2}{4}(1v^2)}1\right)={\displaystyle \underset{0}{\overset{1}{}}}𝑑v{\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{ds}{s}}\left(\text{e}^{\text{i}s\left[m^2+\frac{k^2}{4}(1v^2)\right]}\text{e}^{\text{i}m^2s}\right)`$
$`\stackrel{(\text{11})}{=}{\displaystyle \underset{0}{\overset{1}{}}}𝑑v\mathrm{ln}{\displaystyle \frac{m^2+\frac{k^2}{4}(1v^2)}{m^2}}={\displaystyle \underset{0}{\overset{1}{}}}𝑑v\mathrm{ln}\left(1+{\displaystyle \frac{k^2}{4m^2}}(1v^2)\right)`$
and so obtain another version of $`\mathrm{\Pi }(k^2)`$, namely,
$$\mathrm{\Pi }(k^2)=\frac{g^2}{(4\pi )^2}\frac{1}{2}\underset{0}{\overset{1}{}}𝑑v\mathrm{ln}\left(1+\frac{k^2}{4m^2}(1v^2)\right).$$
(50)
Substituting $`x=\frac{1+v}{2}`$ into Eq. (50) we finally arrive at
$`\mathrm{\Pi }(k^2)`$ $`=`$ $`{\displaystyle \frac{g^2}{2(4\pi )^2}}{\displaystyle \frac{1}{2}}{\displaystyle \underset{1}{\overset{1}{}}}𝑑v\mathrm{ln}\left(1+{\displaystyle \frac{k^2}{4m^2}}(1v^2)\right)`$ (51)
$`=`$ $`{\displaystyle \frac{g^2}{2(4\pi )^2}}{\displaystyle \underset{0}{\overset{1}{}}}𝑑x\mathrm{ln}\left(1+{\displaystyle \frac{k^2}{m^2}}x(1x)\right).`$
It is also instructive to derive the spectral representation of our vacuum polarization diagram. To do this we start from Eq. (49) and substitute $`v=\left(1\frac{4m^2}{M^2}\right)^{1/2}`$. Our end result is then presented as
$$\mathrm{\Pi }(k^2)=k^2\underset{(2m)^2}{\overset{\mathrm{}}{}}𝑑M^2\frac{\sigma (M^2)}{k^2+M^2\text{i}ϵ},$$
(52)
where the so-called spectral measure is given by
$$\sigma (M^2)=\frac{1}{2}\left(\frac{g}{4\pi }\right)^2\frac{\left(1\frac{4m^2}{M^2}\right)^{1/2}}{M^2}.$$
(53)
Finally we write for the modified massless $`\varphi `$-particle propagator $`\overline{\mathrm{\Delta }}_+^\varphi `$:
$$k^2\left[1k^2\underset{(2m)^2}{\overset{\mathrm{}}{}}𝑑M^2\frac{\sigma (M^2)}{k^2+M^2\text{i}ϵ}\right]\overline{\mathrm{\Delta }}_+^\varphi (k)=1$$
(54)
or
$`\overline{\mathrm{\Delta }}_+^\varphi (k)`$ $`=`$ $`{\displaystyle \frac{1}{k^2\text{i}ϵ}}{\displaystyle \frac{1}{1k^2\underset{(2m)^2}{\overset{\mathrm{}}{}}𝑑M^2\frac{\sigma (M^2)}{k^2+M^2\text{i}ϵ}}}`$ (55)
$`=`$ $`{\displaystyle \frac{1}{k^2\text{i}ϵ}}+{\displaystyle \underset{(2m)^2}{\overset{\mathrm{}}{}}}𝑑M^2{\displaystyle \frac{\sigma (M^2)}{k^2+M^2\text{i}ϵ}},`$ (56)
where we expanded the second factor of Eq. (55). We have hereby reproduced the Lehmann-Källen spectral representation:
$$\overline{\mathrm{\Delta }}_+^\varphi (k)=\frac{1}{k^2\text{i}ϵ}+\frac{1}{2}\left(\frac{g}{4\pi }\right)^2\underset{(2m)^2}{\overset{\mathrm{}}{}}\frac{dM^2}{M^2}\left(1\frac{4m^2}{M^2}\right)^{1/2}\frac{1}{k^2+M^2\text{i}ϵ}.$$
(57)
It is interesting to compare expressions Eq. (53) and Eq. (57) with those occurring in scalar QED. There we would find
$$\overline{\mathrm{\Delta }}_{+\mu \nu }^\gamma =\left(g_{\mu \nu }\frac{k_\mu k_\nu }{k^2}\right)\overline{\mathrm{\Delta }}_+^\gamma (k^2),\sigma (M^2)=\frac{1}{3}\left(\frac{e}{4\pi }\right)^2\frac{\left(1\frac{4m^2}{M^2}\right)^{3/2}}{M^2},$$
(58)
$$\overline{\mathrm{\Delta }}_+^\gamma (k^2)=\frac{1}{k^2\text{i}ϵ}+\frac{1}{3}\left(\frac{e}{4\pi }\right)^2\underset{(2m)^2}{\overset{\mathrm{}}{}}\frac{dM^2}{m^2}\left(1\frac{4m^2}{M^2}\right)^{3/2}\frac{1}{k^2+M^2\text{i}ϵ}.$$
(59)
Incidentally, Eq. (51) corresponds to the scalar QED-case:
$$\mathrm{\Pi }(k^2)=\left(\frac{e}{4\pi }\right)^2\underset{0}{\overset{1}{}}𝑑x(2x1)^2\mathrm{ln}\left[1+\frac{k^2}{m^2}x(1x)\right].$$
(60)
## 3 Conclusion
We have presented a one-loop calculation for a simplified model field theory which is based on standard quantum mechanical path integrals. We found that loop-momentum integrals can be avoided and be replaced with simple Feynman path integrals. All the results known from ordinary field theory can thus be obtained with much less labor. Since our calculations have great similarity with those occurring in QED, the reader should now be able to pursue his or her own studies in more realistic models.
## Acknowledgements
I benefited from discussions with R. Shaisultanov, who provided the stimulus for this investigation. I also would like to thank H. Gies for carefully reading the manuscript.
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# 1 Introduction
## 1 Introduction
Ideas of electromagnetic duality lead to a dramatic breakthrough in our understanding of the dynamics of strongly coupled supersymmetric gauge theories. Particularly spectacular results were obtained by Seiberg and Witten in $`𝒩=2`$ supersymmetry where the low energy effective Lagrangians were found exactly .
One of the most important physical outcomes of the Seiberg–Witten theory is the demonstration of confinement of charges via the monopole condensation. The scenario for confinement as a dual Meissner effect was proposed by Mandelstam and ’t Hooft many years ago . However, because the dynamics of monopoles is hard to control in non-supersymmetric gauge theories this picture of confinement remained as an unjustified qualitative scheme.
The breakthrough in this direction was made by Seiberg and Witten . Using the holomorphy imposed by $`𝒩=2`$ supersymmetry they showed that the condensation of monopoles to really occurs near the monopole point on the modular space of the theory once $`𝒩=2`$ supersymmetry is broken down to $`𝒩=1`$ one by the mass term for the adjoint matter.
In this talk I will review the confinement scenario in the Seiberg–Witten theory underlining the basic features of this $`U(1)`$ confinement which distinguish it from those we expect in QCD.
In particular, I am going to discuss extra hadron states arising in the Seiberg–Witten theory which we do not expect in QCD or in $`𝒩=1`$ supersymmetric QCD.
Then I focus on Abrikosov–Nielsen–Olesen (ANO) strings which are responsible for the confinement near the monopole point. We will see that these strings turn out to be too ”thick” and cannot be described by the standard string theory approximation of long and infinitely thin strings. The result of this is that Regge trajectories become linear only at very large values of spin $`j`$.
In the second part of my talk I will review another confinement scenario arising in the Seiberg–Witten theory with the fundamental matter: confinement on Higgs branches . Higgs branch represents a limiting case of the superconductor of type I with vanishing Higgs mass. We will see that in this limit ANO vortices becomes logarithmically ”thick” . Because of this the confining potential is not linear any longer. It behaves as $`L/\mathrm{log}L`$ with the distance between heavy trial charges (monopoles). This peculiar confining regime can occur only in supersymmetric theories.
In the end of my talk I will speculate on the possible ways to avoid at least some of the unwanted features of $`U(1)`$ confinement I am discussing in this talk.
## 2 Confinement as a dual Meissner effect
First let me remind the mechanism of confinement suggested by Mandelstam and ’t Hooft . Consider an Abelian–Higgs model with the action
$$S_{AH}=d^4x\left\{\frac{1}{4g^2}F_{\mu \nu }^2+|_\mu \phi |^2+\lambda (|\phi |^2v^2)^2\right\}.$$
(2.1)
Here $`\phi `$ is a complex scalar field, $`_\mu =_\mu in_eA_\mu `$, where $`n_e`$ is the electric charge of $`\phi `$. We assume the weak coupling $`g^21`$. The scalar field in (2.1) develop VEV
$$|\phi |=v,$$
(2.2)
which breaks $`U(1)`$ gauge group. The photon acquires the mass
$$m_\gamma ^2=2n_e^2g^2v^2,$$
(2.3)
while the mass of the Higgs boson (one real degree of freedom) is
$$m_H^2=4\lambda v^2.$$
(2.4)
Now introduce an infinitely heavy trial monopole and anti-monopole in the vacuum of the Abelian Higgs model. We can think of them as of Dirac monopoles or as of ’t Hooft–Polyakov monopoles of some underlying non-Abelian gauge theory broken down to $`U(1)`$.
Monopoles has quantized values of magnetic flux $`2\pi n/n_e`$. This magnetic flux cannot be absorbed by the vacuum once there are no dynamical magnetic charges in the theory. On the other hand, magnetic field cannot penetrate into the Higgs vacuum. As a result ANO vortex appears connecting monopole with anti-monopole. This vortex can be viewed as a bubble of an unbroken vacuum inside the Higgs vacuum.
It has $`\phi =0`$ along the line connecting monopole and anti-monopole and magnetic flux $`2\pi n/n_e`$, $`n`$ — winding number. ANO vortex is a solution of equations of motion for the Abelian–Higgs model (2.1). It corresponds to a non-trivial map from the infinite circle in the plane orthogonal to the axis of the vortex to the gauge group, $`\pi _1(U(1))=Z`$. For example, for the winding number $`n=1`$ fields $`\phi `$ and $`A_\mu `$ behave at the infinity as
$`\phi ve^{i\alpha },`$
$`A_m\epsilon _{mn}{\displaystyle \frac{x_n}{x^2}},`$ (2.5)
where $`x_n`$, $`n=1,2`$ is the distance from the vortex axis in the plane orthogonal to this axis, while $`\alpha `$ is the polar angle in this plane.
Because the vortex has a fixed energy per unit length (string tension $`T`$) the potential between monopole and anti-monopole at large distances $`L`$ behaves as
$$V(L)=TL.$$
(2.6)
This linear potential means confinement of monopoles.
The ratio of the Higgs mass (2.4) to the photon mass (2.3) is an important parameter characterizing the type of superconductor. If $`m_H>m_\gamma `$ we have the type II superconductor. Different vortices interact via repulsive forces. In particular, the tension of the vortex with winding number $`n`$ is larger than the sum of tensions of $`n`$ vortices with winding numbers $`n=1`$. Therefore, vortices with higher winding numbers, $`n>1`$ are unstable.
In particular, for $`m_Hm_\gamma `$ (London limit) the vortex solution can be found in the analytic form. The string tension for this case was calculated by Abrikosov in 1957 and later on re-obtained by Nielsen and Olesen in 1973 in the framework of the relativistic field theory. The result for $`n=1`$ is
$$T=2\pi v^2\mathrm{ln}\frac{m_H}{m_\gamma }.$$
(2.7)
For the particular interesting case $`m_H=m_\gamma `$ vortices satisfy the first order equations. They saturate the Bogomolny bound and their string tension reads
$$T_n=2\pi v^2n,$$
(2.8)
where $`n`$ is the winding number. In particularly, as is clear from (2.8), different vortices do not interact.
In supersymmetric theories the BPS-saturation means that some of SUSY generators act trivially on the vortex solution . The Bogomolny bound (2.8) coincides with the central charge of SUSY algebra. This means that the classical result (2.8) for the vortex string tension remains exact in the quantum theory.
For $`m_H<m_\gamma `$ we have type I superconductor. In this case vortices attract each other. In particular, vortices with higher $`n`$ are stable . For the case $`m_Hm_\gamma `$ the vortex solution can be found analytically . The string tension in this case has the form
$$T_n=\frac{2\pi v^2}{\mathrm{ln}m_\gamma /m_H},$$
(2.9)
and does not depend on $`n`$ to the leading order in $`\mathrm{ln}m_\gamma /m_H`$.
In the general case the string tension $`T_n`$ is a monotonic function of the ratio $`m_H/m_\gamma `$ . It reaches its Bogomolny bound (2.8) at $`m_H=m_\gamma `$.
To conclude this section let me point out that the main lesson to learn here is that the condensation of electric charges cause the confinement of monopoles. Vice versa, the condensation of monopoles in the dual Abelian Higgs model leads to the confinement of electric charges.
## 3 Monopole condensation
Now consider $`𝒩=2`$ gauge theory. Most of all in this talk I will be talking about the simplest case of the theory with $`SU(2)`$ gauge group studied in . The simplest version of this theory contains one $`𝒩=2`$ vector multiplet. This multiplet on the component level consists of the gauge field $`A_\mu ^a`$, two Weyl fermions $`\lambda _1^{\alpha a}`$ and $`\lambda _2^{\alpha a}`$ $`(\alpha =1,2)`$ and the complex scalar $`\phi ^a`$, where $`a=1,2,3`$ is the color index.
The scalar potential of this theory has a flat direction. Thus the scalar field can develop an arbitrary VEV along this direction breaking $`SU(2)`$ gauge group down to $`U(1)`$. We choose $`\phi ^a=\delta ^{a3}a`$. The complex parameter $`a`$ parameterize the moduli space of the theory (Coulomb branch). The low energy effective theory contains only the photon $`A_\mu =A_\mu ^3`$ and its superpartners: two Weyl fermions $`\lambda _1^3`$, $`\lambda _2^3`$ and the complex scalar $`a`$.
The Coulomb branch can be parameterized by the gauge invariant parameter $`u=\phi ^{a2}/2`$. It has two singular points $`u=\pm 2\mathrm{\Lambda }^2`$ ( $`\mathrm{\Lambda }`$ is the scale of gauge theory <sup>2</sup><sup>2</sup>2We use the Pauli-Villars regularization scheme here.) where monopole or dyon becomes massless. Near, say, the monopole point $`(u=2\mathrm{\Lambda }^2)`$ the effective low energy theory is dual $`𝒩=2`$ QED. This means that the theory has light monopole hypermultiplet interacting with the dual photon multiplet in the same way as ordinary charges interact with the ordinary photon. The action of this dual QED reads
$$S_{eff}=S_g^{eff}+S_m^{eff}.$$
(3.1)
Here the action for the gauge field is
$$S_g^{eff}=d^4x\{d^2\theta d^2\overline{\theta }\frac{1}{g^2}\overline{A}_DA_D+d^2\theta \frac{1}{4g^2}W_D^2+c.c.\},$$
(3.2)
where $`A_D`$ is the dual chiral $`𝒩=1`$ field. Its lowest component $`a_D`$ goes to zero at the monopole point. $`W_D^\alpha `$ $`(\alpha =1,2`$) is the $`𝒩=1`$ chiral field of the dual photon field strength. Together $`A_D`$ and $`W_D`$ form $`𝒩=2`$ vector U(1) multiplet.
The matter-dependent part of the action reads
$`S_m^{eff}={\displaystyle d^4xd^2\theta d^2\overline{\theta }\left[\overline{M}e^VM+\overline{\stackrel{~}{M}}e^V\stackrel{~}{M}\right]}+`$
$`+`$ $`i{\displaystyle d^4xd^2\theta \sqrt{2}\stackrel{~}{M}A_DM}+c.c..`$ (3.3)
Here $`M,\stackrel{~}{M}`$ are two chiral fields of the monopole hypermultiplet. The monopole mass (given by $`m_m^2=2|a_D|^2`$) goes to zero at the monopole point $`a_D=0`$.
Now let us break $`𝒩=2`$ QED (3.1) down to $`𝒩=1`$ adding the mass term for the adjoint matter in the microscopic SU(2) Seiberg–Witten theory
$$S_{\mathrm{mass}}=id^4xd^2\theta \mu \mathrm{\Phi }^{a2}+c.c.$$
(3.4)
where $`\mathrm{\Phi }^a`$ is the $`𝒩=1`$ chiral superfield which contains component fields $`\phi ^a`$ and $`\lambda _2^{\alpha a}`$. Expressed in terms of $`A_D`$ near the monopole point in the effective theory it reads
$$\mu \mathrm{\Phi }^{a2}=\sqrt{2}\xi A_D+\frac{\mu _D}{2}A_D^2+O(A_D^3),$$
(3.5)
where
$`\xi `$ $`=`$ $`2i\mu \mathrm{\Lambda },`$
$`\mu _D`$ $`=`$ $`{\displaystyle \frac{27}{4}}\mu .`$ (3.6)
The coefficients in (3.6) can be read off the Seiberg–Witten exact solution .
Minimizing the superpotential in (3.1),(3.5) with respect to $`M,\stackrel{~}{M}`$ and $`A_D`$ we find that the Coulomb branch shrinks to the point
$$a_D=0,$$
while
$$\stackrel{~}{m}m=\xi ,$$
(3.7)
where $`m,\stackrel{~}{m}`$ are the scalar components of $`M,\stackrel{~}{M}`$. Taking into account the $`D`$-term condition in (3.1)
$$|m|=|\stackrel{~}{m}|$$
(3.8)
we get the monopole condensate
$$|m|=|\stackrel{~}{m}|=\sqrt{|\xi |}.$$
(3.9)
The monopole condensation breaks the U(1) gauge group and ensures confinement of electric charges. ANO strings arising as a result of the monopole condensation connect quarks with anti-quarks (we interpret these states as mesons) or with another quarks with opposite electric charge (we interpret these states as baryons).
Note, that the effective dual QED is in the weak coupling regime at small $`\mu `$. The QED coupling behaves as
$$\frac{8\pi ^2}{g^2}\mathrm{log}\frac{\mu }{\mathrm{\Lambda }}.$$
(3.10)
Moreover at $`\mu \mathrm{\Lambda }`$ we can ignore non-Abelian effects. Note, that $`W`$-boson mass $`m_W^2=2|a|^2`$ is of order of $`\mathrm{\Lambda }`$ at the monopole point, $`m_W\mathrm{\Lambda }`$.
Parameters $`\xi `$ and $`\mu _D`$ in the mass term perturbation (3.5) play quite different role in the effective QED description of the theory. As it is noted in the linear term in (3.5) is the Fayet–Iliopoulos $`F`$-term which do not break $`𝒩=2`$ supersymmetry. I will explain this in more details in the next section.
On the contrary, the mass term for $`A_D`$ in (3.5) proportional to $`\mu _D`$ breaks $`𝒩=2`$ supersymmetry because it shifts the mass of $`A_D`$ away from the photon mass.
## 4 The U(1) confinement versus the QCD–like confinement
Now let me discuss several basic features of the confinement near the monopole point in the Seiberg–Witten theory at small $`\mu `$ and contrast them to those we expect in QCD. It is believed that non-supersymmetric Yang-Mills theory is in the same universality class as $`𝒩=1`$ Yang-Mills theory. The latter can be obtained as a large $`\mu `$ limit of the theory under consideration. The reason is that in this limit the adjoint matter becomes heavy and decouples. Therefore, we also expect the QCD-like confinement at large $`\mu `$ in the theory at hand. Unfortunately we have no control over the theory in this limit.
Now I will show that the $`U(1)`$ confinement in the theory at small $`\mu `$ has several important distinctions from what we expect from the theory at large $`\mu `$.
### 4.1 Higher winding numbers
The first problem I would like to talk about arises already in SU(2) theory . As we discussed before the flux of the ANO vortex is given by its winding number $`n`$, which is an element of $`\pi _1(U(1))=Z`$ (for SU$`(N_c)`$ group it is $`\pi _1(U^{N_c1}(1))=Z^{N_c1})`$. This could produce an extra multiplicity in the hadron spectrum which we do not expect in QCD. In QCD or in the large $`\mu `$ limit of the present theory we expect classification of states under the center of the gauge group, $`Z_2`$ for SU(2) rather than $`Z`$ ($`Z_{N_c}`$ for SU$`(N_c)`$).
Consider as an example ANO vortex with $`n=2`$ in the effective dual QED at small $`\mu `$. This string connects two quarks with two anti-quarks producing an ”exotic” state. Note, that the string with $`n=2`$, in principle, can be broken by $`W`$-boson pair creation but at low energies we can neglect this process because $`W`$-boson is too heavy at small $`\mu `$, $`\mu \mathrm{\Lambda }`$ , ($`m_W\mathrm{\Lambda }`$).
The discussion above of ”exotic” states in the hadron spectrum is based on the purely topological reasoning. Now let us consider the dynamical side of the problem. As we have seen in Sect.2, strings with higher $`n`$ are stable or unstable depending on the type of the superconductor. Namely, in type I superconductor the energy of the vortex with winding number $`n`$ is less than the total energy of $`n`$ vortices with winding numbers $`n=1`$. Therefore vortices with $`n>1`$ are stable and we really have an ”exotic” states in the spectrum.
On the contrary, in type II superconductor vortices with $`n>1`$ are unstable against decay into $`n`$ vortices with winding numbers $`n=1`$. Therefore, ”exotic” states are unstable and actually in the “real world” at strong coupling might be not observable at all.
Thus the natural question arises: what is the type of superconductivity at the monopole point in the Seiberg–Witten theory. This problem is studied in . Let me briefly present the result.
To the leading order in $`\mu /\mathrm{\Lambda }`$ we ignore the mass term for $`A_D`$ in (3.5) and consider only the linear in $`A_D`$ term parameterized by $`\xi `$. As I mentioned before this term is the generalized Fayet–Illiopoulos (FI) term. Let me explain what does this mean. Let us start with $`𝒩=1`$ QED. In $`𝒩=1`$ supersymmetric U(1) gauge theory one can add FI term to the action (we call it FI $`D`$-term here)
$$\xi _3D,$$
(4.1)
where $`D`$ is the $`D`$-component of the gauge field. In $`𝒩=2`$ SUSY theory field $`D`$ belongs to the $`SU_R(2)`$ triplet together with $`F`$-components of the field $`A_D`$, $`F_D`$ and $`\overline{F}_D`$. Namely, let us introduce the triplet $`F_a`$ $`(a=1,2,3)`$ using relations
$`D`$ $`=`$ $`F_3,`$
$`F_D`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(F_1+iF_2)],`$
$`\overline{F}_D`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}(F_1iF_2).`$ (4.2)
Now the generalized FI-term can be written as
$$S_{FI}=id^4x\xi _aF_a.$$
(4.3)
Comparing this with (3.5) we identity
$`\xi `$ $`=`$ $`{\displaystyle \frac{1}{2}}(\xi _1i\xi _2),`$
$`\overline{\xi }`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\xi _1+i\xi _2).`$ (4.4)
For this reason we call the term linear in $`A_D`$ in (3.5) FI $`F`$-term.
It is well known that the $`𝒩=1`$ QED with FI $`D`$-term has BPS ANO string . The reason for this is the following. After the breakdown of $`U(1)`$ gauge group the photon acquires the mass given by (2.3) $`(v^2=2|\xi |`$ in (2.3), see (3.9)). The massive vector $`𝒩=1`$ supermultiplet contains, in particular, one real scalar. This scalar plays the role of the Higgs boson in the Abelian Higgs model (2.1). Thus the BPS condition $`m_H=m_\gamma `$ is imposed by supersymmetry.
Now return to $`𝒩=2`$ supersymmetry and consider dual QED (3.1) with the FI $`F`$-term added. The FI parameter $`\xi ^a`$ explicitly breaks the $`SU_R(2)`$ group. However, the $`𝒩=2`$ supersymmetry remains unbroken. To see this note that FI term is proportional to $`F`$ and $`D`$ components of the vector multiplet which transform as a total derivatives under the $`𝒩=2`$ supersymmetry transformation.
It is clear that the FI $`F`$-term which appears in the Seiberg–Witten theory (see (3.5)) can be obtained by $`SU_R(2)`$ rotation from the FI $`D`$-term. Moreover, masses of particles in $`𝒩=2`$ multiplet do not change under this rotation. In particular, the condition $`m_H=m_\gamma `$ stays intact. Hence, the ANO string is BPS-saturated . Its string tension is given by (2.8). For $`n=1`$
$$T=2\pi (2|\xi |),$$
(4.5)
where we use that $`v^2=2|\xi |`$.
To consider the next-to-leading correction in $`\mu /\mathrm{\Lambda }`$ we switch on the mass term for $`A_D`$ parameterized by $`\mu _D`$ in the effective theory, see (3.5) . It breaks the $`𝒩=2`$ supersymmetry and split the $`𝒩=2`$ multiplet. In particular, it breaks the BPS condition $`m_H=m_\gamma `$. The effect of this term is studied in . The result is that the Higgs mass in the effective Abelian Higgs model appears to be less than the photon mass, $`m_H<m_\gamma `$ and the theory is driven to the type I superconductivity. The string tension is less than its Bogomolny bound,
$$T<2\pi (2|\xi |).$$
(4.6)
In particular, in the large $`\mu _D`$-limit $`\mu _D^2\xi `$ the string tension is found analytically <sup>3</sup><sup>3</sup>3To take this limit we ignore relations (3.6), (3.10) and consider QED (3.1) with the prturbation (3.5) on its own right viewing parameters $`\xi `$ , $`\mu _D`$ and $`g^2`$ as independent ons, assuming only the weak coupling condition $`g^21`$.:
$$T=\frac{2\pi (2|\xi |)}{\mathrm{ln}(\frac{g|\mu _D|}{2\sqrt{|\xi |}})}.$$
(4.7)
As I explained above the result in (4.6) means that we really have an infinite tower of ”exotic” hadron states with higher string fluxes near the monopole point at small $`\mu `$.
Of course, if we increase $`\mu `$ and approach $`\mu \mathrm{\Lambda }`$ the unwanted strings with $`n>1`$ become broken by the $`W`$-bosons production. The only string with $`n=1`$ connecting one quark with one anti-quark<sup>4</sup><sup>4</sup>4Or with another quark with the opposite electric charge. will probably survive. This string is believed to be responsible for the confinement in $`𝒩=1`$ SQCD at large $`\mu `$ . It is called QCD string. This picture of confinement is proposed in refs. within the brane approach.
However, this scenario is hard to implement in the field theory. One reason for this is that at large $`\mu `$ dual QED enters the strong coupling regime and is no longer under control. Another one is probably even more fundamental. The point is that the role of matter fields in the effective QED (3.1) is played by monopoles. As $`\mu `$ approaches $`\mathrm{\Lambda }`$ the inverse mass of the dual photon (2.3) approaches the size of monopole (which is of order of the inverse $`W`$-boson mass $`m_W^1\mathrm{\Lambda }^1`$). Under these conditions we hardly can consider monopoles as a local degrees of freedom and the dual QED effective description breaks down. In particular, we don’t have a field theoretical description of $`n=1`$ string in the region of large $`\mu \mathrm{\Lambda }`$.
### 4.2 p–strings
Another problem of $`U(1)`$ confinement in Seiberg–Witten theory was noticed by Douglas and Shenker . It appears in $`SU(N_c)`$ gauge theories at $`N_c3`$.
Scalar VEV’s breaks the gauge group down to $`U(1)^{N_c1}`$. Hence, $`N_c1`$ different types of ANO vortices arise each one associated with a particular $`U(1)`$ factor. Their string tensions are given by
$$T_p\xi \mathrm{sin}\frac{\pi p}{N_c},$$
(4.8)
where $`p=1,\mathrm{},N_c1`$ numerates different $`U(1)`$ factors.. They are called $`p`$-strings. The number of strings with different string tensions equals to $`[(N_c1)/2]`$.
Therefore it is clear that extra states emerge in the hadron spectrum which we do not expect in QCD . Namely, the number of quark–anti-quark meson states is $`N_c`$. Let me explain this for the case of $`SU(3)`$.
We take three different quarks of $`SU(3)`$ as
$$q_1=\left(\begin{array}{c}1\hfill \\ 0\hfill \\ 0\hfill \end{array}\right),q_2=\left(\begin{array}{c}0\hfill \\ 1\hfill \\ 0\hfill \end{array}\right),q_3=\left(\begin{array}{c}0\hfill \\ 0\hfill \\ 1\hfill \end{array}\right)$$
(4.9)
and choose generators of two $`U(1)`$ groups of the broken $`SU(3)`$ to be
$$\left(\begin{array}{ccc}1& 0& 0\\ 0& 1& 0\\ 0& 0& 0\end{array}\right)\text{ and }\left(\begin{array}{ccc}0& 0& 0\\ 0& 1& 0\\ 0& 0& 1\end{array}\right).$$
(4.10)
Thus, three quarks in (4.9) have the following electric charges with respect to two $`U(1)`$ groups:
$$(1,0);(1,1);(0,1).$$
(4.11)
Now it is clear that $`\stackrel{~}{q}_1q_1`$ meson is formed by 1-string, $`\stackrel{~}{q}_2q_2`$ meson is formed by 2-string and 1-anti-string and $`\stackrel{~}{q}_3q_3`$ mesons is formed by 2-anti-string. In sum we have 3 meson states. Two of them ($`\stackrel{~}{q}_1q_1`$ and $`\stackrel{~}{q}_3q_3`$) are (classically) degenerative while the third one, $`\stackrel{~}{q}_2q_2`$ is two times heavier than the first two <sup>5</sup><sup>5</sup>5This is true if we neglect masses of quarks at the ends of strings (see subsection 4.4 for the discussion of this point).
In general we have $`[(N_c+1)/2]`$ families of $`\stackrel{~}{q}q`$ mesons with different masses . Each family contains two mesons (one family contains one meson if $`N_c`$ is odd). They are classically degenerative but can split in quantum theory.
Let me illustrate this splitting for the simplest example of $`SU(2)`$ gauge group. For $`SU(2)`$ we have two classically degenerative mesons. If we decompose the quark as
$$q=q_+\left(\genfrac{}{}{0pt}{}{1}{0}\right)+q_{}\left(\genfrac{}{}{0pt}{}{0}{1}\right),$$
(4.12)
then these two mesons are $`\stackrel{~}{q}_+q_+`$ and $`\stackrel{~}{q}_{}q_{}`$. Out of these states we can form two different combinations as follows
$`\stackrel{~}{q}q=\stackrel{~}{q}_+q_++\stackrel{~}{q}_{}q_{},`$
$`\stackrel{~}{q}\tau _3q=\stackrel{~}{q}_+q_+\stackrel{~}{q}_{}q_{}.`$ (4.13)
Here $`\tau _3`$ is the color matrix. In gauge invariant notation two states in (4.13) are
$$\stackrel{~}{q}q,\frac{1}{\sqrt{\phi ^2}}\stackrel{~}{q}\phi q.$$
(4.14)
It is clear that classically these states are degenerative but in quantum theory they split. The heavier state acquires a large width and might not be observable at all.
Still we have $`[(N_c+1)/2]`$ different $`\stackrel{~}{q}q`$ meson states instead of one state we expect in the large $`\mu `$ limit.
As it is noted in $`W`$-bosons are charged under different $`U(1)`$ factors and therefore provide a coupling between them. Thus, we expect that at large $`\mu `$, $`\mu \mathrm{\Lambda }`$ most of $`\stackrel{~}{q}q`$-mesons discussed above become unstable and disappear. In fact, it is shown in within the brane approach that as we increase $`\mu `$ all $`\stackrel{~}{q}q`$-meson disappear except one in which quark is connected with anti-quark by the 1-string. This string is shown to become QCD string of ref. at large $`\mu `$.
Unfortunately, as I mentioned before, we still have no description of this transition in field theory.
$`P`$-strings by themselves also survive the large $`\mu `$\- limit . They connect $`p`$ quarks in the $`p`$-index antisymmetric representation with $`p`$-anti-quarks to form an ”exotic” meson. Also they can form baryons.
To conclude this subsection I would like to mention the recent paper in which $`𝒩=2`$ supersymmetry breaking with the first two Casimir operators is considered for the $`SU(N_c)`$ theory. It is claimed that once off-diagonal couplings between different $`U(1)`$ factors are taken into account $`p`$-strings in generic case fail to be BPS saturated even in the limit of zero $`\mu _D`$.
### 4.3 Non-linear Regge trajectories
One of the main motivations to consider hadrons as quarks connected by strings in early days of String Theory was the linear behavior of Regge trajectories. However, as I show now the ANO string at the monopole point of the Seiberg–Witten theory does not produce linear Regge trajectories. More precisely Regge trajectories become linear only at very large $`j`$.
Let me first show the linear behavior of Regge trajectories using very elementary classical arguments.
Consider long ANO string (of the length $`L`$) connecting light quark and anti-quark <sup>6</sup><sup>6</sup>6See next subsection for the discussion on whether quarks can be light near the monopole point.. Let this string to rotate around the axis orthogonal to the string with large spin $`j`$. If $`j`$ is large enough the problem becomes classical and we can apply classical equations of motion to the rotating string. For the linear potential
$$V(L)=TL$$
(4.15)
equation of motion looks like
$$TE\omega ^2L,$$
(4.16)
where $`E`$ is the mass of the string (we neglect the quark masses)
$$EV(L)=TL,$$
(4.17)
while $`\omega `$ is the frequency of the angular rotation. Expressing $`\omega `$ in terms of $`j`$ using $`jE\omega L^2`$ we obtain
$$L^2\frac{j}{T}.$$
(4.18)
Using (4.17) again we get that the meson mass is
$$E^2Tj.$$
(4.19)
The mass square is proportional to $`j`$ at large $`j`$.
More precisely the spectrum of a free Nambu–Goto string goes as
$$E^2=E_0^2+T(j+n).$$
(4.20)
Here $`E_0^2`$ is the intercept and $`n`$ labels the daughter trajectories. At large $`j`$, $`j1`$ our naive estimate (4.19) gives the same result as the exact string spectrum (4.20).
The string tension $`T`$ is given by (4.5) to the leading order in $`\mu `$. Expressed in terms of the photon mass (2.3) (or Higgs mass, $`m_\gamma m_H`$) it reads
$$T=\frac{\pi m_\gamma ^2}{n_e^2g^2}.$$
(4.21)
We see that in terms of photon mass $`T`$ is large in the weak coupling $`g^21`$. This is a typical result for solitonic objects in the semiclassical approximation.
Now let us discuss the region of validity of (4.19) and (4.20). The string theory result (4.20) assumes the approximation of long and thin strings. The transverse size of the ANO vortex is given by $`1/m_\gamma `$. So we need
$$L^2\frac{1}{m_\gamma ^2}.$$
(4.22)
Substituting (4.18) and (4.21) here we get
$$j\frac{1}{g^2}.$$
(4.23)
We see that in fact Regge trajectories become linear only at extremely large $`j`$.
The bound (4.23) is rather restrictive. If $`j`$ is not that large the transverse size of the vortex becomes of order of its length and the string is not developed. Quark and anti-quark are on the border between the stringy regime and the Coulomb regime. At small $`j`$ $`\stackrel{~}{q}q`$-meson more looks like spherically symmetric soliton rather than a string.
We can also see the breakdown of the string description from the string representation for the ANO vortex. It is developed in and for the cases of strings in the type II superconductor, BPS-strings and strings in the type I superconductor respectively. The common feature of these representations is that the leading term of the world sheet action is the Nambu–Goto term
$$S_{\mathrm{string}}=Td^2x\left\{\sqrt{g}+\text{ higher derivatives }\right\},$$
(4.24)
where $`g_{ij}=_ix_\mu _jx_\mu `$ is the induced metric $`(i,j=1,2)`$. Higher derivative corrections in (4.24) contain term important for the string quantization , rigidity term etc. These terms contain powers of $`/m_\gamma `$. For thin strings $`m_\gamma m_H\mathrm{}`$ and these corrections can be neglected in the action (4.24). However, for the ANO vortex in the semiclassical regime $`^2/m_\gamma ^2T/m_\gamma ^21/g^21`$ (see (4.21)). Hence, higher derivative corrections blow up in (4.24) and the string approximation is no longer acceptable. From the string theory point of view this manifest itself as a ”crumpled” string surface .
We see that Regge trajectories are not linear in the wide region of spins $`j\begin{array}{c}<\hfill \\ \hfill \end{array}1/g^2`$. This to be contrasted with perfect linear behavior of Regge trajectories in Nature starting from small $`j`$ <sup>7</sup><sup>7</sup>7 On the contrary, in QCD is hard to talk about linear trajectories at large $`j`$ because higher resonances acquire large widths. (for a recent account see lecture given at this School).
It might seem funny that we are trying to compare some properties of Seiberg–Witten theory with experiment. Still I believe that the linear behavior of Regge trajectories is an important feature of confinement in QCD and we have to reproduce it in a theory with QCD-like confinement.
The main property responsible for this ”disadvantage” of the confinement in the monopole point of the Seiberg–Witten theory is the large value of the string tension (4.21). In the conclusion of this talk I will speculate that this problem as well as some others could be resolved if the string were (almost) tensionless.
### 4.4 Heavy quark-anti-quark states
Another unpleasant consequence of the large value of the string tension (4.21) is that hadrons built of quarks appear to be too heavy.
To show this let me summarize qualitatively the low lying hadron spectrum in our theory. First, there are states with color (magnetic) charge screened by the Higgs mechanism. They are monopoles (described by operators $`\stackrel{~}{M}M`$) and the dual photon with its superpartners. Their masses are of order of $`m_\gamma `$ $`(m_Hm_\gamma `$ at small $`\mu `$).
Second, there are hadrons built of quarks via the confinement mechanism described above. As an example, consider $`\stackrel{~}{q}q`$-meson at $`j1`$. Its mass is of order of
$$m_{\stackrel{~}{q}q}2m_q+\sqrt{T},$$
(4.25)
where $`m_q`$ is the quark mass.
The problem is that the mass in (4.25) is too large (as compared with $`m_\gamma `$) as I will show below. This means that we have light monopole states and the photon (which we can interpret as glueballs), while states built of quarks are heavy. In contrast, in QCD we have light $`\stackrel{~}{q}q`$-states, whereas the candidates for glueballs are much heavier.
First, let us see how small $`m_q`$ in (4.25) can be. So far we discussed the pure gauge theory or the theory with very heavy quarks which we used as a probe for the confinement. Now let us introduce one dynamical flavor of the fundamental matter hypermultiplet (we call it quark).
The quark mass on the Coulomb branch of $`𝒩=2`$ theory is given by
$$m_q=m+\frac{a}{\sqrt{2}},$$
(4.26)
where $`m`$ is the quark mass parameter in the microscopic theory. Quarks become light near the charge singular point on the Coulomb branch, $`a=\sqrt{2}m`$. Hence, in order to have light quarks we have to go near the charge singularity.
On the other hand, once we switch on $`\mu `$ the Coulomb branch shrinks to three singular points: monopole, dyon and charge ones. As we discuss before, in order to have monopole condensation and quark confinement we have to go to the monopole point.
Now to make quarks light near the monopole point we choose $`m`$ to ensure that the charge singularity goes close on the Coulomb branch to the monopole one. The values of $`m`$ corresponding to the colliding of these singularities are called Argyres–Douglas (AD) points . In $`SU(2)`$ theory these points were studied in . On the Coulomb branch of $`𝒩=2`$ theory these points flow in the infrared to a non-trivial conformal field theories.
The value of $`m`$ at which the monopole singularity collides with the charge one is (there are three such points; we choose one of them corresponding to real $`m_{AD}`$)
$$m_{AD}=\frac{3}{2^{4/3}}\mathrm{\Lambda }_1,$$
(4.27)
where $`\mathrm{\Lambda }_1`$ is the scale of the theory with one flavor, $`\mathrm{\Lambda }^4=m\mathrm{\Lambda }_1^3`$ at large $`m`$ .
The problem, however, is that we cannot go directly to the AD-point (4.27). The point is that the monopole condensate vanishes at the AD-point . This means that the AD-point is the point of the quark deconfinement .
Now the question is whether we can go close to the AD-point to make quarks light enough without destroying the monopole condensation and quark confinement. To answer this question let us develop a perturbation theory around AD-point.
In general, the monopole condensate in the theory with one flavor is
$$\stackrel{~}{M}M=2i\mu \left(u_m^22m\mathrm{\Lambda }_1^3\right)^{1/4},$$
(4.28)
where $`u_m`$ is the position of the monopole singularity in the $`u`$-plane given by the Seiberg–Witten curve . In particular, in the AD-point $`u_m^{AD}=\frac{4}{3}m_{AD}^2`$. This value makes the monopole condensate in (4.28) vanish.
Now let us take $`m`$ close to its AD-value
$$m=m_{AD}(1+\epsilon ),$$
(4.29)
where $`\epsilon 1`$. Then extracting $`u_m`$ from the Seiberg–Witten curve near AD point we get
$$\stackrel{~}{M}M\mu \mathrm{\Lambda }_1\epsilon ^{1/4}.$$
(4.30)
The value of the monopole condensate (4.30) set the scale of our effective Abelian Higgs model. In particular, the photon mass and ANO string tension is given by Eqs. (2.3) and (2.8), where the Higgs VEV $`v^2`$ is identified with the monopole condensate (4.30).
Now let us see if we can make quark mass in (4.25) to be small as compared to the string scale $`T^{1/2}\stackrel{~}{M}M^{1/2}`$. According to ref. the anomalous dimension of $`m_q`$ in (4.26) is one, while the anomalous dimension of $`mm_{AD}`$ is $`4/5`$. Using this we conclude that
$$m_q\mathrm{\Lambda }_1\epsilon ^{5/4}.$$
(4.31)
From (4.30) and (4.31) we see that we can always make quarks lighter than the string scale $`T^{1/2}`$ if we take $`\mu `$ not too small, $`\mu \mathrm{\Lambda }_1\epsilon ^{9/4}`$. Of course, we still keep $`\mu \mathrm{\Lambda }_1`$ to ensure the validity of our dual QED description.
Unfortunately, this does not solve the problem of heavy $`\stackrel{~}{q}q`$-states . As we already mentioned, the string scale is much larger or at least of the same order as the photon mass, (see (4.21)). Therefore, the mass of $`\stackrel{~}{q}q`$-meson in (4.25) is much larger or of the same order as the photon and monopole masses in contrast with our expectations about the theory with QCD-like confinement.
## 5 Confinement on Higgs branches
In this section I will consider another confinement scenario in the Seiberg–Witten theory: confinement on the Higgs branch in the theory with matter. I will talk about $`SU(2)`$ theory with $`N_f=2`$ flavors of fundamental matter . The $`N_f=2`$ case is the simplest case of the theory with Higgs branches (for $`N_f=1`$ we do not have Higgs branches).
### 5.1 Review of Higgs branches in SU(2) theory
Let us introduce $`N_f=2`$ fundamental matter hypermultiplets in the $`𝒩=2`$ SU(2) gauge theory. In terms of $`𝒩=1`$ superfields matter dependent part of the microscopic action looks like
$`S_{\mathrm{matter}}={\displaystyle d^4xd^2\theta d^2\overline{\theta }\left[\overline{Q}_Ae^VQ^A+\overline{\stackrel{~}{Q}}^Ae^V\stackrel{~}{Q}_A\right]}+`$
$`+`$ $`i{\displaystyle d^4xd^2\theta \left[\sqrt{2}\stackrel{~}{Q}_A\frac{\tau ^a}{2}Q^A\mathrm{\Phi }^a+m\stackrel{~}{Q}_AQ^A\right]}+\text{ c.c.}`$ (5.1)
Here $`Q^{kA},\stackrel{~}{Q}_{Ak}`$ are matter chiral fields, $`k=1,2`$ and $`A=1,\mathrm{},N_F`$, while $`V`$ is the vector superfield. Thus we have 16 real matter degrees of freedom for $`N_f=2`$.
Consider first the limit of large $`m`$. In this limit the three singularities on the Coulomb branch are easy to understand. Two of them correspond to monopole and dyon singularities of the pure gauge theory. Their positions on the Coulomb branch are given by
$$u_{m,d}=\pm 2m\mathrm{\Lambda }_2\frac{1}{2}\mathrm{\Lambda }_2^2,$$
(5.2)
where $`u=\frac{1}{2}\phi ^{a^2}`$ and $`\mathrm{\Lambda }_2`$ is the scale of the theory with $`N_f=2`$. In the large $`m`$ limit $`u_{m,d}`$ are approximately given by their values in the pure gauge theory $`u_{m,d}\pm 2m\mathrm{\Lambda }_2=\pm 2\mathrm{\Lambda }^2`$, where $`\mathrm{\Lambda }`$ is the scale of $`N_f=0`$ theory.
The third singularity corresponds to the point where charge becomes massless. Let us decompose matter fields as
$$Q^{kA}=\left(\genfrac{}{}{0pt}{}{1}{0}\right)^kQ_+^A+\left(\genfrac{}{}{0pt}{}{0}{1}\right)^kQ_{}^A.$$
(5.3)
From the superpotential in (5.1) we see that the $`Q_+`$ becomes massless at
$$a=\sqrt{2}m.$$
(5.4)
The singular point $`a=+\sqrt{2}m`$ is gauge equivalent to the one in (5.4). In terms of variable $`u`$ (5.4) reads
$$u_c=m^2+\frac{1}{2}\mathrm{\Lambda }^2.$$
(5.5)
Strictly speaking, we have $`2+N_f=4`$ singularities on the Coulomb branch. However two of them coincides for the case of two flavors of matter with the same mass.
The effective theory on the Coulomb branch near charge singularity (5.4) is given by $`𝒩=2`$ QED with light matter fields $`Q_+^A`$, $`\stackrel{~}{Q}_{+A}`$ (8 real degrees of freedom) as well as the photon multiplet.
The charge singularity (5.4),(5.5) is the root of the Higgs branch . To find this branch let us write down $`D`$-term and $`F`$-term conditions which follow from (5.1). $`D`$-term conditions are
$$Q^{kA}\overline{Q}_A\mathrm{}+\overline{\stackrel{~}{Q}}^{kA}\stackrel{~}{Q}_A\mathrm{}=0,$$
(5.6)
while $`F`$-term conditions give (5.4) as well as
$$Q^{kA}\stackrel{~}{Q}_A\mathrm{}=0.$$
(5.7)
Eqs. (5.6),(5.7) have nontrivial solutions for $`N_f2`$. These solutions determines VEV’s for scalar components $`q^{kA}`$, $`\stackrel{~}{q}_{Ak}`$ of fields $`Q^{kA},\stackrel{~}{Q}_{Ak}`$. Dropping heavy components $`q_{}`$ according to decomposition (5.3) and introducing the SU$`{}_{R}{}^{}(2)`$ doublet $`q^{fA}`$ as
$`q^{1A}=q_+^A,`$ $`q^{2A}=\overline{\stackrel{~}{q}}_+^A,`$
$`\overline{q}_{A1}=\overline{q}_A^+,`$ $`\overline{q}_{A2}=\stackrel{~}{q}_A^+,`$ (5.8)
we can rewrite three real conditions in (5.6),(5.7) as
$$\overline{q}_{Ap}(\tau ^a)_f^pq^{fA}=0,a=1,2,3.$$
(5.9)
Eq.(5.9) together with the condition (5.4) determines the Higgs branch (manifold with $`q0`$) which touches the Coulomb branch at the point (5.4).
The low energy theory for boson fields near the root of the Higgs branch looks like
$$S_{\mathrm{boson}}^{\mathrm{root}}=d^4x\{\frac{1}{4g^2}F_{\mu \nu }^2+\overline{}_\mu \overline{q}_{Af}_\mu q^{fA}+\frac{g^2}{8}[\text{ Tr }\overline{q}\tau ^aq)]^2\},$$
(5.10)
where trace is calculated over flavor and SU$`{}_{R}{}^{}(2)`$ indices. Here $`_\mu =_\mu in_eA_\mu `$, $`\overline{}_\mu =_\mu +in_eA_\mu `$, the electric charge $`n_e=1/2`$ for fundamental matter fields.
This is an Abelian Higgs model with last interaction term coming from the elimination of $`D`$ and $`F`$ terms. The QED coupling constant $`g^2`$ is small near the root of the Higgs branch. We include 8 real matter degrees of freedom $`q^{fA}`$ in the theory (5.10) according to the identification (5.8). The rest of matter fields $`q_{}^A`$, $`\stackrel{~}{q}_A^{}`$ (another 8 real degrees of freedom) acquire a large mass $`2m`$ and can be dropped out. The effective theory (5.10) is correct on the Coulomb branch near the root of the Higgs branch (5.4) or on the Higgs branch not far away from the origin $`q=0`$.
It is clear that the last term in (5.10) is zero on the fields $`q`$ which satisfy constraint (5.9). This means that moduli fields which develop VEV’s on the Higgs branch are massless, as it should be. The other fields acquire mass of order $`\overline{q}q^{1/2}`$. It turns out that there are four real moduli fields $`q`$ (out of 8) which satisfy the constraint (5.9) . They correspond to the lowest components of one hypermultiplet.
We can parameterize them as
$$q^{f\dot{A}}(x)=\frac{1}{\sqrt{2}}\sigma _\alpha ^{f\dot{A}}\varphi _\alpha (x)e^{i\alpha (x)}.$$
(5.11)
Here $`\varphi _\alpha (x),\alpha =1\mathrm{}4`$ are four real moduli fields. It is clear that fields (5.11) solve (5.9). The common phase $`\alpha (x)`$ in (5.11) is the U(1) gauge phase. Once $`\varphi _\alpha =v_\alpha 0`$ on the Higgs branch the U(1) group is broken and $`\alpha (x)`$ is eaten by the Higgs mechanism. Say, in the unitary gauge $`\alpha (x)=0`$. In the next subsection we consider vortex solution for the model (5.10). Then $`\alpha (x)`$ is determined by the behavior of the gauge field at the infinity. Substituting (5.11) into (5.10) we get the bosonic part of the effective theory for the massless moduli fields on the Higgs branch near the origin
$$S_{\mathrm{boson}}^{\mathrm{Higgs}}=d^4x\left\{\frac{1}{4g^2}F_{\mu \nu }^2+\overline{}_\mu \overline{q}_\alpha _\mu q_\alpha \right\},$$
(5.12)
where
$$q_\alpha (x)=\varphi _\alpha (x)e^{i\alpha (x)}.$$
(5.13)
Once $`v_\alpha 0`$ we expect monopoles (they are heavy at $`m\mathrm{\Lambda }`$) to confine via formation of vortices which carry the magnetic flux. The peculiar feature of the theory (5.12) is the absence of the Higgs potential. Therefore, the Higgs phase of of the theory in (5.12) is the limiting case of type I superconductor with the vanishing Higgs mass. In the next subsection I will consider the peculiar features of ANO vortices in the model (5.12).
If we increase $`v_\alpha ^2`$ taking $`v_\alpha ^2\begin{array}{c}>\hfill \\ \hfill \end{array}\mathrm{\Lambda }^2`$ we can integrate out massive photon. Then the effective theory is a $`\sigma `$-model for massless fields $`q_\alpha `$ which belong to 4-dimensional Hyper–Kahler manifold, $`R^4/Z_2`$. The metric of this $`\sigma `$-model is flat , there are, however, higher derivative corrections induced by instantons . Here we consider region of Higgs branch with $`v_\alpha ^2\mathrm{\Lambda }_2`$. This determines the scale of the effective Abelian Higgs model (5.12). $`W`$-bosons and other particles which reflect the non-Abelian structure of the underlying microscopic theory are heavy with masses $`\begin{array}{c}>\hfill \\ \hfill \end{array}\mathrm{\Lambda }_2`$ and can be ignored.
To conclude this subsection let us briefly review what happens if we reduce the mass parameter $`m`$. At $`m=\pm \mathrm{\Lambda }`$ the charge singularity (root of the Higgs branch) collides with the monopole (dyon) singularity, see Eqs.(5.2),(5.5). These are Argyres–Douglas points . At these points mutually non-local degrees of freedom (say, charges and monopoles) becomes massless simultaneously. These points are very interesting from the point of view of the monopole confinement on the Higgs branch. Monopoles become dynamical as we approach Argyres–Douglas point, $`m\mathrm{\Lambda }_2`$.
After the collision quantum numbers of particles at singularities change because of monodromies . If we denote quantum numbers as $`(n_m,n_e)_B`$, where $`n_m`$ and $`n_e`$ are magnetic and electric charges of the state, while $`B`$ is its baryon number then at $`m>\mathrm{\Lambda }_2`$ we have charge, monopole and dyon singularities with quantum numbers
$$(0,1/2)_1^2,(\mathrm{1\; 0})_0,(1,1)_0.$$
(5.14)
The superscript for the charge means that we have two flavors of charges. After charge singularity collides with monopole one (at $`m<\mathrm{\Lambda }_2)`$ the quantum numbers of particles at singularities become
$$(1,0)_0^2,(1,1/2)_1,(1,1/2)_1.$$
(5.15)
Now monopole $`(1,0)_0`$ condense on the Higgs branch which emerges from the point (5.5), while dyons $`(1,1/2)_1`$ and $`(1,1/2)_1`$ confine because they carry electric charge. At zero mass, $`m=0`$ two dyon singularities in (5.15) coincide (see (5.2)) and the second Higgs branch appears at the point $`u=1/2\mathrm{\Lambda }_2^2`$. This restores the global symmetry from SU$`(N_f=2)`$ in the massive theory to SO$`(2N_f=4)`$ at $`m=0`$ .
### 5.2 ANO string on the Higgs branch
Now let us focus on confinement of monopoles on the Higgs branch and consider ANO vortices in the model (5.12) <sup>8</sup><sup>8</sup>8 Let us take $`m>\mathrm{\Lambda }_2`$ to avoid confusion with notation of dyon states.. This is done in .
Without loss of generality we take VEV’s of $`q_\alpha `$ $`v_\alpha =(v,0,0,0)`$. Moreover, we drop fields $`q_2,q_3`$ and $`q_4`$ from (5.12) because they are irrelevant for the purpose of finding classical vortex solutions. Thus, we arrive at the Abelian Higgs model (2.1) with $`\lambda =0`$ and identification $`q_1=\phi `$.
Following consider first the model (2.1) with small $`\lambda `$, so that $`m_Hm_\gamma `$ (see (2.3), (2.4) for $`n_e=1/2`$). Then we take the limit $`m_H0`$.
To the leading order in $`\mathrm{log}m_\gamma /m_H`$ the vortex solution has the following structure . The electromagnetic field is confined in a core with the radius
$$R_g^2\frac{1}{m_\gamma ^2}\mathrm{ln}^2\frac{m_\gamma }{m_H}.$$
(5.16)
The scalar field is close to zero inside the core. Instead, outside the core, the electromagnetic field is vanishingly small. At intermediate distances
$$R_gr\frac{1}{m_H}$$
(5.17)
($`r`$ is the distance from the center of vortex in (1,2) plane) the scalar field satisfy the free equation of motion. Its solution reads
$$\phi (r)=v\left(1\frac{\mathrm{ln}1/rm_H}{\mathrm{ln}1/R_gm_H}\right).$$
(5.18)
At large distances $`r1/m_H`$ $`\phi `$ approaches its VEV as $`\phi v\mathrm{exp}(m_Hr)`$.
The main contribution to the string tension comes from the logarithmically large region (5.17), where scalar field is given by (5.18). The result for the string tension is (as we already mentioned, see (2.9))
$$T=\frac{2\pi v^2}{\mathrm{ln}m_\gamma /m_H}.$$
(5.19)
It comes from the kinetic energy of the scalar field in (2.1) (”surface” energy).
The results in (5.16),(5.19) mean that if we naively take the limit $`m_H0`$ the string becomes infinitely thick and its tension goes to zero . This means that there are no strings in the limit $`m_H=0`$. The ANO vortex becomes a vacuum state with ”twisted” boundary conditions which ensure the magnetic flux $`2\pi /n_e`$. <sup>9</sup><sup>9</sup>9The author is indebted to A. Vainshtein for this interpretation. The absence of ANO strings in theories with flat Higgs potential was noticed in . ANO strings on the Higgs branch was also discussed in from the brane point of view. However it was not noticed in that the vortex becomes infinitely ”thick” and its tension goes to zero.
One might think that the absence of ANO strings means that there is no confinement on Higgs branches. As we will see now this is not the case .
So far we have considered infinitely long ANO strings. However the setup for the confinement problem is slightly different . We have to consider monopole–anti-monopole pair at large but finite separation $`L`$. Our aim is to take the limit $`m_H0`$. To do so let us consider ANO string of the finite length $`L`$ within the region
$$\frac{1}{m_\gamma }L\frac{1}{m_H}.$$
(5.20)
Then it turns out that $`1/L`$ plays the role of the $`IR`$-cutoff in Eqs. (5.16) and (5.19) instead of $`m_H`$ .
The reason for this is easy to understand. In fact, the reason for the absence of the vortex solution is quite clear from (5.18). If we put $`m_H`$ to zero the scalar field $`\phi `$ cannot reach its VEV at infinity because of its logarithmic behavior. This was noticed in .
Now consider vortex of finite length $`L`$ and look at the behavior of $`\phi `$ at large distances $`r`$, $`Lr1/m_H`$. In this region the problem becomes three dimensional. The solution of the free equation of motion in three dimensions reads
$$\phi v\frac{1}{|x|},$$
(5.21)
where $`|x|`$ is three dimensional distance from the vortex. The solution (5.21) perfectly goes to its VEV $`v`$ at infinity. To put it in the other way at distances $`|x|L`$ from the vortex the two dimensional problem transforms into the three dimensional one. Namely, the scalar field has logarithmic behavior at $`r`$ within the bounds $`R_grL`$, whereas at $`|x|L`$ it acquires $`1/|x|`$ behavior given by (5.21). We see that $`L`$ really plays the role of the $`IR`$-cutoff for the logarithmic behavior of $`\phi `$. Now we can take the limit $`m_H0`$.
The result for the electromagnetic core of the vortex becomes
$$R_g^2\frac{1}{m_\gamma ^2}\mathrm{ln}^2m_\gamma L,$$
(5.22)
while its string tension is given by
$$T=\frac{2\pi v^2}{\mathrm{ln}m_\gamma L}.$$
(5.23)
We see that the ANO string becomes ”thick” but still its transverse size $`R_g`$ is much less than its length $`L`$, $`R_gL`$. As a result the potential between heavy well separated monopole and anti-monopole is still confining but is no longer linear in $`L`$. It behaves as
$$V(L)=2\pi v^2\frac{L}{\mathrm{ln}m_\gamma L}.$$
(5.24)
As soon as the potential $`V(L)`$ is an order parameter which distinguishes different phases of a theory (see, for example, review ) we conclude that we have a new confining phase on the Higgs branch of the Seiberg–Witten theory. It is clear that this phase can arise only in supersymmetric theories because we do not have Higgs branches without supersymmetry.
Unfortunately, the confinement on Higgs branches cannot play a role of a model for QCD-like confinement. In particular, the confining potential (5.24) gives rise to the following behavior of Regge trajectories at ultra-large $`j`$, $`j1/g^2`$ (see (4.18))
$$E^2v^2\frac{j}{\mathrm{ln}(g^2j)}.$$
(5.25)
We see that Regge trajectories never becomes linear. At $`j\begin{array}{c}<\hfill \\ \hfill \end{array}1/g^2`$ they are non-linear due to the reason we have discussed in subsection 4.3. At $`j1/g^2`$ they are still non-linear because of the logarithmic factor in (5.25).
It is worth note also that because of type I superconductivity on the Higgs branch we have a tower of “exotic” hadron states corresponding to higher winding numbers of strings.
## 6 Outlook: tensionless strings
We have discussed two confinement scenarios in the Seiberg–Witten theory: the confinement of quarks in the monopole point upon breaking $`𝒩=2`$ down to $`𝒩=1`$ and the confinement of monopoles on the Higgs branch in $`𝒩=2`$ theory. The latter one does not looks like confinement in QCD, while the former one is more promising.
Still as we have discussed in Section 4 it has several unwanted features. Basically these unwanted features fall into two categories. First, is the presence of extra states in the hadron spectrum. It is a reflection of $`U(1)`$ nature of the confinement in Seiberg–Witten theory. The second is associated with the large value of the string tension (4.21) as compared with the squire of the photon or Higgs mass.
One can believe that the first problem disappears if we increase $`\mu `$ and go to the limit of $`𝒩=1`$ QCD. Now I would like to speculate that the second problem also might disappear in the same limit.
One might think that once we increase $`\mu `$ the dual QED coupling $`g^2`$ increases and becomes large, $`g^21`$. If this happens the string tension becomes small, see (4.21). As a result Regge trajectories becomes linear already at $`j1`$. Moreover, hadrons built of quarks becomes light and can be visible in the low energy spectrum. Instead monopoles and photon (to be interpreted as glueballs) become relatively heavy. Of course, at strong coupling we loose control over the theory and cannot use eq. (4.21) relating the string tension to the photon mass. The string is not a BPS state and its tension is not protected from corrections. Moreover, we do not have exact formula for the photon mass in the strong coupling as well. Still we can speculate that there is a region in the parameter space where the string tension becomes much less then the squire of the photon mass.
In fact, the conjecture that ANO strings might become tensionless in the strong coupling limit of QED was suggested in Ref. in order to find the field theory explanation of tensionless $`M`$-theory strings which arise when 5-branes approach each other.
Note, that we have to find the regime such that $`T`$ goes to zero, while $`m_\gamma `$ and $`m_H`$ (which control the transverse size of ANO string) stay finite.
For example, the string on a Higgs branch discussed in the previous section does not do the job. Its tension goes to zero but its transverse size $`R_g`$ becomes infinite (see (5.16)). Thus this string disappears. Another example is the string in the monopole point at the $`AD`$-value of the quark mass (see subsection 4.4). As the monopole condensate vanishes the string tension goes to zero, however the string transverse size might become infinite. It is not clear if this string can serve as a QCD string.
To have a QCD-like confinement we need a regime in which string remains a string (its size does not grow) while its tension becomes small.
### Acknowledgments
The author is grateful to the organizers of XXXIV PNPI Winter school for the opportunity to present this lecture. Also the author would like to thank A. Hanany, M. Shifman, M. Strassler, and, in particularly, A. Vainshtein for useful discussions. The author is also grateful to the Theoretical Physics Institute at the University of Minnesota for support. This work is also supported by Russian Foundation for Basic Research under grant No.99-02-16576.
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# Moment-angle complexes and combinatorics of simplicial manifolds
## 1. Introduction
The classical Dehn–Sommerville equations for simplicial convex $`n`$-dimensional polytope $`P^n`$ give a set of linear relations among the numbers $`f_i`$ of $`i`$-dimensional faces of $`P^n`$. The integer vector $`f(P^n)=(f_0,f_1,\mathrm{},f_{n1})`$ is called the $`f`$-vector of $`P^n`$. We put $`f_1=1`$. The Dehn–Sommerville equations were established by Dehn for $`n5`$, and by Sommerville in the general case in 1927 (see \[So\]). Their original form is as follows
(1)
$$f_k=\underset{j=k}{\overset{n1}{}}(1)^{n1j}\left(\genfrac{}{}{0pt}{}{j+1}{k+1}\right)f_j,k=0,\mathrm{},n1.$$
Define the $`h`$-vector from the equation
(2)
$$h_0t^n+\mathrm{}+h_{n1}t+h_n=(t1)^n+f_0(t1)^{n1}+\mathrm{}+f_{n1}.$$
Obviously, the $`f`$-vector and the $`h`$-vector determine each other by means of linear equations (note that $`h_0=1`$). For instance,
(3)
$$h_k=\underset{i=0}{\overset{k}{}}(1)^{ki}\left(\genfrac{}{}{0pt}{}{ni}{ki}\right)f_{i1}.$$
The notion of $`h`$-vector gives rise to the simplest and the most elegant form of the Dehn–Sommerville equations (1) for simplicial polytopes:
(4)
$$h_i=h_{ni},i=0,\mathrm{},n.$$
The Dehn–Sommerville equations can be generalized quite widely. In \[Kl\] V. Klee reproved the Dehn–Sommerville equations in the form (1) in a more general context of Eulerian manifolds. In particular, it turns out that equations (1) hold for any simplicial manifold (i.e. triangulated topological manifold) of dimension $`n1`$. Analogues of equations (1) were obtained by Bayer and Billera \[BaBi\] (for Eulerian partially ordered sets) and Chen and Yan \[CY\] (for arbitrary polyhedra).
It follows directly from (4) that the affine hull of $`f`$-vectors $`(f_0,\mathrm{},f_{n1})`$ of simplicial polytopes in $`n`$-dimensional space is an $`\left[\frac{n}{2}\right]`$-dimensional plane. The same is true for the affine hull of vectors $`(b_0,b_1,\mathrm{},b_n)`$ of Betti numbers $`b_i=dimH_i(M)`$ of orientable connected (simplicial) $`n`$-dimensional manifolds $`M^n`$ due to the Poincaré duality:
(5)
$$b_i(M^n)=b_{ni}(M^n),i=0,\mathrm{},n.$$
This similarity between $`f`$-vectors of simplicial polytopes and Betti numbers of orientable manifolds was pointed out by Klee in \[Kl\]. Moreover, it turns out that the “combinatorial” duality (4) can be interpreted in terms of the “topological” duality (5). Given an $`n`$-dimensional simplicial polytope $`P`$ (or, more generally, a complete simplicial fan $`\mathrm{\Sigma }`$), one may construct the toric variety $`M_P`$ (or $`M_\mathrm{\Sigma }`$) of (real) dimension $`2n`$ (see \[Da\]\[Fu\]). This variety is not necessarily a manifold, however its homology (with rational coefficients in general) satisfies the Poincaré duality, and its even Betti numbers equal the components of $`h`$-vector of $`P`$: $`b_{2i}(M_P)=h_i(P)`$. This gives a “topological” proof of the Dehn–Sommerville equations (4).
The dual (or polar) to any simplicial polytope $`P`$ is a simple polytope, which we denote $`P^{}`$. Given an $`n`$-dimensional simple polytope $`P^{}`$ with $`m`$ facets, Davis and Januszkiewicz defined in \[DJ\] a manifold $`𝒵_P^{}`$ of dimension $`m+n`$ acted on by the torus $`T^m`$. This manifold depends only on the combinatorial type (i.e., the face lattice) of a polytope, and the orbit space for the $`T^m`$-action is combinatorially $`P^{}`$. The manifolds $`𝒵_P^{}`$ establish the bridge between topology of manifolds and combinatorics of polytopes (or more general objects, such as simplicial spheres). Various connections between topology of $`𝒵_P^{}`$ and toric geometry, symplectic geometry, subspace arrangements, and combinatorial theory of $`f`$-vectors were studied in \[BP1\], \[BP2\], \[BP3\]. Any smooth toric variety (or Hamiltonian $`T^n`$-manifold) defined by a simple polytope of combinatorial type $`P^{}`$ is the quotient of $`𝒵_P^{}`$ by a freely acting toric subgroup $`T^{mn}T^m`$. This is also true for topological analogues of smooth toric varieties, which we call quasitoric manifolds, also introduced in \[DJ\]. For more information about quasitoric manifolds and their topology see \[BaBe\], \[BR1\], \[BR2\], \[Ma\], \[Pa1\], \[Pa2\]. We showed in \[BP2\], \[BP3\] that the cohomology algebra of $`𝒵_P^{}`$ with coefficients in any field $`𝐤`$ is isomorphic to the $`Tor`$-algebra $`Tor_{𝐤[v_1,\mathrm{},v_m]}(𝐤(P),𝐤)`$, where $`P`$ is polar to $`P^{}`$, $`m=f_0`$ is the number of vertices of $`P`$, and $`𝐤(P)`$ is the Stanley–Reisner face ring of $`P`$. The Koszul resolution gives then a very simple model $`H[𝐤(P)\mathrm{\Lambda }[u_1,\mathrm{},u_m],d]`$ for the cohomology algebra of $`𝒵_P^{}`$. In particular, the cohomology of $`𝒵_P^{}`$ acquires a bigraded algebra structure, and the bigraded Poincaré duality implies the Dehn–Sommerville equations (4).
The boundary complex of a simplicial polytope is a simplicial sphere. However, now it is well known that not any triangulation of a topological sphere can be obtained in such way. First examples of such “non-polytopal” spheres were found by Grünbaum, and the smallest non-polytopal sphere is of dimension 3 with 8 vertices (the so-called Barnette sphere, see \[Ba\]). However, the Dehn–Sommerville relations in the form (4) still hold for any simplicial sphere (see \[St2, § II.6\]). In \[BP4\] we extend our approach to manifolds $`𝒵_P^{}`$ defined by simple polytopes to the case of arbitrary simplicial complex $`K`$. We endow the cone $`cone(K)`$ with a structure of cubical complex and embed it into the boundary complex of $`m`$-dimensional unit cube $`I^m`$ (here $`m=f_0(K)`$ is the number of vertices of $`K`$). Then we view $`I^m`$ as the orbit space for the diagonal action of the torus $`T^m=S^1\times \mathrm{}\times S^1`$ on the unit poly-disk $`(D^2)^m`$ in $`m`$-dimensional complex space $`^m`$. Hence, the above cubical embedding $`cone(K)I^m`$ is covered by a $`T^m`$-equivariant embedding $`𝒵_K(D^2)^m`$, where $`𝒵_K`$ is a cellular complex canonically decomposed into the union of blocks $`(D^2)^n\times T^{mn}`$ with $`n=dimK+1`$. We call this $`𝒵_K`$ the moment-angle complex associated to the simplicial complex $`K`$. If $`K`$ is a simplicial sphere, then $`𝒵_K`$ is a manifold. (If, moreover, $`K`$ is polytopal, i.e. $`K=P`$, then $`𝒵_K`$ coincides with the above described manifold $`𝒵_P^{}`$, where $`P^{}`$ is polar to $`P`$.) However, in general, $`𝒵_K`$ has more complicated structure. In \[BP4\] we showed that our moment-angle complex $`𝒵_K`$ is homotopy equivalent to the complement of complex coordinate subspace arrangement defined by $`K`$, and its cohomology algebra is isomorphic, as in the polytopal case, to the $`Tor`$-algebra $`Tor_{𝐤[v_1,\mathrm{},v_m]}(𝐤(K),𝐤)`$. We note that the Betti numbers of the complement of a real coordinate subspace arrangement were calculated in terms of resolution of the Stanley–Reisner ring in \[GPW\]. In the case when $`K`$ is a simplicial sphere, the bigraded Poincaré duality in the cohomology algebra of $`𝒵_K`$ gives a “topological proof” of the Dehn–Sommerville equations (4) for simplicial spheres. In this paper we construct a cellular chain complex that calculates the homology of $`𝒵_K`$ and gives a very transparent characterization of homology classes of $`𝒵_K`$ (and of the complement of a coordinate subspace arrangement). This chain complex is dual to a certain cochain subcomplex of the Koszul complex $`[𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m],d]`$ for the $`Tor`$-algebra $`Tor_{𝐤[v_1,\mathrm{},v_m]}(𝐤(K),𝐤)`$.
From the topological viewpoint it is very interesting to study the case when $`K`$ is a simplicial manifold. The moment-angle complex $`𝒵_K`$ here is not a manifold, however, its singularities are tractable. It turns out that $`𝒵_K`$ contains the product $`|cone(K)|\times T^m`$ of (polyhedron corresponding to) the cone over $`K`$ and torus $`T^m`$. The closure $`W_K=\overline{𝒵_K|cone(K)|\times T^m}`$ of the complement of this singular part is a manifold with boundary $`|K|\times T^m`$. This $`W_K`$ is homotopically equivalent to another moment-angle complex $`𝒲_K`$ which covers equivariantly the restriction to $`Kcone(K)`$ of the cubical embedding $`cone(K)I^m`$. We construct an appropriate cellular decomposition of $`𝒲_K`$, which allows to calculate the homology efficiently. The Poincaré duality for manifolds with boundary gives in this case
$$H_i(W_K)H^{m+ni}(W_K,|K|\times T^m),i=0,\mathrm{},m+n.$$
This duality again regards the bigraded structure. As a consequence, we obtain a topological proof of the Dehn–Sommerville equations (1) for simplicial manifolds in the following nice form:
(6)
$$h_{ni}h_i=(1)^i\left(\chi (K^{n1})\chi (S^{n1})\right)\left(\genfrac{}{}{0pt}{}{n}{i}\right),i=0,1,\mathrm{},n.$$
where $`dimK^{n1}=n1`$, and $`\chi ()`$ denotes the Euler number. We have $`\chi (K^{n1})=f_0f_1+\mathrm{}+(1)^{n1}f_{n1}=1+(1)^{n1}h_n`$ and $`\chi (S^{n1})=1+(1)^{n1}`$. This generalizes equations (4) to the case of arbitrary simplicial manifold. In particular, if $`K`$ is an odd-dimensional simplicial manifold, one has $`h_{ni}=h_i`$. In \[Pa, (7.11)\] Pachner proved by means of his bistellar flip theorem that the value $`h_{ni}h_i`$ is a topological invariant of a PL-manifold (i.e. is independent on a PL-triangulation). The formula (6) calculates this invariant exactly for any simplicial (not necessarily PL) manifold.
Our moment-angle complexes enable to reformulate many combinatorial statements and hypotheses concerning $`f`$-vectors in topological terms. We have already mentioned the Dehn–Sommerville relations for both simplicial spheres and simplicial manifolds, however this is only the first and the simplest example. The most intriguing open problem here is the so-called $`g`$-theorem (or McMullen’s inequalities) for simplicial spheres \[St3\]. It includes the Generalized Lower Bound hypothesis for simplicial spheres, which asserts the monotonicity property for the $`h`$-vector:
(7)
$$h_0h_1h_2\mathrm{}h_{[\frac{n}{2}]}.$$
For polytopal spheres $`g`$-theorem was proved by Stanley \[St1\] (necessity) and Billera, Lee \[BiLe\] (sufficiency) in 1980. The first inequality $`h_0h_1`$ is equivalent to $`1mn`$, which is obvious. The second ($`h_1h_2`$, $`n4`$) is equivalent to the lower bound $`f_1nf_0\left(\genfrac{}{}{0pt}{}{n+1}{2}\right)`$ for the number of edges, which is also known for simplicial spheres. The next inequality $`h_2h_3`$ is still open (for simplicial spheres). For the history of $`g`$-theorem and related questions see \[St2\]\[St3\]\[Zi\]. As we mentioned in \[BP2\], the inequality $`h_1h_2`$ is equivalent to the upper bound $`b_3(𝒵_K)\left(\genfrac{}{}{0pt}{}{mn}{2}\right)`$ for the third Betti number of manifold $`𝒵_K`$ (we note that $`𝒵_K`$ is always 2-connected). Other inequalities from (7) also acquire such topological interpretation, however, in general case in terms of bigraded Betti numbers of $`𝒵_K`$. We hope that such inequalities can be deduced from the equivariant topology of $`𝒵_K`$.
The authors wish to express special thanks to Oleg Musin for stimulating discussions and helpful comments, in particular, for drawing our attention to the results of \[CY\] and \[Pa\].
## 2. Cubical complexes determined by a simplicial complex
Let $`K^{n1}`$ be an $`(n1)`$-dimensional simplicial complex on the vertex set $`[m]=\{1,\mathrm{},m\}`$ (hence, any simplex of $`K`$ has at most $`n`$ vertices). As usual, we view $`K`$ as a set of subsets of $`[m]`$, that is, $`K2^{[m]}`$. We would assume that $`K`$ is not the $`m`$-simplex $`\mathrm{\Delta }^m=2^{[m]}`$, so $`n<m`$. If $`I=\{i_1,\mathrm{},i_k\}[m]`$ is a simplex of $`K`$, then we would write $`IK`$. By definition, the barycentric subdivision of $`K`$ is the simplicial complex $`bs(K)`$ whose simplices are chains $`I_1I_2\mathrm{}I_p`$ of embedded simplices of $`K`$. In particular, the vertices of $`bs(K)`$ are in one to one correspondence with simplices of $`K`$ of all dimensions. We denote the geometric realization of $`K`$ (as a polyhedron) by $`|K|`$. Let $`I^m`$ be the standard unit $`m`$-dimensional cube in $`^m`$:
$$I^m=\{(y_1,\mathrm{},y_m)^m:\mathrm{\hspace{0.33em}0}y_i1,i=1,\mathrm{},m\}.$$
Every face of $`I^m`$ has the form
(8)
$$F_{IJ}=\{(y_1,\mathrm{},y_m)I^m:y_i=0\text{ for }iI,y_j=1\text{ for }jJ\},$$
were $`IJ`$ are two (possibly empty) subsets of $`[m]`$. Now assign to each subset $`I=\{i_1,\mathrm{},i_k\}[m]`$ the vertex $`v_I=F_{II}`$ of the cube $`I^m`$. One has $`v_I=(\epsilon _1,\mathrm{},\epsilon _m)`$, where $`\epsilon _i=0`$ if $`iI`$ and $`\epsilon _i=1`$ otherwise. Viewing $`I`$ as a vertex of the barycentric subdivision $`bs(\mathrm{\Delta }^m)`$ of an $`m`$-simplex, we see that the correspondence $`Iv_I`$ extends to a piecewise linear embedding $`i_c`$ of the polyhedron $`|bs(\mathrm{\Delta }^m)|`$ into the (boundary complex of) unit cube $`I^m`$. Under this embedding, the vertices of $`\mathrm{\Delta }^m`$ are mapped to the vertices $`(1,\mathrm{},1,0,1,\mathrm{},1)`$ of the cube $`I^m`$, while the point in the interior of $`|\mathrm{\Delta }^m|`$ (viewed as a vertex of $`bs(\mathrm{\Delta }^m)`$) is mapped to the vertex $`(0,\mathrm{},0)`$ of $`I^m`$. Hence, the whole image of $`|bs(\mathrm{\Delta }^m)|`$ is the set of $`m`$ facets of $`I^m`$ meeting at the vertex $`(0,\mathrm{},0)`$. Moreover, given a pair $`I`$, $`J`$ of non-empty subsets of $`[m]`$ such that $`IJ`$, all simplices of $`bs(\mathrm{\Delta }^m)`$ of the form $`I=I_1I_2\mathrm{}I_k=J`$ are mapped to the same face $`F_{IJ}`$ of $`I^m`$ (see (8)).
Our simplicial complex $`K^{n1}`$ can be viewed as a subcomplex in $`\mathrm{\Delta }^m`$. Hence, the above constructed map $`i_c:|bs(\mathrm{\Delta }^m)|I^m`$ embeds $`|bs(K)|`$ piecewise linearly into the boundary complex of $`I^m`$. The image $`i_c(|bs(K)|)`$ is a certain $`(n1)`$-dimensional cubical complex, which we denote $`cub(K)`$. Thus, we have proved the following statement.
###### Proposition 2.1.
There is a piecewise linear embedding $`i_c`$ of $`|bs(\mathrm{\Delta }^m)|`$ into the boundary complex of $`I^m`$ such that for any simplicial complex $`K\mathrm{\Delta }^m`$ on $`m`$ vertices the image $`i_c(|bs(K)|)=:cub(K)`$ is the union of faces (8) corresponding to all pairs $`IJ`$ of embedded simplices of $`K`$. $`\mathrm{}`$
###### Remark.
Cubes of the cubical subdivision $`cub(K)`$ of the polyhedron $`|K|`$ are formed by simplices of the barycentric subdivision $`bs(K)`$. This cubical subdivision was employed in some previous researches for different purposes (see, e.g., \[DJ, p. 434\]). The above cubical embedding $`i_c:cub(K)I^m`$ was used previously in \[SS\] to study which cubical complexes can be embedded into the standard cubical lattice.
The map $`i_c:|bs(\mathrm{\Delta }^m)|I^m`$ can be extended to a piecewise linear map $`cone(i_c):|cone(bs(\mathrm{\Delta }^m))|I^m`$ by taking the vertex of the cone to the vertex $`(1,\mathrm{},1)`$ of the cube $`I^m`$. (Note that the cone over the barycentric subdivision of a $`k`$-simplex is identified with the standard triangulation of a $`(k+1)`$-cube.) Now the image of $`|cone(bs(\mathrm{\Delta }^m))|`$ is the whole $`I^m`$, so $`cone(i_c)`$ is a PL homeomorphism. The image $`i_c\left(|cone(bs(K))|\right)I^m`$ is another, this time $`n`$-dimensional, cubical subcomplex of $`I^m`$, which we denote $`cc(K)`$. This cubical complex is explicitly described by the following proposition.
###### Proposition 2.2.
For any simplicial complex $`K`$ on $`m`$ vertices there is a piecewise linear embedding of the polyhedron $`|cone(bs(K))|`$ into the boundary complex of $`I^m`$ such that its image $`cc(K)`$ is the union of faces
(9)
$$F_J=\{(y_1,\mathrm{},y_m)I^m:y_j=1\text{ for }jJ\}I^m$$
and all their subfaces, where $`J`$ ranges over all simplices of $`K`$. $`\mathrm{}`$
According to (8), $`F_J=F_\mathrm{}J`$. Hence, any subface of $`F_J`$ is $`F_{IJ}`$ for some (possibly empty) $`IJ`$. It follows that
$$cub(K)=\underset{I\mathrm{}}{\underset{JK}{}}F_{IJ},cc(K)=\underset{JK}{}F_{IJ}.$$
###### Remark.
Viewed as a topological space, $`|cub(K)|`$ is homeomorphic to $`|K|`$, while $`|cc(K)|`$ is homeomorphic to $`|cone(K)|`$. Viewing $`cone(K)`$ as a simplicial complex, one may construct the cubical complex $`cub\left(cone(K)\right)`$, which is also homeomorphic to $`|cone(K)|`$. However, as cubical complexes, $`cc(K)`$ and $`cub\left(cone(K)\right)`$ differ.
The cubical complex $`cc(K)`$ was introduced in \[BP2\] and studied in \[BP3\], \[BP4\] in connection with simple (and simplicial) polytopes and subspace arrangements.
###### Example 2.3.
The cubical complex $`cub(K)`$ for $`K`$ a disjoint union of 3 vertices and the boundary complex of a 2-simplex is shown on Figure 1 a) and b) respectively. The corresponding cubical complexes $`cc(K)`$ are indicated on Figure 2 a) and b).
## 3. Equivariant moment-angle complexes
Let $`(D^2)^m`$ denote the unit poly-disk in the complex space:
$$(D^2)^m=\{(z_1,\mathrm{},z_m)^m:|z_i|1,i=1,\mathrm{},m\}.$$
The unit cube $`I^m`$ can be viewed as the orbit space for the standard action of $`m`$-dimensional torus $`T^m`$ on $`(D^2)^m`$ by coordinatewise rotations. The orbit map $`\rho :(D^2)^mI^m`$ can be given by $`(z_1,\mathrm{},z_m)(|z_1|^2,\mathrm{},|z_m|^2)`$. For each face $`F_{IJ}`$ of $`I^m`$ (see (8)) define
(10)
$$\begin{array}{c}B_{IJ}:=\rho ^1(F_{IJ})\hfill \\ \hfill =\{(z_1,\mathrm{},z_m)(D^2)^m:z_i=0\text{ for }iI,|z_j|=1\text{ for }jJ\}.\end{array}$$
It follows that if $`\mathrm{\#}I=i`$, $`\mathrm{\#}J=j`$, then $`B_{IJ}(D^2)^{ji}\times T^{mj}`$, where disk factors $`D^2(D^2)^{ji}`$ correspond to elements from $`JI`$, while circle factors $`S^1T^{mj}`$ correspond to elements from $`[m]J`$. Introducing the polar coordinate system in $`(D^2)^m`$, we see that $`B_{IJ}`$ is parametrized by $`(ji)`$ radial (or moment) and $`(mi)`$ angle coordinates. Here we come to the following definition.
###### Definition 3.1.
Let $`C`$ be a cubical subcomplex of $`I^m`$. The moment-angle complex $`ma(C)`$ corresponding to $`C`$ is the $`T^m`$-invariant decomposition of $`\rho ^1(|C|)`$ to “moment-angle” blocks $`B_{IJ}`$ (see (10)) corresponding to faces $`F_{IJ}|C|I^m`$. Hence, $`ma(C)`$ is defined from the commutative diagram
$$\begin{array}{ccc}ma(C)& & (D^2)^m\\ & & \rho & & \\ |C|& & I^m\end{array}.$$
It follows that the torus $`T^m`$ acts on $`ma(C)`$ with orbit space $`|C|`$.
The moment-angle complexes corresponding to the introduced above cubical complexes $`cub(K)`$ and $`cc(K)`$ (see propositions 2.1 and 2.2) will be denoted $`𝒲_K`$ and $`𝒵_K`$ respectively. Thus, we have
(11) $`\begin{array}{ccc}𝒲_K& & (D^2)^m\\ & & \rho & & \\ |cub(K)|& & I^m\end{array}`$ and $`\begin{array}{ccc}𝒵_K& & (D^2)^m\\ & & \rho & & \\ |cc(K)|& & I^m\end{array},`$
where horizontal arrows are embeddings, while vertical ones are orbit maps for certain $`T^m`$-actions. It follows that $`dim𝒵_K=m+n`$ and $`dim𝒲_K=m+n1`$.
Let us consider the cellular decomposition of $`D^2`$ with one 2-cell $`D`$, two 1-cells $`I`$, $`T`$, and two 0-cells 0, 1 (see Figure 3). It defines a $`T^m`$-invariant cellular decomposition of the poly-disk $`(D^2)^m`$ with $`5^m`$ cells. Each cell of this decomposition is the product of cells of 5 different types: $`D_i`$, $`I_i`$, $`0_i`$, $`T_i`$, and $`1_i`$, $`i=1,\mathrm{},m`$. We will encode cells of $`(D^2)^m`$ by “words” of type $`D_II_J0_LT_P1_Q`$, where $`I,J,L,P,Q`$ are pairwise disjoint subsets of $`[m]`$ such that $`IJLPQ=[m]`$. Sometimes we would drop the last factor $`1_Q`$, so in our notations $`D_II_J0_LT_P=D_II_J0_LT_P1_{[m]IJLP}`$. It follows that the closure of $`D_II_J0_LT_P1_Q`$ is homeomorphic to the product of $`\mathrm{\#}I`$ disks, $`\mathrm{\#}J`$ segments, and $`\mathrm{\#}P`$ circles. The constructed cellular decomposition of $`(D^2)^m`$ allows to identify moment-angle complexes as certain cellular subcomplexes in $`(D^2)^m`$.
###### Lemma 3.2.
For any cubical subcomplex $`C`$ of $`I^m`$ the corresponding moment-angle complex $`ma(C)`$ is a $`T^m`$-invariant cellular subcomplex of $`(D^2)^m`$.
###### Proof.
Since $`ma(C)`$ is a union of “moment-angle” blocks $`B_{IJ}`$ (see (10)), and each $`B_{IJ}`$ is obviously $`T^m`$-invariant, the whole $`ma(C)`$ is also $`T^m`$-invariant. In order to show that $`ma(C)`$ is a cellular subcomplex of $`(D^2)^m`$ (with respect to the above constructed cellular decomposition) we just mention that $`B_{IJ}`$ is the closure of cell $`D_{JI}I_{\mathrm{}}0_IT_{[m]J}1_{\mathrm{}}`$. ∎
## 4. Cohomology of $`𝒵_K`$, subspace arrangements, and numbers of faces of $`K`$
Here we study the moment-angle complex $`𝒵_K`$ corresponding to the cubical complex $`cc(K)I^m`$ (see (11)). Remember that $`K`$ is an $`(n1)`$-dimensional simplicial complex on $`m`$ vertices, and $`cc(K)`$ topologically is the cone over $`K`$. By definition, $`𝒵_K`$ is a union of certain moment-angle blocks $`B_{IJ}(D^2)^m`$ with $`I=\mathrm{}`$ (see Proposition 2.2). In analogy with (9), we put
$$B_J:=B_\mathrm{}J=\rho ^1(F_J)=\{(z_1,\mathrm{},z_m)(D^2)^m:|z_j|=1\text{ for }jJ\}.$$
Hence, $`𝒵_K=_{JK}B_J`$, where $`B_J(D^2)^j\times T^{mj}`$, $`j=\mathrm{\#}J`$. This remark allows to simplify the cellular decomposition constructed in the previous section for $`ma(C)`$ in the case when $`ma(C)=𝒵_K`$. To do this we replace the union of cells $`0`$, $`I`$, $`D`$ (see Figure 3) by one 2-dimensional cell. To simplify notations we denote this 2-cell by $`D`$ throughout this section. The resulting $`T^m`$-invariant cellular decomposition of $`(D^2)^m`$ now has $`3^m`$ cells, each of which is the product of 3 different types of cells: $`D_i`$, $`T_i`$, and $`1_i`$, $`i=1,\mathrm{},m`$. In this section we encode these cells of $`(D^2)^m`$ as $`D_IT_P1_Q`$, where $`I,P,Q`$ are pairwise disjoint subsets of $`[m]`$ such that $`IPQ=[m]`$. Usually we would denote the cell $`D_IT_P1_Q`$ just by $`D_IT_P`$, so $`D_IT_P=D_IT_P1_{[m]IP}`$. Since $`B_J=B_\mathrm{}J`$ is the closure of cell $`D_JT_{[m]J}1_{\mathrm{}}`$, the moment-angle complex $`𝒵_K`$ is a $`T^m`$-invariant cellular subcomplex of $`(D^2)^m`$ (with respect to the new $`3^m`$-cell decomposition). The complex $`𝒵_K`$ consists of all cells $`D_IT_P(D^2)^m`$ such that $`I`$ is a simplex of $`K`$.
###### Remark.
Note that for general $`C`$ the moment-angle complex $`ma(C)`$ is not a cellular subcomplex for the $`3^m`$-cell decomposition of $`(D^2)^m`$.
The cohomology ring of $`𝒵_K`$ was described in \[BP2\]\[BP3\] (in the case when $`K`$ is a polytopal sphere) and in \[BP4\] (for general $`K`$). Before going further, we review some results of these papers.
Throughout the rest of this paper we work over some field $`𝐤`$, referred to as the ground field. Let $`𝐤[v_1,\mathrm{},v_m]`$ be the polynomial algebra, and $`\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$ the exterior algebra over $`𝐤`$ on $`m`$ generators. We make both algebras graded by putting $`\mathrm{deg}(v_i)=2`$, $`\mathrm{deg}(u_i)=1`$.
###### Definition 4.1.
The face ring (or the Stanley–Reisner ring) $`𝐤(K)`$ of simplicial complex $`K`$ is the quotient ring $`𝐤[v_1,\mathrm{},v_m]/I`$, where the ideal $`I`$ is generated by all square-free monomials $`v_{i_1}\mathrm{}v_{i_s}`$, $`1i_1<\mathrm{}<i_sm`$, such that $`\{i_1,\mathrm{},i_s\}`$ is not a simplex of K.
For any subset $`I=\{i_1,\mathrm{},i_k\}[m]`$ denote by $`L_I`$ the coordinate plane in $`^m`$ consisting of points whose $`i_1,\mathrm{},i_k`$ coordinates vanish:
(12)
$$L_I=\{(z_1,\mathrm{},z_m)^m:z_{i_1}=\mathrm{}=z_{i_k}=0\}.$$
Each simplicial complex $`K`$ on $`m`$ vertices defines a complex coordinate subspace arrangement $`𝒜(K)`$. The latter is the set of all planes $`L_I`$ such that $`I`$ is not a simplex of $`K`$:
$$𝒜(K)=\{L_I:IK\}.$$
Define the support of $`𝒜(K)`$ as $`|𝒜(K)|=_{IK}L_I^m`$ and the complement $`U(K)=^m|𝒜(K)|`$, that is
(13)
$$U(K)=^m\underset{IK}{}L_I.$$
It can be easily seen that the complement of any coordinate subspace arrangement in $`^m`$ (i.e., the complement of any set of planes (12)) is $`U(K)`$ for some $`K`$. Note that $`U(K)`$ is invariant with respect to the standard $`T^m`$-action on $`^m`$. The following lemma establishes the connection between moment-angle complexes and complements of coordinate subspace arrangements.
###### Lemma 4.2.
The complement $`U(K)`$ is $`T^m`$-equivariantly homotopy equivalent to the moment-angle complex $`𝒵_K`$.
###### Proof.
See \[BP4, Lemma 2.13\]
The next theorem describes the cohomology algebra of $`𝒵_K`$.
###### Theorem 4.3.
The following isomorphisms of graded algebras holds:
(14)
$$H^{}\left(𝒵_K\right)Tor_{𝐤[v_1,\mathrm{},v_m]}(𝐤(K),𝐤)H[𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m],d],$$
where $`𝐤(K)=𝐤[v_1,\mathrm{},v_m]/I`$ is the face ring, and the differential $`d`$ is defined by $`d(v_i)=0`$, $`d(u_i)=v_i`$, $`i=1,\mathrm{},m`$.
###### Proof.
See \[BP4, theorems 3.2 and 3.3\]
Due to Lemma 4.2, isomorphism (14) also holds for the cohomology algebra of the complement of a coordinate subspace arrangement in $`^m`$. The first isomorphism of (14) is proved by applying the Eilenberg–Moore spectral sequence to some $`T^m`$-bundles. The second isomorphism follows from the Koszul complex for the $`𝐤[v_1,\mathrm{},v_m]`$-module $`𝐤(K)`$.
###### Remark.
As we mentioned in the introduction, the Betti numbers of the complement of a real coordinate subspace arrangement were calculated in terms of resolution of the Stanley–Reisner ring in \[GPW\]. The latter paper also contains the formulation of the multiplicative isomorphism (14) for complex coordinate subspace arrangements (see \[GPW, Thm. 3.6\]) with a reference to yet unpublished paper by Babson and Chan. We note also, that, as it was observed in \[GPW\], there is no isomorphism between the cohomology algebra of a real coordinate subspace arrangement and the corresponding Stanley–Reisner $`Tor`$-algebra.
The $`Tor`$-algebra from (14) is naturally a bigraded algebra with $`bideg(v_i)=(0,2)`$, $`bideg(u_i)=(1,2)`$, and differential $`d`$ adding $`(1,0)`$ to bidegree. Since the differential does not change the second grading, the whole algebra is decomposed into the sum of differential algebras consisting of elements with fixed second degree. Below we deduce some important consequences of this bigraded structure.
###### Remark.
Note that according to our agreement the first degree in the $`Tor`$-algebra is non-positive. This corresponds to numerating the terms of Koszul $`𝐤[v_1,\mathrm{},v_m]`$-free resolution of $`𝐤`$ by non-positive integers. In such notations the Koszul complex $`[𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m],d]`$ becomes a cochain complex, and $`Tor_{𝐤[v_1,\mathrm{},v_m]}(𝐤(K),𝐤)`$ is its cohomology, not the homology as usually regarded. This is the standard trick used for applying the Eilenberg–Moore spectral sequence, see \[Sm\]. It also explains why we write $`Tor_{𝐤[v_1,\mathrm{},v_m]}^,(𝐤(K),𝐤)`$ instead of usual $`Tor_,^{𝐤[v_1,\mathrm{},v_m]}(𝐤(K),𝐤)`$.
Following \[BP2\], define the subcomplex $`𝒞^{}(K)`$ of the cochain complex $`[𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m],d]`$ (see (14)) as follows. As a $`𝐤`$-module, $`𝒞^{}(K)=_{q=0}^m𝒞^q(K)`$, where $`𝒞^q(K)`$ is generated by monomials $`u_{j_1}\mathrm{}u_{j_q}`$ and $`v_{i_1}\mathrm{}v_{i_p}u_{j_1}\mathrm{}u_{j_q}`$ such that $`\{i_1,\mathrm{},i_p\}`$ is a simplex of $`K`$ and $`\{i_1,\mathrm{},i_p\}\{j_1,\mathrm{},j_q\}=\mathrm{}`$. Since $`d(u_i)=v_i`$, we have $`d\left(𝒞^q(K)\right)𝒞^{q+1}(K)`$ and, therefore, $`𝒞^{}(K)`$ is a cochain subcomplex. Moreover, $`𝒞^{}(K)`$ inherits the bigraded module structure from $`𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$, with differential $`d`$ adding $`(1,0)`$ to bidegree. Hence, we have the additive inclusion (i.e., the monomorphism of bigraded modules) $`i:𝒞^{}(K)𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$. Finally, $`𝒞^{}(K)`$ can be viewed as an algebra in obvious way, however this time this is not a subalgebra of $`𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$ (since, for instance, $`v_1^2=0`$ in $`𝒞^{}(K)`$ but not in $`𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$). However, we have the multiplicative projection (i.e., the epimorphism of bigraded algebras) $`j:𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m]𝒞^{}(K)`$. The additive inclusion $`i`$ and the multiplicative projection $`j`$ obviously satisfy $`ji=id`$.
###### Lemma 4.4.
Cochain complexes $`[𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m],d]`$ and $`[𝒞^{}(K),d]`$ have same cohomology. Hence, the following isomorphism of bigraded $`𝐤`$-modules holds:
$$H[𝒞^{}(K),d]Tor_{k[v_1,\mathrm{},v_m]}^{}(𝐤(K),𝐤).$$
###### Proof.
See \[BP2, Lemma 5.3\]. ∎
In the sequel we denote (square-free) monomials $`v_{i_1}\mathrm{}v_{i_p}u_{j_1}\mathrm{}u_{j_q}𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$ by $`v_Iu_J`$, where $`I=\{i_1,\mathrm{},i_p\}`$, $`J=\{j_1,\mathrm{},j_q\}`$ are multiindices.
Now we recall our cellular decomposition of $`𝒵_K`$, whose cells are $`D_IT_J`$, where $`I,J[m]`$, $`I`$ is a simplex of $`K`$, and $`IJ=\mathrm{}`$. Let $`𝒞_{}(𝒵_K)`$ and $`𝒞^{}(𝒵_K)`$ denote the corresponding cellular chain and cochain complex respectively. Both complexes (or differential algebras) $`𝒞^{}(𝒵_K)`$ and $`𝒞^{}(K)`$ have same cohomology $`H^{}(𝒵_K)`$. Cochain complex $`𝒞^{}(𝒵_K)`$ has basis consisting of elements $`(D_IT_J)^{}`$ dual to $`D_IT_J𝒞_{}(𝒵_K)`$ (the latter is viewed as a cellular chain). Note that the cochain algebra $`𝒞^{}(𝒵_K)`$ is multiplicatively generated by the elements $`T_i^{}`$, $`D_i^{}`$, $`i,j=1,\mathrm{},m`$, (of dimension 1 and 2 respectively), while $`𝒞^{}(K)`$ is multiplicatively generated by $`u_i`$, $`v_i`$, $`i,j=1,\mathrm{},m`$. The following theorem shows that these two algebras are the same.
###### Theorem 4.5.
The correspondence $`v_Iu_J(D_IT_J)^{}`$ establishes a canonical isomorphism of differential graded algebras $`𝒞^{}(K)`$ and $`𝒞^{}(𝒵_K)`$.
###### Proof.
It follows directly from the definitions of $`𝒞^{}(K)`$ and $`𝒞^{}(𝒵_K)`$ that the constructed map is an isomorphism of graded algebras. So, it remains to prove that it commutes with differentials. Let $`d`$, $`d_c`$ and $`_c`$ denote the differentials in $`𝒞^{}(K)`$, $`𝒞^{}(𝒵_K)`$ and $`𝒞_{}(𝒵_K)`$ respectively. Since $`d(v_i)=0`$, $`d(u_i)=v_i`$, we need to show that $`d_c(D_i^{})=0`$, $`d_c(T_i^{})=D_i^{}`$. We have $`_c(D_i)=T_i`$, $`_c(T_i)=0`$. Since any 2-cell of $`𝒵_K`$ is either $`D_j`$ or $`T_{jk}`$, $`kj`$, it follows that
$$(d_cT_i^{},D_j)=(T_i^{},_cD_j)=(T_i^{},T_j)=\delta _{ij},(d_cT_i^{},T_{jk})=(T_i^{},_cT_{jk})=0,$$
where $`\delta _{ij}=1`$ if $`i=j`$ and $`\delta _{ij}=0`$ otherwise. Hence, $`d(T_i^{})=D_i^{}`$. Further, since any 3-cell of $`𝒵_K`$ is either $`D_jT_k`$ or $`T_{j_1j_2j_3}`$, it follows that
$`(d_cD_i^{},D_jT_k)=(D_i^{},_c(D_jT_k))=(D_i^{},T_{jk})=0,`$
$`(d_cD_i^{},T_{j_1j_2j_3})=(D_i^{},_cT_{j_1j_2j_3})=0.`$
Hence, $`d_c(D_i^{})=0`$. ∎
In the sequel we would not distinguish cochain complexes $`𝒞^{}(K)`$ and $`𝒞^{}(𝒵_K)`$ and their elements $`u_i`$ and $`T_i^{}`$, $`v_i`$ and $`D_i^{}`$. The above theorem provides two methods for calculating the (co)homology of $`𝒵_K`$: either by means of the differential (bi)graded algebra $`[𝒞^{}(K),d]`$, where $`𝒞^{}(K)𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$ (as modules), or using the cellular chain complex $`[𝒞_{}(𝒵_K),_c]`$.
Now we recall that the algebra $`[𝒞^{}(K),d]`$ is bigraded. Theorem 4.5 shows that the cellular chain complex $`[𝒞_{}(𝒵_K),_c]`$ can be also made bigraded by setting
(15)
$$bideg(D_i)=(0,2),bideg(T_i)=(1,2),bideg(1_i)=(0,0).$$
The differential $`_c`$ adds $`(1,0)`$ to bidegree, and the cellular homology of $`𝒵_K`$ also acquires a bigraded structure. Let us assume now that the ground field $`𝐤`$ is of zero characteristic (e.g., $`𝐤=`$ is the field of rational numbers). Define the bigraded Betti numbers
(16)
$$b_{q,2p}(𝒵_K)=dimH_{q,2p}[𝒞_{}(𝒵_K),_c],q,p=0,\mathrm{},m.$$
Theorem 4.5 and Lemma 4.4 show that
$$b_{q,2p}(𝒵_K)=dimTor_{𝐤[v_1,\mathrm{},v_m]}^{q,2p}(𝐤(K),𝐤)$$
(i.e., $`b_{q,2p}(𝒵_K)`$ is the dimension of $`(q,2p)`$-th bigraded component of the cohomology algebra $`H[𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m],d]`$). For the ordinary Betti numbers $`b_k(𝒵_K)`$ holds
$$b_k(𝒵_K)=\underset{q+2p=k}{}b_{q,2p}(𝒵_K),k=0,\mathrm{},m+n.$$
Below we describe some basic properties of bigraded Betti numbers (16).
###### Lemma 4.6.
Let $`K^{n1}`$ be an $`(n1)`$-dimensional simplicial complex with $`m=f_0`$ vertices and $`f_1`$ edges, and let $`𝒵_K`$ be the corresponding moment-angle complex, $`dim𝒵_K=m+n`$. Then
* $`b_{0,0}(𝒵_K)=b_0(𝒵_K)=1`$, $`b_{0,2p}(𝒵_K)=0`$ if $`p>0`$;
* $`b_{q,2p}=0`$ if $`p>m`$ or $`q>p`$;
* $`b_1(𝒵_K)=b_2(𝒵_K)=0`$;
* $`b_3(𝒵_K)=b_{1,4}(𝒵_K)=\left(\genfrac{}{}{0pt}{}{f_0}{2}\right)f_1`$;
* $`b_{q,2p}(𝒵_K)=0`$ if $`qp>0`$ or $`pq>n`$;
* $`b_{m+n}(𝒵_K)=b_{(mn),2m}(𝒵_K)`$.
###### Proof.
We use the cochain complex $`𝒞^{}(K)𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$ for calculations. The module $`𝒞^{}(K)`$ has basis consisting of monomials $`v_Iu_J`$ with $`IK`$ and $`IJ=\mathrm{}`$. Since $`bidegv_i=(0,2)`$, $`bidegu_j=(1,2)`$, the bigraded component $`𝒞^{q,2p}(K)`$ is generated by monomials $`v_Iu_J`$ with $`\mathrm{\#}I=pq`$ and $`\mathrm{\#}J=q`$. In particular, $`𝒞^{q,2p}(K)=0`$ if $`p>m`$ or $`q>p`$, whence the assertion (b) follows. To prove (a) we mention that $`𝒞^{0,0}(K)`$ is generated by 1, while any $`v_I𝒞^{0,2p}(K)`$, $`p>0`$, is a coboundary, hence, $`H^{q,2p}(K)=0`$, $`p>0`$.
Now we are going to prove the assertion (e). Since any $`v_Iu_J𝒞^{q,2p}(K)`$ has $`IK`$, and any simplex of $`K`$ is at most $`(n1)`$-dimensional, it follows that $`𝒞^{q,2p}(K)=0`$ for $`pq>n`$. It follows from (b) that $`b_{q,2p}(𝒵_K)=0`$ for $`q>p`$, so it remains to prove that $`b_{q,2q}(𝒵_K)=0`$ for $`q>0`$. The module $`𝒞^{q,2q}(K)`$ is generated by monomials $`u_J`$, $`\mathrm{\#}J=q`$. Since $`d(u_i)=v_i`$, it follows easily that there no cocycles in $`𝒞^{q,2q}(K)`$. Hence, $`H^{q,2q}(𝒵_K)=0`$.
To prove (c) we mention that $`H^1(𝒵_K)=H^{1,2}(K)`$ and $`H^2(𝒵_K)=H^{2,4}(𝒵_K)`$ (this follows from (a) and (b)). Hence, the assertion (c) follows from (e).
As for (d), it follows from (a), (b) and (e) that $`H^3(𝒵_K)=H^{1,4}(𝒵_K)`$. The module $`𝒞^{1,4}(K)`$ is generated by monomials $`v_iu_j`$, $`ij`$. We have $`d(v_iu_j)=v_iv_j`$ and $`d(u_iu_j)=v_iu_jv_ju_i`$. It follows that $`v_iu_j`$ is a cocycle if and only if $`\{i,j\}`$ is not a 1-simplex in $`K`$; in this case two cocycles $`v_iu_j`$ and $`v_ju_i`$ are cohomological. The assertion (d) now follows easily.
The remaining assertion (f) follows from the fact that the monomial $`u_Iv_J𝒞^{}(K)`$ of maximal total degree $`m+n`$ necessarily has $`\mathrm{\#}I+\mathrm{\#}J=m`$, $`\mathrm{\#}J=n`$, $`\mathrm{\#}I=mn`$. ∎
The above lemma shows that non-zero bigraded Betti numbers $`b_{r,2p}(𝒵_K)`$, $`r0`$ appear only in the “strip” bounded by the lines $`r=(m1)`$, $`r=1`$, $`p+r=1`$ and $`p+r=n`$ in the second quadrant (see Figure 4 (a)).
Let us consider now the bigraded cellular chain complex $`[𝒞_,(𝒵_K),_c]`$. The homogeneous component $`𝒞_{q,2p}(𝒵_K)`$ consists of cellular chains $`D_IT_J`$ with $`IK`$, $`\mathrm{\#}I=pq`$, $`\mathrm{\#}J=q`$. It follows that
(17)
$$dim𝒞_{q,2p}(𝒵_K)=f_{pq1}\left(\genfrac{}{}{0pt}{}{mp+q}{q}\right)$$
(with usual agreement $`\left(\genfrac{}{}{0pt}{}{i}{j}\right)=0`$ if $`i<j`$ or $`j<0`$), where $`f_i`$ is the number of $`i`$-simplices of $`K^{n1}`$ and $`f_1=1`$. Since the differential $`_c`$ does not change the second degree, i.e.
$$_c:𝒞_{q,2p}(𝒵_K)𝒞_{q1,2p}(𝒵_K),$$
the chain complex $`𝒞_,(𝒵_K)`$ splits into the sum of chain complexes as follows:
$$[𝒞_,(𝒵_K),_c]=\underset{p=0}{\overset{m}{}}[𝒞_{,2p}(𝒵_K),_c].$$
The similar decomposition holds also for the cellular cochain complex $`[𝒞^,(𝒵_K),d_c][𝒞^,(K),d]`$. Let $`\chi _p(𝒵_K)`$ denote the Euler characteristic of complex $`[𝒞_{,2p}(𝒵_K),_c]`$, i.e.
(18)
$$\chi _p(𝒵_K)=\underset{q=0}{\overset{m}{}}(1)^qdim𝒞_{q,2p}(𝒵_K)=\underset{q=0}{\overset{m}{}}(1)^qb_{q,2p}(𝒵_K).$$
Define the generating polynomial $`\chi (𝒵_K,t)`$ as
$$\chi (𝒵_K,t)=\underset{p=0}{\overset{m}{}}\chi _p(𝒵_K)t^{2p}.$$
The following theorem calculates this polynomial in terms of the numbers of faces of $`K`$. It was firstly proved in \[BP2\] in the particular case of polytopal $`K`$.
###### Theorem 4.7.
For any $`(n1)`$-dimensional simplicial complex $`K`$ with $`m`$ vertices holds
(19)
$$\chi (𝒵_K,t)=(1t^2)^{mn}(h_0+h_1t^2+\mathrm{}+h_nt^{2n}),$$
where $`(h_0,h_1,\mathrm{},h_n)`$ is the $`h`$-vector of $`K`$ (see (2)).
###### Proof.
It follows from (18) and (17) that
(20)
$$\chi _p(𝒵_K)=\underset{j=0}{\overset{m}{}}(1)^{pj}f_{j1}\left(\genfrac{}{}{0pt}{}{mj}{pj}\right),$$
Then
(21)
$$\begin{array}{c}\chi (𝒵_K,t)=\underset{p=0}{\overset{m}{}}\chi _p(K)t^{2p}=\underset{p=0}{\overset{m}{}}\underset{j=0}{\overset{m}{}}t^{2j}t^{2(pj)}(1)^{pj}f_{j1}\left(\genfrac{}{}{0pt}{}{mj}{pj}\right)\hfill \\ \hfill =\underset{j=0}{\overset{m}{}}f_{j1}t^{2j}(1t^2)^{mj}=(1t^2)^m\underset{j=0}{\overset{n}{}}f_{j1}(t^21)^j.\end{array}$$
Denote $`h(t)=h_0+h_1t+\mathrm{}+h_nt^n`$. Then it follows from (2) that
$$t^nh(t^1)=(t1)^n\underset{i=0}{\overset{n}{}}f_{i1}(t1)^i.$$
Substituting above $`t^2`$ for $`t`$, we finally obtain from (21)
$$\frac{\chi (𝒵_K,t)}{(1t^2)^m}=\frac{t^{2n}h(t^2)}{(t^21)^n}=\frac{h(t^2)}{(1t^2)^n},$$
which is equivalent to (19). ∎
The above theorem allows to express the numbers of faces of a simplicial complex in terms of bigraded Betti numbers of the corresponding moment-angle complex $`𝒵_K`$. The first important corollary of this is as follows.
###### Corollary 4.8.
For any simplicial complex $`K`$ the Euler number of the corresponding moment-angle complex $`𝒵_K`$ is zero.
###### Proof.
We have
$$\chi (𝒵_K)=\underset{p,q=0}{\overset{m}{}}(1)^{q+2p}b_{q,2p}(𝒵_K)=\underset{p=0}{\overset{m}{}}\chi _p(𝒵_K)=\chi (𝒵_K,1)$$
Now the statement follows from (19). ∎
###### Remark.
Another way to prove the above corollary is to mention that the diagonal subgroup $`S^1T^m`$ always acts freely on the moment-angle complex $`𝒵_K`$ (see \[BP2\]). Hence, there exists a principal $`S^1`$-bundle $`𝒵_K𝒵_K/S^1`$, which implies $`\chi (𝒵_K)=0`$.
###### Corollary 4.9.
The Euler number of the complement of a complex coordinate subspace arrangement is zero.
###### Proof.
This follows from the previous corollary and Lemma 4.2. ∎
By definition (see Proposition 2.2), the cubical complex $`cc(K)`$ always contains the vertex $`(1,\mathrm{},1)I^m`$. Hence, the torus $`T^m=\rho ^1(1,\mathrm{},1)`$ is contained in $`𝒵_K`$. Here $`\rho :(D^2)^mI^m`$ is the orbit map for the $`T^m`$-action (see (11)).
###### Lemma 4.10.
The inclusion $`T^m=\rho ^1(1,\mathrm{},1)𝒵_K`$ is a cellular map homotopical to the map to a point, i.e. the torus $`T^m=\rho ^1(1,\mathrm{},1)`$ is a contractible cellular subcomplex of $`𝒵_K`$.
###### Proof.
To prove that $`T^m=\rho ^1(1,\mathrm{},1)`$ is a cellular subcomplex of $`𝒵_K`$ we just mention that this $`T^m`$ is the closure of the $`m`$-cell $`D_{\mathrm{}}T_{[m]}𝒵_K`$. So, it remains to prove that $`T^m`$ is contractible within $`𝒵_K`$. To do this we show that the embedding $`T^m(D^2)^m`$ is homotopical to the map to the point $`(1,\mathrm{},1)T^m(D^2)^m`$. On the first step we note that $`𝒵_K`$ contains the cell $`D_1T_{2,\mathrm{},m}`$, whose closure contains $`T^m`$ and is homeomorphic to $`D^2\times T^{m1}`$. Hence, our $`T^m`$ can be contracted to $`1\times T^{m1}`$ within $`𝒵_K`$. On the second step we note that $`𝒵_K`$ contains the cell $`D_2T_{3,\mathrm{},m}`$, whose closure contains $`1\times T^{m1}`$ and is homeomorphic to $`D^2\times T^{m2}`$. Hence, $`1\times T^{m1}`$ can be contracted to $`1\times 1\times T^{m2}`$ within $`𝒵_K`$, and so on. On the $`k`$th step we note that $`𝒵_K`$ contains the cell $`D_kT_{k+1,\mathrm{},m}`$, whose closure contains $`1\times \mathrm{}\times 1\times T^{mk+1}`$ and is homeomorphic to $`D^2\times T^{mk}`$. Hence, $`1\times \mathrm{}\times 1\times T^{mk+1}`$ can be contracted to $`1\times \mathrm{}\times 1\times T^{mk}`$ within $`𝒵_K`$. We end up at the point $`1\times \mathrm{}\times 1`$ to which the whole torus $`T^m`$ can be contracted. ∎
###### Corollary 4.11.
For any simplicial complex $`K`$ the moment-angle complex $`𝒵_K`$ is simply connected.
###### Proof.
Indeed, the 1-skeleton of our cellular decomposition of $`𝒵_K`$ is contained in the torus $`T^m=\rho ^1(1,\mathrm{},1)`$. ∎
The cohomology of cellular pair $`(𝒵_K,T^m)`$ also can be calculated by means of the cochain complex $`𝒞^{}(K)`$. The cellular cochain subcomplex $`𝒞^{}(T^m)𝒞^{}(K)`$ consists of monomials $`u_I`$ (i.e. monomials that do not contain $`v_i`$’s). This, of course, is just the exterior algebra $`\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$. Hence,
(22)
$$𝒞^{}(𝒵_K,T^m)=𝒞^{}(𝒵_K)/\mathrm{\Lambda }[u_1,\mathrm{},u_m]$$
(as complexes, not as algebras). We can also introduce relative bigraded Betti numbers
(23)
$$b_{q,2p}(𝒵_K,T^m)=dimH^{q,2p}[𝒞^{}(𝒵_K,T^m),d],q,p=0,\mathrm{},m,$$
define the $`p`$th relative Euler characteristic $`\chi _p(𝒵_K,T^m)`$ as the Euler number of complex $`𝒞^{,2p}(𝒵_K,T^m)`$:
(24)
$$\chi _p(𝒵_K,T^m)=\underset{q=0}{\overset{m}{}}(1)^qdim𝒞^{q,2p}(𝒵_K,T^m)=\underset{q=0}{\overset{m}{}}(1)^qb_{q,2p}(𝒵_K,T^m),$$
and define the generating polynomial $`\chi (𝒵_K,T^m,t)`$ as
$$\chi (𝒵_K,T^m,t)=\underset{p=0}{\overset{m}{}}\chi _p(𝒵_K,T^m)t^{2p}.$$
We will use the following theorem in the next section.
###### Theorem 4.12.
For any $`(n1)`$-dimensional simplicial complex $`K`$ with $`m`$ vertices holds
(25)
$$\chi (𝒵_K,T^m,t)=(1t^2)^{mn}(h_0+h_1t^2+\mathrm{}+h_nt^{2n})(1t^2)^m.$$
###### Proof.
Since $`𝒞^{}(T^m)=\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$ and $`bidegu_i=(1,2)`$, we have
$$dim𝒞^q(T^m)=dim𝒞^{q,2q}(T^m)=\left(\genfrac{}{}{0pt}{}{m}{q}\right).$$
It follows from (22), (18) and (24) that
$$\chi _p(𝒵_K,T^m)=\chi _p(𝒵_K)(1)^pdim𝒞^{p,2p}(T^m).$$
Hence,
$`\chi (𝒵_K,T^m,t)`$ $`=\chi (𝒵_K,t){\displaystyle \underset{p=0}{\overset{m}{}}}(1)^p\left(\genfrac{}{}{0pt}{}{m}{p}\right)t^{2p}`$
$`=(1t^2)^{mn}(h_0+h_1t^2+\mathrm{}+h_nt^{2n})(1t^2)^m,`$
by (19). ∎
At the end of this section we review the most important additional properties of $`𝒵_K`$ in the case when $`|K|S^{n1}`$, i.e. $`K`$ is a simplicial sphere.
###### Lemma 4.13.
If $`K`$ is a simplicial sphere, i.e. $`|K|=S^{n1}`$, then $`𝒵_K`$ is an $`(m+n)`$-dimensional (closed) manifold. $`\mathrm{}`$
In \[DJ, p. 434\] the authors considered the manifold $`𝒵`$ defined for any simple $`n`$-polytope $`P^{}`$ with $`m`$ facets as $`𝒵=(T^m\times P^{})/`$, where $``$ is a certain equivalence relation. We showed in \[BP2\] that if $`K`$ is a polytopal sphere, i.e. $`K=P^n`$ for some simplicial polytope $`P^n`$, then our moment-angle complex $`𝒵_K`$ coincides with the manifold $`𝒵`$ defined by simple polytope $`P^{}`$ dual to $`P`$. For the case of general simplicial sphere $`K`$, see \[BP3\]\[BP4\].
###### Theorem 4.14.
Let $`K`$ be an $`(n1)`$-dimensional simplicial sphere, and let $`𝒵_K`$ be the corresponding moment-angle manifold. The fundamental cohomological class of $`𝒵_K`$ is represented by any monomial $`\pm v_Iu_J𝒞(K)`$ of bidegree $`((mn),2m)`$ such that $`I`$ is an $`(n1)`$-simplex of $`K`$ and $`IJ=\mathrm{}`$. The choice of sign depends on the orientation of $`𝒵_K`$.
###### Proof.
We have $`dim𝒵_K=m+n`$. It follows from Lemma 4.6 (f) that $`H^{m+n}(𝒵_K)=H^{(mn),2m}(𝒵_K)`$. By definition, the module $`𝒞^{(mn),2m}(𝒵_K)`$ is spanned by monomials $`v_Iu_J`$ such that $`IK^{n1}`$, $`\mathrm{\#}I=n`$, $`J=[m]I`$, and all these monomials are cocycles. Suppose that $`I,I^{}`$ are two $`(n1)`$-simplices of $`K^{n1}`$ sharing a common $`(n2)`$-face. Then the corresponding cocycles $`v_Iu_J`$, $`v_I^{}u_J^{}`$, where $`J=[m]I`$, $`J^{}=[m]I^{}`$, are cohomological (up to sign). Indeed, let $`v_Iu_J=v_{i_1}\mathrm{}v_{i_n}u_{j_1}\mathrm{}u_{j_{mn}}`$, $`v_I^{}u_J^{}=v_{i_1}\mathrm{}v_{i_{n1}}v_{j_1}u_{i_n}u_{j_2}\mathrm{}u_{j_{mn}}`$. Since any $`(n2)`$-face of $`K`$ is contained in exactly two $`(n1)`$-faces, the identity
$$\begin{array}{c}d(v_{i_1}\mathrm{}v_{i_{n1}}u_{i_n}u_{j_1}u_{j_2}\mathrm{}u_{j_{mn}})\hfill \\ \hfill =v_{i_1}\mathrm{}v_{i_n}u_{j_1}\mathrm{}u_{j_{mn}}v_{i_1}\mathrm{}v_{i_{n1}}v_{j_1}u_{i_n}u_{j_2}\mathrm{}u_{j_{mn}}\end{array}$$
holds in $`𝒞(K)𝐤(K)\mathrm{\Lambda }[u_1,\mathrm{},u_m]`$, hence, $`v_Iu_J`$ and $`v_I^{}u_J^{}`$ are cohomological. Since $`K^{n1}`$ is a simplicial sphere, any two $`(n1)`$-simplices of $`K^{n1}`$ can be connected by a chain of simplices such that any two successive simplices share a common $`(n2)`$ face. Thus, any two cocycles in $`𝒞^{(mn),2m}(𝒵_K)`$ are cohomological, and we can take any one as a representative for the fundamental cohomological class of $`𝒵_K`$ (after a proper choice of sign). ∎
###### Remark.
In the proof of the above theorem we have used two combinatorial properties of $`K^{n1}`$. The first one is that any $`(n2)`$-face is contained in exactly two $`(n1)`$-faces, and the second one is that any two $`(n1)`$-simplices can be connected by a chain of simplices such that any two successive simplices share a common $`(n2)`$-face. Both properties hold for any simplicial manifold. Hence, for any simplicial manifold $`K^{n1}`$ one has $`b_{m+n}(𝒵_K)=b_{(mn),2m}(𝒵_K)=1`$ and the generator of $`H^{m+n}(𝒵_K)`$ can be chosen as in Theorem 4.14.
###### Corollary 4.15.
The Poincaré duality for the moment angle manifold $`𝒵_K`$ defined by a simplicial sphere $`K^{n1}`$ regards the bigraded structure in the (co)homology, i.e.
$$H^{q,2p}(𝒵_K)H_{(mn)+q,2(mp)}(𝒵_K).$$
In particular,
(26)
$$b_{q,2p}(𝒵_K)=b_{(mn)+q,2(mp)}(𝒵_K).\mathrm{}$$
###### Corollary 4.16.
Let $`K^{n1}`$ be an $`(n1)`$-dimensional simplicial sphere, and let $`𝒵_K`$ be the corresponding moment-angle complex, $`dim𝒵_K=m+n`$. Then
* $`b_{q,2p}(𝒵_K)=0`$ if $`qmn`$, with only exception $`b_{(mn),2m}=1`$;
* $`b_{q,2p}(𝒵_K)=0`$ if $`pqn`$, with only exception $`b_{(mn),2m}=1`$. $`\mathrm{}`$
It follows that if $`K^{n1}`$ is a simplicial sphere, then non-zero bigraded Betti numbers $`b_{r,2p}(𝒵_K)`$, $`r0`$, $`rmn`$, appear only in the “strip” bounded by the lines $`r=(mn1)`$, $`r=1`$, $`p+r=1`$ and $`p+r=n1`$ in the second quadrant (see Figure 4 (b)). Compare this with Figure 4 (a) corresponding to the case of general $`K`$.
It follows from (18) and (26) that for any simplicial sphere $`K`$ holds
$$\chi _p(𝒵_K)=(1)^{mn}\chi _{mp}(𝒵_K).$$
From this and (19) we get
$$\begin{array}{c}\frac{h_0+h_1t^2+\mathrm{}+h_nt^{2n}}{(1t^2)^n}=(1)^{mn}\frac{\chi _m+\chi _{m1}t^2+\mathrm{}+\chi _0t^{2m}}{(1t^2)^m}\hfill \\ \hfill =(1)^n\frac{\chi _0+\chi _1t^2+\mathrm{}+\chi _mt^{2m}}{(1t^2)^m}=(1)^n\frac{h_0+h_1t^2+\mathrm{}+h_nt^{2n}}{(1t^2)^n}\\ \hfill =\frac{h_0t^{2n}+h_1t^{2(n1)}+\mathrm{}+h_n}{(1t^2)^n}.\end{array}$$
Hence, $`h_i=h_{ni}`$. Thus, we have deduced the Dehn–Sommerville equations as a corollary of the bigraded Poincaré duality (26).
The identity (19) also allows to interpret different inequalities for the numbers of faces of simplicial spheres (resp. simplicial manifolds) in terms of topological invariants (bigraded Betti numbers) of the corresponding moment-angle manifolds (resp. complexes) $`𝒵_K`$.
###### Example 4.17.
It follows from Lemma 4.6 that for any $`K`$ holds
$`\chi _0(K)=1,`$
$`\chi _1(K)=0,`$
$`\chi _2(K)=b_{1,4}(𝒵_K)=b_3(𝒵_K),`$
$`\chi _3(K)=b_{2,6}(𝒵_K)b_{1,6}(𝒵_K)`$
(note that $`b_4(𝒵_K)=b_{2,6}(𝒵_K)`$, while $`b_5(𝒵_K)=b_{1,6}(𝒵_K)+b_{3,8}(𝒵_K)`$). Now, identity (19) shows that
$`h_0=1,`$
$`h_1=mn,`$
$`h_2=\left(\genfrac{}{}{0pt}{}{mn+1}{2}\right)b_3(𝒵_K),`$
$`h_3=\left(\genfrac{}{}{0pt}{}{mn+2}{3}\right)(mn)b_{1,4}(𝒵_K)+b_{2,6}(𝒵_K)b_{1,6}(𝒵_K).`$
It follows that the inequality $`h_1h_2`$, $`n4`$, from the Generalized Lower Bound hypothesis (7) for simplicial spheres is equivalent to the following inequality:
(27)
$$b_3(𝒵_K)\left(\genfrac{}{}{0pt}{}{mn}{2}\right).$$
The next inequality $`h_2h_3`$, $`n6`$, from (7) is equivalent to the following inequality for the bigraded Betti numbers of $`𝒵_K`$:
(28)
$$\left(\genfrac{}{}{0pt}{}{mn+1}{3}\right)(mn1)b_{1,4}(𝒵_K)+b_{2,6}(𝒵_K)b_{1,6}(𝒵_K)0.$$
Thus, we see that the combinatorial Generalized Lower Bound inequalities are interpreted as “topological” inequalities for the (bigraded) Betti numbers of a certain manifold. So, one can try to use topological methods (such as the equivariant topology or Morse theory) to prove inequalities like (27) or (28). Such topological approach to the hypotheses like $`g`$-theorem or Generalized Lower Bound has an advantage of being independent on whether the simplicial sphere $`K`$ is polytopal or not. Indeed, all known proofs of the necessity of $`g`$-theorem for simplicial polytopes (including the original one by Stanley \[St2\], McMullen’s proof \[McM\], and the recent proof by Timorin \[Ti\]) follow the same scheme. Namely, the numbers $`h_i`$, $`i=1,\mathrm{},n`$, are interpreted as the dimensions of graded components $`A^i`$ of a certain graded algebra $`A`$ satisfying the Hard Lefschetz Theorem. The latter means that there is an element $`\omega A^1`$ such that the multiplication by $`\omega `$ defines a monomorphism $`A^iA^{i+1}`$ for $`i<\left[\frac{n}{2}\right]`$. This implies $`h_ih_{i+1}`$ for $`i<\left[\frac{n}{2}\right]`$. However, such element $`\omega `$ is lacking for non-polytopal $`K`$, which means that one should develop a new technique for proving the $`g`$-theorem (or, better to say, $`g`$-conjecture) for simplicial spheres. Certainly, it may happen that the $`g`$-theorem fails to be true for simplicial spheres, however, many recent efforts of computer-aided seek for counter examples were unsuccessful (see, e.g., \[BjLu\]).
It can be easily seen that our moment-angle complex $`𝒵_K`$ is a manifold provided that the cone $`cone(K)`$ is non-singular. This is equivalent to the condition that the suspension $`\mathrm{\Sigma }|K|`$ is a (topological) manifold, which implies (due to the suspension isomorphism and Poincaré duality in the homology) that $`|K|`$ is a homology sphere. An important class of simplicial homology spheres is known in combinatorics as Gorenstein\* complexes (see, e.g., \[St3\] for definition). As it was pointed out by Stanley in \[St3\], the Gorenstein\* complexes are the most general objects appropriate for generalizing the $`g`$-theorem (they include polytopal spheres, PL-spheres and simplicial spheres as particular cases). In our terms, the Gorenstein\* complexes $`K`$ (see \[St1, p.75\]) can be characterized by the condition that the $`Tor`$-algebra $`Tor_{𝐤[v_1,\mathrm{},v_m]}(𝐤(K),𝐤)`$ satisfies the bigraded duality (26), i.e. $`𝒵_K`$ is a Poincaré duality space (not necessarily a manifold). In particular, the Dehn–Sommerville relations $`h_i=h_{ni}`$ continue to hold for Gorenstein\* complexes.
## 5. Homology of $`𝒲_K`$ and generalized Dehn–Sommerville equations
Here we assume that $`K^{n1}`$ is a triangulation of a manifold, i.e., a simplicial manifold. In this case the moment-angle complex $`𝒵_K`$ is not a manifold, however, its singularities can be easily treated. Indeed, $`|cc(K)|`$ is homeomorphic to $`|cone(K)|`$, and the vertex of the cone is the point $`p=(1,\mathrm{},1)|cc(K)|I^m`$. Let $`U_\epsilon (p)|cc(K)|`$ be a small neighbourhood of $`p`$ in $`|cc(K)|`$. Then the closure of $`U_\epsilon (p)`$ is also homeomorphic to $`|cone(K)|`$. It follows from the definition of $`𝒵_K`$ (see (11)) that $`U_\epsilon (T^m)=\rho ^1\left(U_\epsilon (p)\right)𝒵_K`$ is a small invariant neighbourhood of the torus $`T^m=\rho ^1(p)`$ in $`𝒵_K`$. Here $`\rho :(D^2)^mI^m`$ is the orbit map. Then for small $`\epsilon `$ the closure of the neighbourhood $`U_\epsilon (T^m)`$ is homeomorphic to $`|cone(K)|\times T^m`$. Taking $`U_\epsilon (T^m)`$ away from $`𝒵_K`$ we obtain a manifold with boundary, which we denote $`W_K`$. Hence, we have
$$W_K=\overline{𝒵_K|cone(K)|\times T^m},W_K=|K|\times T^m.$$
Note that since the neighbourhood $`U_\epsilon (T^m)`$ is $`T^m`$-invariant, the torus $`T^m`$ acts on $`W_K`$.
###### Theorem 5.1.
The manifold with boundary $`W_K`$ is equivariantly homotopy equivalent to the moment-angle complex $`𝒲_K`$ (see (11)). There is a canonical relative isomorphism of pairs $`(W_K,W_K)(𝒵_K,T^m)`$.
###### Proof.
To prove the first assertion we construct homotopy equivalence $`|cc(K)|U_\epsilon (p)|cub(K)|`$ (see Proposition 2.1) as it is shown on Figure 5. This homotopy equivalence is covered by a $`T^m`$-invariant homotopy equivalence $`W_K=𝒵_KU_\epsilon (T^m)𝒲_K`$, as needed. The second assertion follows easily from the definition of $`W_K`$.
According to Lemma 3.2, the moment-angle complex $`𝒲_K(D^2)^m`$ has a cellular structure with 5 different cell types $`D_i`$, $`I_i`$, $`0_i`$, $`T_i`$, $`1_i`$, $`i=1,\mathrm{},m`$, (see Figure 3). The homology of $`𝒲_K`$ (and of $`W_K`$) can be calculated by means of the corresponding cellular chain complex, which we denote $`[𝒞_{}(𝒲_K),_c]`$. In comparison with the moment-angle complex $`𝒵_K`$ studied in the previous section the complex $`𝒲_K`$ has more types of cells (remember that $`𝒵_K`$ has only 3 cell types $`D_i`$, $`T_i`$, $`1_i`$). However, the wonderful thing is that the cellular chain complex $`[𝒞_{}(𝒲_K),_c]`$ can be canonically made bigraded. Namely, the following statement holds (compare with (15)).
###### Lemma 5.2.
Put
(29)
$$bidegD_i=(0,2),bidegT_i=(1,2),bidegI_i=(1,0),$$
$$bideg0_i=bideg1_i=(0,0),i=1,\mathrm{},m.$$
This makes the cellular chain complex $`[𝒞_{}(𝒲_K),_c]`$ a bigraded differential module with differential $`_c`$ adding $`(1,0)`$ to bidegree. The original grading of $`𝒞_{}(𝒲_K)`$ by dimensions of cells corresponds to the total degree (i.e., the dimension of a cell equals the sum of its two degrees).
###### Proof.
We need only to check that the differential $`_c`$ adds $`(1,0)`$ to bidegree. This follows from (29) and
$$_cD_i=T_i,_cI_i=1_i0_i,_cT_i=_c1_i=_c0_i=0.$$
Note that, unlike bigraded structure in $`𝒞_{}(𝒵_K)`$, elements of $`𝒞_,(𝒲_K)`$ may have positive first degree (due to the positive first degree of $`I_i`$’s). However, as in the case of $`𝒵_K`$, the differential $`_c`$ does not change the second degree, which allows to split the bigraded complex $`𝒞_,(𝒲_K)`$ to the sum of complexes $`𝒞_{,2p}(𝒲_K)`$, $`p=0,\mathrm{},m`$.
In the same way as we have done this for $`𝒵_K`$ and for the pair $`(𝒵_K,T^m)`$ define the bigraded Betti numbers
(30)
$$b_{q,2p}(𝒲_K)=dimH_{q,2p}[𝒞_,(𝒲_K),_c],mqm,\mathrm{\hspace{0.33em}0}pm$$
(note that $`q`$ may be both positive and negative), the $`p`$th Euler characteristic $`\chi _p(𝒲_K)`$ as the Euler number of complex $`𝒞_{,2p}(𝒲_K)`$:
(31)
$$\chi _p(𝒲_K)=\underset{q=m}{\overset{m}{}}(1)^qdim𝒞_{q,2p}(𝒲_K)=\underset{q=m}{\overset{m}{}}(1)^qb_{q,2p}(𝒲_K),$$
and the generating polynomial $`\chi (𝒲_K,t)`$ as
$$\chi (𝒲_K,t)=\underset{p=0}{\overset{m}{}}\chi _p(𝒲_K)t^{2p}.$$
The following theorem provides the exact formula for generating polynomial $`\chi (𝒲_K,t)`$ and is analogous to theorems 4.7 and 4.12
###### Theorem 5.3.
For any $`(n1)`$-dimensional simplicial complex $`K`$ with $`m`$ vertices holds
$`\chi (𝒲_K,t)`$ $`=(1t^2)^{mn}(h_0+h_1t^2+\mathrm{}+h_nt^{2n})+\left(\chi (K)1\right)(1t^2)^m`$
$`=(1t^2)^{mn}(h_0+h_1t^2+\mathrm{}+h_nt^{2n})+(1)^{n1}h_n(1t^2)^m,`$
where $`\chi (K)=f_0f_1+\mathrm{}+(1)^{n1}f_{n1}=1+(1)^{n1}h_n`$ is the Euler number of $`K`$.
###### Proof.
Lemma 3.2 shows that every moment-angle complex $`ma(C)`$ is a cellular subcomplex of $`(D^2)^m`$, and each cell is the product of cells of 5 different types: $`D_i`$, $`I_i`$, $`0_i`$, $`T_i`$, and $`1_i`$, $`i=1,\mathrm{},m`$. These products are encoded by words $`D_II_J0_LT_P1_Q`$, where $`I,J,L,P,Q`$ are pairwise disjoint subsets of $`[m]`$ such that $`IJKPQ=[m]`$. In the case $`ma(C)=𝒲_K=ma\left(cub(K)\right)`$ the definition of $`cub(K)`$ (see Proposition 2.1) shows that the cell $`D_II_J0_LT_P1_Q`$ belongs to $`𝒲_K`$ if and only if the following two conditions are satisfied:
1. The set $`IJL`$ is a simplex of $`K^{n1}`$.
2. $`\mathrm{\#}L1`$.
Let $`c_{ijlpq}(𝒲_K)`$ denote the number of cells $`D_II_J0_LT_P1_Q𝒲_K`$ with $`i=\mathrm{\#}I`$, $`j=\mathrm{\#}J`$, $`l=\mathrm{\#}L`$, $`p=\mathrm{\#}P`$, $`q=\mathrm{\#}Q`$, $`i+j+l+p+q=m`$. It follows that
(32)
$$c_{ijlpq}(𝒲_K)=f_{i+j+l1}\left(\genfrac{}{}{0pt}{}{i+j+l}{i}\right)\left(\genfrac{}{}{0pt}{}{j+l}{l}\right)\left(\genfrac{}{}{0pt}{}{mijl}{p}\right),$$
where $`(f_0,\mathrm{},f_{n1})`$ is the $`f`$-vector of $`K`$ (we also assume $`f_1=1`$ and $`f_k=0`$ for $`k<1`$ or $`k>n1`$). By (29),
$$bideg(D_II_J0_LT_P1_Q)=(jp,2(i+p)).$$
Now we calculate $`\chi _r(𝒲_K)`$ as it is defined by (31), using (32):
$$\chi _r(𝒲_K)=\underset{i+p=r,l1}{\underset{i,j,l,p}{}}(1)^{jp}f_{i+j+l1}\left(\genfrac{}{}{0pt}{}{i+j+l}{i}\right)\left(\genfrac{}{}{0pt}{}{j+l}{l}\right)\left(\genfrac{}{}{0pt}{}{mijl}{p}\right).$$
Substituting $`s=i+j+l`$, $`i=rp`$ above, we obtain
$`\chi _r(𝒲_K)`$ $`={\displaystyle \underset{l1}{\underset{l,s,p}{}}}(1)^{srl}f_{s1}\left(\genfrac{}{}{0pt}{}{s}{rp}\right)\left(\genfrac{}{}{0pt}{}{sr+p}{l}\right)\left(\genfrac{}{}{0pt}{}{ms}{p}\right)`$
$`={\displaystyle \underset{s,p}{}}\left((1)^{sr}f_{s1}\left(\genfrac{}{}{0pt}{}{s}{rp}\right)\left(\genfrac{}{}{0pt}{}{ms}{p}\right){\displaystyle \underset{l1}{}}(1)^l\left(\genfrac{}{}{0pt}{}{sr+p}{l}\right)\right)`$
Since
$$\underset{l1}{}(1)^l\left(\genfrac{}{}{0pt}{}{sr+p}{l}\right)=\{\begin{array}{cc}\hfill 1,& s>rp,\hfill \\ \hfill 0,& srp,\hfill \end{array}$$
we obtain
$`\chi _r(𝒲_K)`$ $`={\displaystyle \underset{s>rp}{\underset{s,p}{}}}(1)^{sr}f_{s1}\left(\genfrac{}{}{0pt}{}{s}{rp}\right)\left(\genfrac{}{}{0pt}{}{ms}{p}\right)`$
$`={\displaystyle \underset{s,p}{}}(1)^{rs}f_{s1}\left(\genfrac{}{}{0pt}{}{s}{rp}\right)\left(\genfrac{}{}{0pt}{}{ms}{p}\right)+{\displaystyle \underset{s}{}}(1)^{rs}f_{s1}\left(\genfrac{}{}{0pt}{}{ms}{rs}\right).`$
The second sum in the above formula is exactly $`\chi _r(𝒵_K)`$ (see (20)). To calculate the first sum, we mention that
$$\underset{p}{}\left(\genfrac{}{}{0pt}{}{s}{rp}\right)\left(\genfrac{}{}{0pt}{}{ms}{p}\right)=\left(\genfrac{}{}{0pt}{}{m}{r}\right)$$
(this follows from calculating the coefficient of $`\alpha ^r`$ in both sides of $`(1+\alpha )^s(1+\alpha )^{ms}=(1+\alpha )^m`$). Hence,
$$\chi _r(𝒲_K)=\underset{s}{}(1)^{rs}f_{s1}\left(\genfrac{}{}{0pt}{}{m}{r}\right)+\chi _r(𝒵_K)=(1)^r\left(\genfrac{}{}{0pt}{}{m}{r}\right)\left(\chi (K)1\right)+\chi _r(𝒵_K),$$
since $`_s(1)^sf_{s1}=\chi (K)1`$ (remember that $`f_1=1`$). Finally, using (19), we calculate
$`\chi (𝒲_K,t)={\displaystyle \underset{r=0}{\overset{m}{}}}\chi _r(𝒲_K)t^{2r}={\displaystyle \underset{r=0}{\overset{m}{}}}(1)^r\left(\genfrac{}{}{0pt}{}{m}{r}\right)\left(\chi (K)1\right)t^{2r}+{\displaystyle \underset{r=0}{\overset{m}{}}}\chi _r(𝒵_K)t^{2r}`$
$`=\left(\chi (K)1\right)(1t^2)^m+(1t^2)^{mn}(h_0+h_1t^2+\mathrm{}+h_nt^{2n}).`$
Suppose now that $`K`$ is an orientable simplicial manifold. It is easy to see that then $`W_K`$ is also orientable. Hence, there are relative Poincaré duality isomorphisms:
(33)
$$H_k(W_K)H^{m+nk}(W_K,_cW_K),k=0,\mathrm{},m.$$
###### Corollary 5.4.
The following relations hold for the $`h`$-vector $`(h_0,h_1,\mathrm{},h_n)`$ of any simplicial manifold $`K^{n1}`$:
$$h_{ni}h_i=(1)^i\left(\chi (K^{n1})\chi (S^{n1})\right)\left(\genfrac{}{}{0pt}{}{n}{i}\right),i=0,1,\mathrm{},n,$$
where $`\chi (S^{n1})=1+(1)^{n1}`$ is the Euler number of an $`(n1)`$-sphere.
###### Proof.
Theorem 5.1 shows that $`H^{m+nk}(W_K,_cW_K)=H^{m+nk}(𝒵_K,T^m)`$ and $`H_k(W_K)=H_k(𝒲_K)`$. Moreover, it can be seen in the same way as in Corollary 4.15 that relative Poincaré duality isomorphisms (33) regard the bigraded structures in the (co)homology of $`𝒲_K`$ and $`(𝒵_K,T^m)`$. It follows that
$$b_{q,2p}(𝒲_K)=b_{(mn)+q,2(mp)}(𝒵_K,T^m).$$
Hence,
$$\chi _p(𝒲_K)=(1)^{mn}\chi _{mp}(𝒵_K,T^m),$$
and
(34) $`\chi (𝒲_K,t)`$ $`=(1)^{mn}{\displaystyle \underset{p}{}}\chi _{mp}(𝒵_K,T^m)t^{2p}`$
$`=(1)^{mn}t^{2m}\chi (𝒵_K,T^m,\frac{1}{t}).`$
Using (25), we calculate
$$\begin{array}{c}(1)^{mn}t^{2m}\chi (𝒵_K,T^m,\frac{1}{t})\hfill \\ \hfill =(1)^{mn}t^{2m}(1t^2)^{mn}(h_0+h_1t^2+\mathrm{}+h_nt^{2n})\\ \hfill (1)^{mn}t^{2m}(1t^2)^m\\ \hfill =(1t^2)^{mn}(h_0t^{2n}+h_1t^{2n2}+\mathrm{}+h_n)+(1)^{n1}(1t^2)^m.\end{array}$$
Substituting the formula for $`\chi (𝒲_K,t)`$ from Theorem 5.3 and the above expression into formula (34), we obtain
$$\begin{array}{c}(1t^2)^{mn}(h_0+h_1t^2+\mathrm{}+h_nt^{2n})+\left(\chi (K)1\right)(1t^2)^m\hfill \\ \hfill =(1t^2)^{mn}(h_0t^{2n}+h_1t^{2n2}+\mathrm{}+h_n)+(1)^{n1}(1t^2)^m.\end{array}$$
Calculating the coefficient of $`t^{2i}`$ in both sides after dividing the above identity by $`(1t^2)^{mn}`$, we obtain $`h_{ni}h_i=(1)^i\left(\chi (K^{n1})\chi (S^{n1})\right)\left(\genfrac{}{}{0pt}{}{n}{i}\right)`$, as needed. ∎
Corollary 5.4 generalize the Dehn–Sommerville equations (4) for simplicial spheres. If $`|K|=S^{n1}`$ or $`(n1)`$ is odd, Corollary 5.4 gives just $`h_{ni}=h_i`$.
###### Corollary 5.5.
The following relations hold for any simplicial manifold $`K^{n1}`$:
$$h_{ni}h_i=(1)^i(h_n1)\left(\genfrac{}{}{0pt}{}{n}{i}\right),i=0,1,\mathrm{},n.$$
###### Proof.
Since $`\chi (K^{n1})=1+(1)^{n1}h_n`$, $`\chi (S^{n1})=1+(1)^{n1}`$, we have
$$\chi (K^{n1})\chi (S^{n1})=(1)^{n1}(h_n1)=(h_n1)$$
(the coefficient $`(1)^{n1}`$ can be dropped since for odd $`(n1)`$ the left side is zero). ∎
###### Corollary 5.6.
For any $`(n1)`$-dimensional simplicial manifold the numbers $`h_{ni}h_i`$, $`i=0,1,\mathrm{},n`$, are homotopy invariants. In particular, they are independent on a triangulation.
In the particular case of PL-manifolds the topological invariance of numbers $`h_{ni}h_i`$ was firstly observed by Pachner in \[Pa, (7.11)\].
###### Example 5.7.
Consider triangulations of the 2-torus $`T^2`$, so $`n=3`$, $`\chi (T^2)=0`$. Since for any $`K^{n1}`$ holds $`\chi (K^{n1})=1+(1)^{n1}h_n`$, in our case we have $`h_3=1`$. Then Corollary 5.4 gives
$$h_3h_0=2,h_2h_1=6.$$
For instance, the triangulation on Figure 6 has $`f_0=9`$ vertices, $`f_1=27`$ edges and $`f_2=18`$ triangles. The corresponding $`h`$-vector is $`(1,6,12,1)`$.
## 6. Concluding remarks
The main goal of the present paper was to establish new connections between topology of manifolds and cellular complexes and combinatorics by means of the notion of a moment-angle complex, introduced by the authors in previous papers \[BP2\]\[BP3\]\[BP4\]. As we have seen, the combinatorics of simplicial manifolds and related objects (polytopes, simplicial spheres, simplicial complexes, coordinate subspace arrangements) can be effectively described by means of topological invariants of bigraded equivariant moment-angle complexes. One of the main properties of a moment-angle complex is the existence of a torus action all of whose isotropy subgroups are coordinate ones. This, in particular, allows to introduce an additional grading to the cohomology ring of the moment-angle complex. On this point one can observe that there is a natural $`/2`$-analogue of almost all constructions presented in this paper. The first step is to replace the torus $`T^m`$ by its “real analogue”, namely, the group $`(/2)^m`$. Then the unit cube $`I^m=[0,1]^m`$ is the orbit space for the action of $`(/2)^m`$ on the bigger cube $`[1,1]^m`$, the “real analogue” of the poly-disk $`(D^2)^m^m`$. Now, starting from any cubical subcomplex $`CI^m`$ one can construct another $`(/2)^m`$-symmetrical cubical complex embedded into $`[1,1]^m^m`$, just in the same way as we did it in Definition 3.1. In particular, for any simplicial complex $`K`$ on $`m`$ vertices one can construct $`(/2)^m`$-symmetrical cubical complexes $`𝒵_K^{}`$, $`𝒲_K^{}`$, the “real analogues” of moment-angle complexes $`𝒵_K`$, $`𝒲_K`$, see (11). The analogue of Lemma 4.2 holds for real coordinate subspace arrangements. Namely, the complement $`U\left(𝒜^{}(K)\right)`$ of the real coordinate subspace arrangement $`𝒜^{}(K)`$ defined by $`K`$ in the same way as in (13) is $`(/2)^m`$-equivariantly homotopy equivalent to $`𝒵_K^{}`$. However, the situation with the cohomology algebra of $`𝒵_K^{}`$ is more subtle: as we have already mentioned, the analogue of Theorem 4.3 does not hold for $`𝒵_K^{}`$, at least for $``$-coefficients (this is a usual thing in topology: the cohomology of “real” objects is more complicated than that of “complex” ones). At the same time $`𝒵_K^{}`$ is an $`m`$-dimensional manifold provided that $`K`$ is a simplicial sphere. So, for any simplicial sphere $`K`$ with $`m`$ vertices we have a $`(/2)^m`$-symmetric manifold with $`(/2)^m`$-invariant cubical complex structure. This class of cubical manifolds may be useful in the cubical analogue of the combinatorial theory of $`f`$-vectors of simplicial complexes.
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# Quantum Teleportation and Beam Splitting
## 1 Introduction
Following the previous paper , we further discuss the non-perfect teleportation. The notion of non-perfect teleportation is introduced in to construct a handy (i.e., physically more realizable) teleportation, although its mathematics becomes a little more complicated. For the completeness of the present paper, we quickly review the meaning of the teleportation and some basic facts of Fock space in this section. Then we dicuss the perfect teleportation in very general (more general than one given in ) scheme with our previous results, and we state the main theorem obtained in for non-perfect teleportation, both in the section 2. The main results of this paper are presented in the section 3, where we discuss the difference among three models, i.e., the perfect model, the non-perfect one given in and that discussed in the present paper. The proofs of the main results are given in the section 4.
### 1.1 Quantum teleportation
The study of quantum teleportation was started by the paper as a part of quantum cryptography , whose scheme can be mathematically expressed in the following steps :
A girl named Alice has an unknown quantum state $`\rho `$ on (a $`N`$–dimensional) Hilbert space $`_1`$ and she was asked to teleport it to a boy named Bob.
For this purpose, we need two other Hilbert spaces $`_2`$ and $`_3`$, $`_2`$ is attached to Alice and $`_3`$ is attached to Bob. Prearrange a so-called entangled state $`\sigma `$ on $`_2_3`$ having certain correlations and prepare an ensemble of the combined system in the state $`\rho \sigma `$ on $`_1_2_3`$.
One then fixes a family of mutually orthogonal projections $`(F_{nm})_{n,m=1}^N`$ on the Hilbert space $`_1_2`$ corresponding to an observable $`F:=\underset{n,m}{}z_{n,m}F_{nm}`$, and for a fixed one pair of indices $`n,m`$, Alice performs a first kind incomplete measurement, involving only the $`_1_2`$ part of the system in the state $`\rho \sigma `$, which filters the value $`z_{nm}`$, that is, after measurement on the given ensemble $`\rho \sigma `$ of identically prepared systems, only those where $`F`$ shows the value $`z_{nm}`$ are allowed to pass. According to the von Neumann rule, after Alice’s measurement, the state becomes
$$\rho _{nm}^{(123)}:=\frac{(F_{nm}\mathrm{𝟏})\rho \sigma (F_{nm}\mathrm{𝟏})}{\mathrm{tr}_{123}(F_{nm}\mathrm{𝟏})\rho \sigma (F_{nm}\mathrm{𝟏})}$$
where $`\mathrm{tr}_{123}`$ is the full trace on the Hilbert space $`_1_2_3`$.
Bob is informed which measurement was done by Alice. This is equivalent to transmit the information that the eigenvalue $`z_{nm}`$ was detected. This information is transmitted from Alice to Bob without disturbance and by means of classical tools.
Making only partial measurements on the third part on the system in the state $`\rho _{nm}^{(123)}`$ means that Bob will control a state $`\mathrm{\Lambda }_{nm}(\rho )`$ on $`_3`$ given by the partial trace on $`_1_2`$ of the state $`\rho _{nm}^{(123)}`$ (after Alice’s measurement)
$`\mathrm{\Lambda }_{nm}(\rho )`$ $`=`$ $`\mathrm{tr}_{12}\rho _{nm}^{(123)}`$
$`=`$ $`\mathrm{tr}_{12}{\displaystyle \frac{(F_{nm}\mathrm{𝟏})\rho \sigma (F_{nm}\mathrm{𝟏})}{\mathrm{tr}_{123}(F_{nm}\mathrm{𝟏})\rho \sigma (F_{nm}\mathrm{𝟏})}}`$
Thus the whole teleportation scheme given by the family $`(F_{nm})`$ and the entangled state $`\sigma `$ can be characterized by the family $`(\mathrm{\Lambda }_{nm})`$ of channels from the set of states on $`_1`$ into the set of states on $`_3`$ and the family $`(p_{nm})`$ given by
$$p_{nm}(\rho ):=\mathrm{tr}_{123}(F_{nm}\mathrm{𝟏})\rho \sigma (F_{nm}\mathrm{𝟏})$$
of the probabilities that Alice’s measurement according to the observable $`F`$ will show the value $`z_{nm}`$.
The teleportation scheme works perfectly with respect to a certain class $`𝔖`$ of states $`\rho `$ on $`_1`$ if the following conditions are fulfilled.
For each $`n,m`$ there exists a unitary operator $`v_{nm}:_1_3`$ such that
$$\mathrm{\Lambda }_{nm}(\rho )=v_{nm}\rho v_{nm}^{}(\rho 𝔖)$$
$$\underset{nm}{}p_{nm}(\rho )=1(\rho 𝔖)$$
(E1) means that Bob can reconstruct the original state $`\rho `$ by unitary keys $`\{v_{nm}\}`$ provided to him.
(E2) means that Bob will succeed to find a proper key with certainty.
Such a teleportation process can be classified into two cases , i.e., weak teleportation and general teleportation, in which the solutions of the teleportation in each case and the conditions of the uniqueness of unitary key were discussed. The solution of the weak teleportation is a triple $`\{\sigma ^{(23)},F^{(12)},U\}`$such that
$$\mathrm{\Lambda }^{}\rho ^{(1)}=U^{}\rho ^{(1)}U$$
holds for any state $`\rho ^{(1)}𝒮(_1)`$ . Once a weak solution of a teleportation problem is given, we can construct the general solution for all $`n,m`$ above .
In , we considered a teleportation model where the entangled state $`\sigma `$ is given by the splitting of a superposition of certain coherent states, although this model doesn’t work perfectly, that is, neither (E2) nor (E1) hold. In the same paper, we estimated the difference between the perfect teleportation and this non-perfect teleportation by adding a further step in the teleportation scheme:
Bob will perform a measurement on his part of the system according to the projection
$$F_+:=\mathrm{𝟏}|\mathrm{exp}(0)><\mathrm{exp}(0)|$$
where $`|\mathrm{exp}(0)><\mathrm{exp}(0)|`$ denotes the vacuum state (the coherent state with density $`0`$).
Then our new teleportation channels (we denote it again by $`\mathrm{\Lambda }_{nm}`$) have the form
$$\mathrm{\Lambda }_{nm}(\rho ):=\mathrm{tr}_{12}\frac{(F_{nm}F_+)\rho \sigma (F_{nm}F_+)}{\mathrm{tr}_{123}(F_{nm}F_+)\rho \sigma (F_{nm}F_+)}$$
and the corresponding probabilities are
$$p_{nm}(\rho ):=\mathrm{tr}_{123}(F_{nm}F_+)\rho \sigma (F_{nm}F_+)$$
For this teleportation scheme, (E1) is fulfilled but (E2) is not, about which we review in the next section.
### 1.2 Basic Notions and Notations
We collect some basic facts concerning the (symmetric) Fock space in a way adapted to the language of counting measures. For details we refer to .
Let $`G`$ be an arbitrary complete separable metric space. Further, let $`\mu `$ be a locally finite diffuse measure on $`G`$, i.e. $`\mu (B)<+\mathrm{}`$ for bounded measurable subsets of $`G`$ and $`\mu (\{x\})=0`$ for all singletons $`xG`$.
We denote the set of all finite counting measures on $`G`$ by $`M=M(G)`$. Since $`\phi M`$ can be written in the form $`\phi =\underset{j=1}{\overset{n}{}}\delta _{x_j}`$ for some $`n=0,1,2,\mathrm{}`$ and $`x_jG`$ with the Dirac measure $`\delta _x`$ corresponding to $`xG`$, the elements of $`M`$ can be interpreted as finite (symmetric) point configurations in $`G`$. We equip $`M`$ with its canonical $`\sigma `$–algebra $`𝔚`$ (cf. , ) and we consider the $`\sigma `$–finite measure $`F`$ by setting
$$F(Y):=𝒳_Y(O)+\underset{n1}{}\frac{1}{n!}\underset{G^n}{}𝒳_Y\left(\underset{j=1}{\overset{n}{}}\delta _{x_j}\right)\mu ^n(d[x_1,\mathrm{},x_n])(Y𝔚),$$
where $`𝒳_Y`$ denotes the indicator function of a set $`Y`$ and $`O`$ represents the empty configuration, i. e., $`O(G)=0`$.
Since $`\mu `$ was assumed to be diffuse one easily checks that $`F`$ is concentrated on the set of a simple configurations (i.e., without multiple points)
$$\widehat{M}:=\{\phi M|\phi (\{x\})1\text{ for all }xG\}$$
$`=(G):=L^2(M,𝔚,F)`$ is called the (symmetric) Fock space over $`G`$.
In it was proved that $``$ and the Boson Fock space $`\mathrm{\Gamma }(L^2(G))`$ in the usual definition are isomorphic. For each $`\mathrm{\Phi }`$ with $`\mathrm{\Phi }0`$ we denote by $`|\mathrm{\Phi }>`$ the corresponding normalized vector
$$|\mathrm{\Phi }>:=\frac{\mathrm{\Phi }}{\mathrm{\Phi }}.$$
Further, $`|\mathrm{\Phi }><\mathrm{\Phi }|`$ denotes the corresponding one–dimensional projection describing a pure state given by the normalized vector $`|\mathrm{\Phi }>`$. Now, for each $`n1`$ let $`^n`$ be the $`n`$–fold tensor product of the Hilbert space $``$, which can be identified with $`L^2(M^n,F^n)`$.
For a given function $`g:G`$ the function $`\mathrm{exp}(g):M`$ defined by
$$\mathrm{exp}(g)(\phi ):=\{\begin{array}{ccc}1\hfill & \text{ if }\hfill & \phi =0\hfill \\ _{xG,\phi \left(\left\{x\right\}\right)>0}g(x)\hfill & & otherwise\hfill \end{array}$$
is called exponential vector generated by $`g`$.
Observe that $`\mathrm{exp}(g)`$ if and only if $`gL^2(G)`$ and one has in this case
$`\mathrm{exp}(g)^2=e^{g^2}`$ and $`|\mathrm{exp}(g)>=e^{\frac{1}{2}g^2}\mathrm{exp}(g)`$. The projection $`|\mathrm{exp}(g)><\mathrm{exp}(g)|`$ is called the coherent state corresponding to $`gL^2(G)`$. In the special case $`g0`$ we get the vacuum state
$$|\mathrm{exp}(0)>=𝒳_{\{0\}}.$$
The linear span of the exponential vectors of $``$ is dense in $``$, so that bounded operators and certain unbounded operators can be characterized by their actions on exponential vectors.
The operator $`D:\mathrm{dom}(D)^2`$ given on a dense domain $`\mathrm{dom}(D)`$ containing the exponential vectors from $``$ by
$$D\psi (\phi _1,\phi _2):=\psi (\phi _1+\phi _2)(\psi \mathrm{dom}(D),\phi _1,\phi _2M)$$
is called compound Malliavin derivative. On exponential vectors $`\mathrm{exp}(g)`$ with $`gL^2(G),`$ one gets immediately
$$D\mathrm{exp}(g)=\mathrm{exp}(g)\mathrm{exp}(g)$$
(1)
The operator $`S:\mathrm{dom}(S)`$ given on a dense domain $`\mathrm{dom}(S)^2`$ containing tensor products of exponential vectors by
$$S\mathrm{\Phi }(\phi ):=\underset{\stackrel{~}{\phi }\phi }{}\mathrm{\Phi }(\stackrel{~}{\phi },\phi \stackrel{~}{\phi })(\mathrm{\Phi }\mathrm{dom}(S),\phi M)$$
is called compound Skorohod integral. One gets
$$D\psi ,\mathrm{\Phi }_^2=\psi ,S\mathrm{\Phi }_{}(\psi \mathrm{dom}(D),\mathrm{\Phi }\mathrm{dom}(S))$$
(2)
$$S(\mathrm{exp}(g)\mathrm{exp}(h))=\mathrm{exp}(g+h)(g,hL^2(G))$$
(3)
For more details we refer to .
Let $`T`$ be a linear operator on $`L^2(G)`$ with $`T1`$. Then the operator $`\mathrm{\Gamma }(T)`$ called second quantization of $`T`$ is the (uniquely determined) bounded operator on $``$ fulfilling
$$\mathrm{\Gamma }(T)\mathrm{exp}(g)=\mathrm{exp}(Tg)(gL^2(G)).$$
Clearly, it holds
$`\mathrm{\Gamma }(T_1)\mathrm{\Gamma }(T_2)`$ $`=`$ $`\mathrm{\Gamma }(T_1T_2)`$ (4)
$`\mathrm{\Gamma }(T^{})`$ $`=`$ $`\mathrm{\Gamma }(T)^{}`$
It follows that $`\mathrm{\Gamma }(T)`$ is an unitary operator on $``$ if $`T`$ is an unitary operator on $`L^2(G)`$.
In we proved.
###### LEMMA 1.1
Let $`K_1,K_2`$ be linear operators on $`L^2(G)`$ with property
$$K_1^{}K_1+K_2^{}K_2=\mathrm{𝟏}.$$
(5)
Then there exists exactly one isometry $`\nu _{K_1,K_2}`$ from $``$ to $`^2=`$ with
$$\nu _{K_1,K_2}\mathrm{exp}(g)=\mathrm{exp}(K_1g)\mathrm{exp}(K_2g)(gL^2(G))$$
(6)
Further it holds
$$\nu _{K_1,K_2}=(\mathrm{\Gamma }(K_1)\mathrm{\Gamma }(K_2))D$$
(7)
(at least on $`\mathrm{dom}(D)`$ but one has the unique extension).
The adjoint $`\nu _{K_1,K_2}^{}`$ of $`\nu _{K_1,K_2}`$ is characterized by
$$\nu _{K_1,K_2}^{}(\mathrm{exp}(h)\mathrm{exp}(g))=\mathrm{exp}(K_1^{}h+K_2^{}g)(g,hL^2(G))$$
(8)
and it holds
$$\nu _{K_1,K_2}^{}=S(\mathrm{\Gamma }(K_1^{})\mathrm{\Gamma }(K_2^{}))$$
(9)
From $`K_1,K_2`$ we get a transition expectation $`\xi _{K_1K_2}:`$, using $`\nu _{K_1,K_2}`$ and the lifting $`\xi _{K_1K_2}^{}`$ may be interpreted as a certain splitting (cf. ). The property (5) implies
$$K_1g^2+K_2g^2=g^2(gL^2(G))$$
(10)
Let $`U`$, $`V`$ be unitary operators on $`L^2(G)`$. If operators $`K_1,K_2`$ satisfy (5), then the pair $`\widehat{K}_1=UK_1,\widehat{K}_2=VK_2`$ fulfill (5).
Here we explain fundamental scheme of beam splitting . We define an isometric operator $`V_{\alpha ,\beta }`$ for coherent vectors such that
$$V_{\alpha ,\beta }|\mathrm{exp}(g)=|\mathrm{exp}(\alpha g)|\mathrm{exp}(\beta g)$$
with $`\alpha ^2+\beta ^2=1`$. This beam splitting is a useful mathematical expression for optical communication and quantum measurements . As one example, take $`\alpha =\beta =\frac{1}{2}`$ in the above formula and let $`K_1=K_2`$ be the following operator of multiplication on $`L^2(G)`$
$$K_1g=\frac{1}{\sqrt{2}}g=K_2g(gL^2(G))$$
We put
$$\nu :=\nu _{K_1,K_2},$$
then we obtain
$$\nu \mathrm{exp}(g)=\mathrm{exp}\left(\frac{1}{\sqrt{2}}g\right)\mathrm{exp}(\frac{1}{\sqrt{2}}g)(gL^2(G)).$$
(11)
Another example is given by taking $`K_1`$ and $`K_2`$ as the projections to the corresponding subspaces $`_1,_2`$ of the orthogonal sum $`L^2(G)=_1_2`$.
In we used the first example in order to describe a teleportation model where Bob performs his experiments on the same ensemble of the systems like Alice. Further we used a special case of the second example in order to describe a teleportation model where Bob and Alice are spatially separated (cf. section 5 of ).
## 2 Previous results on teleportation
Let us review some results obtained in . We fix an ONS $`\{g_1,\mathrm{},g_N\}L^2(G)`$, operators $`K_1,K_2`$ on $`L^2(G)`$ with (5), an unitary operator $`T`$ on $`L^2(G)`$, and $`d>0`$. We assume
$$TK_1g_k=K_2g_k(k=1,\mathrm{},N),$$
(12)
$$K_1g_k,K_1g_j=0(kj;k,j=1\mathrm{},N),$$
(13)
Using (11) and (12) we get
$$K_1g_k^2=K_2g_k^2=\frac{1}{2}.$$
(14)
From (12) and (13) we get
$$K_2g_k,K_2g_j=0(kj;k,j=1,\mathrm{},N).$$
(15)
The state of Alice asked to teleport is of the type
$$\rho =\underset{s=1}{\overset{N}{}}\lambda _s|\mathrm{\Phi }_s\mathrm{\Phi }_s|,$$
(16)
where
$$|\mathrm{\Phi }_s=\underset{j=1}{\overset{N}{}}c_{sj}|\mathrm{exp}(aK_1g_j)\mathrm{exp}(0)\left(\underset{j}{}|c_{sj}|^2=1;s=1,\mathrm{},N\right)$$
(17)
and $`a=\sqrt{d}`$. One easily checks that $`(|\mathrm{exp}(aK_1g_j)\mathrm{exp}(0))_{j=1}^N`$ and $`(|\mathrm{exp}aK_2g_j)\mathrm{exp}(0))_{j=1}^N`$ are ONS in $``$. The set $`\left\{\mathrm{\Phi }_s;s=1,\mathrm{},N\right\}`$ makes the $`N`$-dimensional Hilbert space $`_1`$ defining an input state teleported by Alice. We may include the vaccum state $`|\mathrm{exp}(0)`$ to define $`_1,`$ however we take the $`N`$-dimensional Hilbert space $`_1`$ as above because of computational simplicity.
In order to achieve that $`(|\mathrm{\Phi }_s)_{s=1}^N`$ is still an ONS in $``$ we assume
$$\underset{j=1}{\overset{N}{}}\overline{c}_{sj}c_{kj}=0(jk;j,k=1,\mathrm{},N).$$
(18)
Denote $`c_s=[c_{s1,\mathrm{},}c_{sN}]^N`$, then $`(c_s)_{s=1}^N`$ is an CONS in $`^N`$.
Let $`(b_n)_{n=1}^N`$ be a sequence in $`^N`$,
$$b_n=[b_{n1,\mathrm{},}b_{nN}]$$
with properties
$$|b_{nk}|=1(n,k=1,\mathrm{},N),$$
(19)
$$b_n,b_j=0(nj;n,j=1,\mathrm{},N).$$
(20)
Now, for each $`m,n(=1,\mathrm{},N),`$ we have unitary operators $`U_m,B_n`$ on $``$ given by
$$B_n|\mathrm{exp}(aK_1g_j)\mathrm{exp}(0)=b_{nj}|\mathrm{exp}(aK_1g_j)\mathrm{exp}(0)(j=1,\mathrm{},N)$$
(21)
$$U_m|\mathrm{exp}(aK_1g_j)\mathrm{exp}(0)=|\mathrm{exp}(aK_1g_{jm})\mathrm{exp}(0)(j=1,\mathrm{},N)$$
(22)
### 2.1 A perfect teleportation
Then Alice’s measurements are performed with projection
$$F_{nm}=|\xi _{nm}\xi _{nm}|(n,m=1,\mathrm{},N)$$
(23)
given by
$$|\xi _{nm}=\frac{1}{\sqrt{N}}\underset{j=1}{\overset{N}{}}b_{nj}|\mathrm{exp}(aK_1g_j)\mathrm{exp}(0)>|\mathrm{exp}(aK_1g_{jm})\mathrm{exp}(0),$$
(24)
where $`jm:=j+m(\mathrm{mod}N)`$.
One easily checks that $`(|\xi _{nm})_{n,m=1}^N`$ is an ONS in $`^2`$. Further, the state vector $`|\xi `$ of the entangled state $`\sigma =|\xi \xi |`$ is given by
$$|\xi =\frac{1}{\sqrt{N}}\underset{k}{}|\mathrm{exp}(aK_1g_k)\mathrm{exp}(0)|\mathrm{exp}(aK_2g_k)\mathrm{exp}(0).$$
(25)
In we proved the following theorem.
###### THEOREM 2.1
For each $`n,m=1,\mathrm{},N`$, define a channel $`\mathrm{\Lambda }_{nm}`$ by
$$\mathrm{\Lambda }_{nm}(\rho ):=\mathrm{tr}_{12}\frac{\left(F_{nm}1\right)(\rho \sigma )\left(F_{nm}\mathrm{𝟏}\right)}{\mathrm{tr}_{123}\left(F_{nm}1\right)\left(\rho \sigma \right)\left(F_{nm}\mathrm{𝟏}\right)}(\rho \text{ normal state on })$$
(26)
Then we have for all states $`\rho `$ on $`M`$ with (16) and (17)
$$\mathrm{\Lambda }_{nm}(\rho )=\left(\mathrm{\Gamma }(T)U_mB_n^{}\right)\rho \left(\mathrm{\Gamma }(T)U_mB_n^{}\right)^{}$$
(27)
###### REMARK 2.2
Using the operators $`B_n,U_m,\mathrm{\Gamma }(T),`$ the projections $`F_{nm}`$ are given by unitary transformations of the entangled state $`\sigma `$ :
$`F_{nm}`$ $`=`$ $`\left(B_nU_m\mathrm{\Gamma }(T^{})\right)\sigma \left(B_nU_m\mathrm{\Gamma }(T^{})\right)^{}`$
or
$`|\xi _{nm}`$ $`=`$ $`\left(B_nU_m\mathrm{\Gamma }(T^{})\right)|\xi .`$
If Alice performs a measurement according to the following selfadjoint operator
$$F=\underset{n,m=1}{\overset{N}{}}z_{nm}F_{nm}$$
with $`\{z_{nm}|n,m=1,\mathrm{},N\}𝐑\{0\},`$ then she will obtain the value $`z_{nm}`$ with probability $`1/N^2`$. The sum over all this probabilities is $`1`$, so that the teleportation model works perfectly.
Before stating some fundamental results in for non-perfect case, we note that our perfect teleportation is obviously treated in general finite Hilbert spaces $`_k\left(k=1,2,3\right)`$ same as usual one . Moreover, our teleportation scheme can be a bit generalized by introducing the entagled state $`\sigma _{12}`$ on $`_1_2`$ defining the projections $`\left\{F_{nm}\right\}`$ by the unitary operators $`B_n,U_m.`$ We here discuss the perfect teleportation on general Hilbert spaces $`_k\left(k=1,2,3\right).`$ Let $`\left\{\xi _j^k;j=1,\mathrm{},N\right\}`$ be CONS of the Hilbert space $`_k\left(k=1,2,3\right).`$ Define the entangled states $`\sigma _{12}`$ and $`\sigma _{23}`$ on $`_1_2`$ and $`_2_3,`$ respectively, such as
$$\sigma _{12}=|\xi _{12}\xi _{12}|,\text{ }\sigma _{23}=|\xi _{23}\xi _{23}|$$
with $`\xi _{12}\frac{1}{\sqrt{N}}_{j=1}^N\xi _j^1\xi _j^2`$ and $`\xi _{23}\frac{1}{\sqrt{N}}_{j=1}^N\xi _j^2\xi _j^3.`$ By a sequence $`\left\{b_n=[b_{n1,\mathrm{},}b_{nN}];n=1,\mathrm{},N\right\}`$ in $`^N`$ with the properties (19) and (20), we define the unitary operator $`B_n`$ and $`U_m`$ such as
$$B_n\xi _j^1b_{nj}\xi _j^1(n,j=1,\mathrm{},N)\text{ and}U_m\xi _j^2\xi _{jm}^2(n,j=1,\mathrm{},N)$$
with $`jmj+m`$ (mod $`N).`$ Then the set $`\left\{F_{nm};n,m=1,\mathrm{},N\right\}`$ of the projections of Alice is given by
$$F_{nm}=\left(B_nU_m\right)\sigma _{12}\left(B_nU_m\right)^{}$$
and the teleportation channels $`\left\{\mathrm{\Lambda }_{nm}^{};n,m=1,\mathrm{},N\right\}`$ are defined as
$$\mathrm{\Lambda }_{nm}(\rho ):=\mathrm{tr}_{12}\frac{\left(F_{nm}1\right)(\rho \sigma _{23})\left(F_{nm}\mathrm{𝟏}\right)}{\mathrm{tr}_{123}\left(F_{nm}1\right)\left(\rho \sigma _{23}\right)\left(F_{nm}\mathrm{𝟏}\right)}(\rho \text{ normal state on }_1).$$
Finally the unitary keys $`\left\{W_{nm};n,m=1,\mathrm{},N\right\}`$ of Bob are given as
$$W_{nm}\xi _j^1=\overline{b}_{nj}\xi _{jm}^3,\text{ }(n,m=1,\mathrm{},N)$$
by which we obtain the perfect teleportation
$$\mathrm{\Lambda }_{nm}(\rho )=W_{nm}\rho W_{nm}^{}.$$
The above perfect teleportation is unique in the sense of unitary equivalence.
### 2.2 A non–perfect teleportation
We will review a non-perfect teleportation model in which the probability teleporting the state from Alice to Bob is less than $`1`$ and it depends on the density parameter $`d`$ (may be the energy of the beams) of the coherent vector.There, when $`d=a^2`$ tends to infinity, the probability tends to $`1`$. Thus the model can be considered as asymptotically perfect.
Take the normalized vector
$`|\eta :=`$ $`{\displaystyle \frac{\gamma }{\sqrt{N}}}{\displaystyle \underset{k=1}{\overset{N}{}}}|\mathrm{exp}(ag_k)`$ (29)
$`\text{with }\gamma :=`$ $`\left({\displaystyle \frac{1}{1+(N1)e^d}}\right)^{\frac{1}{2}}=\left({\displaystyle \frac{1}{1+(N1)e^{a^2}}}\right)^{\frac{1}{2}}`$
and we replace in (26) the entangled state $`\sigma `$ by
$`\stackrel{~}{\sigma }:=`$ $`|\stackrel{~}{\xi }\stackrel{~}{\xi }|`$ (30)
$`\stackrel{~}{\xi }:=`$ $`\nu _{K_1,K_2}(\eta )={\displaystyle \frac{\gamma }{\sqrt{N}}}{\displaystyle \underset{k=1}{\overset{N}{}}}|\mathrm{exp}(aK_1g_k)|\mathrm{exp}(aK_2g_k)`$
Then for each $`n,m=1,\mathrm{},N,`$ we get the channels on any normal state $`\rho `$ on $``$ such as
$`\stackrel{~}{\mathrm{\Lambda }}_{nm}(\rho ):=`$ $`\mathrm{tr}_{12}{\displaystyle \frac{\left(F_{nm}\mathrm{𝟏}\right)\left(\rho \stackrel{~}{\sigma }\right)\left(F_{nm}\mathrm{𝟏}\right)}{\mathrm{tr}_{123}\left(F_{nm}\mathrm{𝟏}\right)\left(\rho \stackrel{~}{\sigma }\right)\left(F_{nm}\mathrm{𝟏}\right)}}`$ (31)
$`\mathrm{\Theta }_{nm}(\rho ):=`$ $`\mathrm{tr}_{12}{\displaystyle \frac{\left(F_{nm}F_+\right)\left(\rho \stackrel{~}{\sigma }\right)\left(F_{nm}F_+\right)}{\mathrm{tr}_{123}\left(F_{nm}F_+\right)\left(\rho \stackrel{~}{\sigma }\right)\left(F_{nm}F_+\right)}},`$ (32)
where $`F_+=\mathrm{𝟏}|\mathrm{exp}(0)\mathrm{exp}(0)|,`$ i.e.., $`F_+`$ is the projection onto the space $`_+`$ of configurations having no vacuum part;
$$_+:=\{\psi |\mathrm{exp}(0)\mathrm{exp}(0)|\psi =0\}$$
One easily checks that
$$\mathrm{\Theta }_{nm}(\rho )=\frac{F_+\stackrel{~}{\mathrm{\Lambda }}_{nm}(\rho )F_+}{\mathrm{tr}\left(F_+\stackrel{~}{\mathrm{\Lambda }}_{nm}(\rho )F_+\right)}$$
(33)
that is, after receiving the state $`\stackrel{~}{\mathrm{\Lambda }}_{nm}(\rho )`$ from Alice, Bob has to omit the vacuum.
From Theorem 2.1 it follows that for all $`\rho `$ with (16) and (17)
$$\mathrm{\Lambda }_{nm}(\rho )=\frac{F_+\mathrm{\Lambda }_{nm}(\rho )F_+}{\mathrm{tr}(F_+\mathrm{\Lambda }_{nm}(\rho )F_+)}.$$
This is not true if we replace $`\mathrm{\Lambda }_{nm}`$ by $`\stackrel{~}{\mathrm{\Lambda }}_{nm}`$, namely, in general it does not hold
$$\mathrm{\Theta }_{nm}(\rho )=\stackrel{~}{\mathrm{\Lambda }}_{nm}(\rho )$$
In we proved the following theorem.
###### THEOREM 2.3
For all states $`\rho `$ on $``$ with (16) and (17) and each pair $`n,m(=1,\mathrm{},N),`$ we have
$$\mathrm{\Theta }_{nm}(\rho )=\left(\mathrm{\Gamma }\left(T\right)U_mB_n^{}\right)\rho \left(\mathrm{\Gamma }\left(T\right)U_mB_n^{}\right)^{}\text{ or }\mathrm{\Theta }_{nm}(\rho )=\mathrm{\Lambda }_{nm}(\rho )$$
(34)
and
$$\underset{n,m}{}p_{nm}(\rho )=\underset{n,m}{}\mathrm{tr}_{123}\left(F_{nm}F_+\right)\left(\rho \stackrel{~}{\sigma }\right)\left(F_{nm}F_+\right)=\frac{\left(1e^{\frac{d}{2}}\right)^2}{1+(N1)e^d}.$$
(35)
That is, the model works only asymptotically perfect in the sense of condition (E2). With other words, in the case of high density (or energy) of the considered beams the model works perfectly.
## 3 Main results
The tools of the teleportation model considered in section 2.1 are the entangled state $`\sigma `$ and the family of projections $`(F_{nm})_{n,m=1}^N`$. In order to have a more handy model, in section 2.2. we have replaced the entangled state $`\sigma `$ by another entangled state $`\stackrel{~}{\sigma }`$ given by the splitting of a superposition of certain coherent states (30). In addition now we are going to replace the projectors $`F_{nm}`$ by projectors $`\stackrel{~}{F}_{nm}`$ defined as follows.
$$\stackrel{~}{F}_{nm}:=\left(B_nU_m\mathrm{\Gamma }(T)^{}\right)\stackrel{~}{\sigma }\left(B_nU_m\mathrm{\Gamma }(T)^{}\right)^{}$$
(36)
In order to make this definition precise we assume, in addition to (22 ), that is holds:
$$U_m\mathrm{exp}(0)=\mathrm{exp}(0)(m=1,\mathrm{},N)$$
Together with (22) that implies
$$U_m|\mathrm{exp}(aK_1g_j)=|\mathrm{exp}(aK_1g_{jm})(m,j=1,\mathrm{},N)$$
(37)
Formally we have the same relation between $`\stackrel{~}{\sigma }`$ and $`\stackrel{~}{F}_{nm}`$ like the relation between $`\sigma `$ and $`F_{nm}`$ (cf. Remark 2.2). Further for each pair $`n,m=1,\mathrm{},N`$ we define channels on normal states on $``$ such as
$$\stackrel{~}{\mathrm{\Theta }}_{nm}(\rho ):=\mathrm{tr}_{12}\frac{\left(\stackrel{~}{F}_{nm}F_+\right)\left(\rho \stackrel{~}{\sigma }\right)\left(\stackrel{~}{F}_{nm}F_+\right)}{\stackrel{~}{p}_{nm}(\rho )}$$
(38)
where
$$\stackrel{~}{p}_{nm}(\rho ):=\mathrm{tr}_{123}\left(\stackrel{~}{F}_{nm}F_+\right)\left(\rho \stackrel{~}{\sigma }\right)\left(\stackrel{~}{F}_{nm}F_+\right)$$
(39)
(cf. (33), and (34)).
In section 4, we will prove the following theorem.
###### THEOREM 3.1
For each state $`\rho `$ on $``$ with (16), and (17), each pair $`n,m(=1,\mathrm{},N)`$ and each bounded operator $`A`$ on $``$ it holds
$$|\mathrm{tr}\left(\stackrel{~}{\mathrm{\Theta }}_{nm}(\rho )A\right)\mathrm{tr}\left(\mathrm{\Lambda }_{nm}(\rho )A\right)|\frac{2e^{\frac{d}{2}}}{\left(1e^{\frac{d}{2}}\right)}\left(N^2+N\sqrt{N}+N\right)$$
(40)
$$\left|\stackrel{~}{p}_{nm}(\rho )\frac{1}{N^2}\right|e^{\frac{d}{2}}\left(\frac{14}{N^2}+2+\frac{2}{\sqrt{N}}\right)$$
(41)
From Theorem 2.1 and $`e^{\frac{d}{2}}0\left(d+\mathrm{}\right),`$ the theorem 3.1 means that our modified teleportation model works asymptotically perfect (case of high density or energy) in the sense of conditions (E1), and (E2).
In order to obtain a deeper understanding of the whole procedure we are going to discuss another representation of the projectors $`\stackrel{~}{F}_{nm}`$ and of the channels $`\stackrel{~}{\mathrm{\Theta }}_{nm}`$. Starting point is again the normalized vector $`|\eta `$ given by (29). From (14) we obtain
$$O_\sqrt{2}K_1g_k^2=g_k^2,$$
(42)
where $`O_f`$ denotes the operator of multiplication corresponding to the number (or function) $`f`$
$$O_f\psi :=f\psi \left(\psi L_2(G)\right)$$
(43)
Furthermore (13) implies
$$O_fK_1g_k,O_fK_1g_j=0(kj)$$
(44)
From (42), and (44) follows that we have a normalized vector $`|\stackrel{~}{\eta }`$ given by
$$|\stackrel{~}{\eta }:=\mathrm{\Gamma }\left(O_\sqrt{2}K_1\right)|\eta =\frac{\gamma }{\sqrt{N}}\underset{k=1}{\overset{N}{}}|\mathrm{exp}\left(a\sqrt{2}K_1g_k\right)$$
(45)
###### REMARK 3.2
In the case of Example LABEL:def7 we have
$$|\stackrel{~}{\eta }=|\eta $$
Now let $`V`$ be the unitary operator on $``$ characterized by
$`V\left(\mathrm{exp}(f_1)\mathrm{exp}(f_2)\right)`$ (46)
$`=`$ $`\mathrm{exp}\left({\displaystyle \frac{1}{\sqrt{2}}}\left(f_1f_2\right)\right)\mathrm{exp}\left({\displaystyle \frac{1}{\sqrt{2}}}\left(f_1+f_2\right)\right)\left(f_1,f_2L_2(G)\right)`$
On easily checks
$`V^{}\left(\mathrm{exp}(f_1)\mathrm{exp}(f_2)\right)`$ (47)
$`=`$ $`\mathrm{exp}\left({\displaystyle \frac{1}{\sqrt{2}}}\left(f_1+f_2\right)\right)\mathrm{exp}\left({\displaystyle \frac{1}{\sqrt{2}}}\left(f_2f_1\right)\right)\left(f_1,f_2L_2(G)\right)`$
###### REMARK 3.3
$`V`$ describes a certain exchange procedure of particles (or energy) between two systems or beams (cf. )
Now, using (12), (30), (45), and (47), resp. (46) one gets
$`\stackrel{~}{\xi }`$ $`=`$ $`\nu _{K_1,K_2}(\eta )=(\mathrm{𝟏}\mathrm{\Gamma }(T))V^{}\left(|\mathrm{exp}(0)|\stackrel{~}{\eta }\right)`$ (48)
$`\stackrel{~}{\xi }`$ $`=`$ $`(\mathrm{𝟏}\mathrm{\Gamma }(T))V\left(|\stackrel{~}{\eta }|\mathrm{exp}(0)\right)`$ (49)
and it follows
$`\stackrel{~}{\sigma }`$ $`=`$ $`|\stackrel{~}{\xi }\stackrel{~}{\xi }|`$ (50)
$`=`$ $`(\mathrm{𝟏}\mathrm{\Gamma }(T))V^{}\left(|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|\right)\left((\mathrm{𝟏}\mathrm{\Gamma }(T))V^{}\right)^{}`$
$$\stackrel{~}{\sigma }=(\mathrm{𝟏}\mathrm{\Gamma }(T))V\left(|\stackrel{~}{\eta }\stackrel{~}{\eta }||\mathrm{exp}(0)\mathrm{exp}(0)|\right)((\mathrm{𝟏}\mathrm{\Gamma }(T))V)^{}$$
(51)
From the definition of $`\stackrel{~}{F}_{nm}`$ (36) and (50) it follows
$$\stackrel{~}{F}_{nm}=\left(B_nU_m\right)V^{}\left(|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|\right)\left(\left(B_nU_m\right)V^{}\right)^{}$$
(52)
Using (51), and (52) we obtain
$`(\stackrel{~}{F}_{nm}F_+)(\rho \stackrel{~}{\sigma })(\stackrel{~}{F}_{nm}F_+)(n,m=1,\mathrm{},N)`$ (53)
$`=`$ $`\left(X_{nm}\mathrm{𝟏}\right)W_{nm}\left(\rho |\stackrel{~}{\eta }\stackrel{~}{\eta }||\mathrm{exp}(0)\mathrm{exp}(0)|\right)W_{nm}^{}\left(X_{nm}\mathrm{𝟏}\right)^{}`$
where
$$\begin{array}{cc}\hfill X_{nm}& :=(B_nU_m)V^{}(n,m=1,\mathrm{},N)\hfill \\ \hfill W_{nm}& :=\left(|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|F_+\right)\left(V\mathrm{𝟏}\right)\left(B_n^{}U_m^{}\mathrm{\Gamma }(T)\right)\left(\mathrm{𝟏}V\right)\hfill \end{array}$$
(54)
$`X_{nm}`$ and consequently $`X_{nm}\mathrm{𝟏}`$ are unitary operators. For that reason we get from (53)
$`\mathrm{tr}_{123}(\stackrel{~}{F}_{nm}F_+)(\rho \stackrel{~}{\sigma })(\stackrel{~}{F}_{nm}F_+)(n,m=1,\mathrm{},N)`$ (55)
$`=`$ $`\mathrm{tr}_{123}W_{nm}\left(\rho |\stackrel{~}{\eta }\stackrel{~}{\eta }|\mathrm{exp}(0)\mathrm{exp}(0)|\right)W_{nm}^{}`$
and
$`\mathrm{tr}_{12}\left(\stackrel{~}{F}_{nm}F_+\right)\left(\rho \stackrel{~}{\sigma }\right)\left(\stackrel{~}{F}_{nm}F_+\right)`$ (56)
$`=`$ $`\mathrm{tr}_{12}W_{nm}\left(\rho |\stackrel{~}{\eta }\stackrel{~}{\eta }||\mathrm{exp}(0)\mathrm{exp}(0)|\right)W_{nm}^{}`$
Now from (38), (39), (55) and (56) it follows
$`\stackrel{~}{p}_{nm}(\rho )`$ $`=`$ $`\mathrm{tr}_{123}W_{nm}\left(\rho |\stackrel{~}{\eta }\stackrel{~}{\eta }||exp(0)\mathrm{exp}(0)|\right)W_{nm}^{}`$ (57)
$`\stackrel{~}{\mathrm{\Theta }}_{nm}(\rho )`$ $`=`$ $`tr_{12}{\displaystyle \frac{W_{nm}\left(\rho |\stackrel{~}{\eta }\stackrel{~}{\eta }||\mathrm{exp}(0)\mathrm{exp}(0)|\right)W_{nm}^{}}{\mathrm{tr}_{123}W_{nm}\left(\rho |\stackrel{~}{\eta }\stackrel{~}{\eta }||\mathrm{exp}(0)\mathrm{exp}(0)|\right)W_{nm}^{}}}`$ (58)
According to (57,58) and (54), the procedure of the special teleportation model can be expressed in the following steps.
| Step 0 –initial state | $`s_{\mathrm{in}}(\rho )=\rho |\stackrel{~}{\eta }\stackrel{~}{\eta }||\mathrm{exp}(0)\mathrm{exp}(0)|`$ |
| --- | --- |
| $`\rho `$–the unknown state | $`|`$ |
| Alice want to teleport | $`|`$ |
| $`|\mathrm{exp}(0)\mathrm{exp}(0)|`$–vacuum state, | $`|`$ |
| Bobs state at the beginning | $`|`$ |
| | $`|`$ |
| Step 1 –Transformation according to | $`\mathrm{𝟏}V`$ |
| that means: splitting of | $`|`$ |
| the state $`|\stackrel{~}{\eta }\stackrel{~}{\eta }|`$ | $`|`$ |
| | $`|`$ |
| Step 2 –Transformation according to | $`B_n^{}U_m^{}\mathrm{\Gamma }(T)`$ |
| | $`|`$ |
| | $`|`$ |
| Step 3 –Transformation according to | $`V\mathrm{𝟏}`$ |
| exchange of particles (or energy) | $`|`$ |
| between the first and the second | $`|`$ |
| part of the system | $`|`$ |
| | $`|`$ |
| Step 4 –measurement according to | $`|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|F_+`$ |
| checking for | $`|`$ |
| \- first part in the vacuum? | $`|`$ |
| \- in the third part is no vacuum? | $`|`$ |
| \- second part reconstructed? | $`|`$ |
| | $``$ |
| Final state $`s_{\mathrm{fin}}(\rho )`$ | $`=\frac{W_{nm}\left(s_{\mathrm{in}}(\rho )\right)W_{nm}^{}}{\mathrm{tr}_{123}W_{nm}\left(s_{\mathrm{in}}(\rho )\right)W_{nm}^{}}`$ |
Now from (57) we get $`\stackrel{~}{\mathrm{\Theta }}_{nm}(\rho )=\mathrm{tr}_{12}s_{\mathrm{fin}}(\rho )`$. Thus theorem 3.1 means that in the case of high density (or energy) $`d`$ we have approximately ($`\rho `$ with (16), and (17))
$$\mathrm{tr}_{12}s_{\mathrm{fin}}(\rho )=\left(\mathrm{\Gamma }(T)U_mB_n^{}\right)\rho \left(\mathrm{\Gamma }(T)U_mB_n^{}\right)^{}$$
The proof of theorem 3.1 shows that we have even more, namely it holds (approximately)
$$s_{\mathrm{fin}}(\rho )=|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|\left(\mathrm{\Gamma }(T)U_mB_n^{}\right)\rho \left(\mathrm{\Gamma }(T)U_mB_n^{}\right)^{}$$
(59)
Adding in our scheme the following step
Step 5 –Transformation $`\mathrm{𝟏}\mathrm{𝟏}\left(\mathrm{\Gamma }(T)U_mB_n^{}\right)^{}`$
(that means Bob uses the key provided to him)
Then $`s_{\mathrm{fin}}(\rho )`$ will change into the new final state
$$|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|\rho $$
Summarizing one can describe the effect of the procedure (for large $`d`$!) as follows: At the beginning Alice has (e. g., can control) a state $`\rho `$, and Bob has the vacuum state (e. g., can control nothing). After the procedure Bob has the state $`\rho `$ and Alice has the vacuum. Furthermore the teleportation mechanism is ready for the next teleportation (e. g. $`|\stackrel{~}{\eta }\stackrel{~}{\eta }|`$ is reproduced in the course of teleportation).
We have considered three different models (cf. sections 2.1, 2.2, 2.3). Each of them is a special case of a more general concept we are going to describe in the following:
Let $`H_1,H_2`$ be $`N`$–dimensional subspaces of $`_+`$ such that $`\mathrm{\Gamma }(T)`$ maps $`H_1`$ onto $`H_2`$, and $`H_1`$ is invariant with respect to the unitary transformations $`B_n`$, $`U_m`$ $`(n,m=1,\mathrm{},N)`$.
Further let $`\sigma _1`$, $`\sigma _2`$ be projections of the type
$$\sigma _k=|\xi _k\xi _k|,\xi _k(H_1_0)(H_2_0)(k=1,2)$$
where $`_0`$ is the orthogonal complement of $`_+`$, e. g., $`_0`$ is the one-dimensional subspace of $``$ spanned by the vacuum vector $`|\mathrm{exp}(0)`$.
Now for each $`n,m=1,\mathrm{},N`$ and each pair $`\sigma _1`$, $`\sigma _2`$ we define a channel $`\mathrm{\Omega }_{nm}^{\sigma _1,\sigma _2}`$ from the set of all normal states $`\rho `$ on $`H_1`$ into the set of all normal states on $`_+`$
$$\mathrm{\Omega }_{nm}^{\sigma _1\sigma _2}(\rho ):=\mathrm{tr}_{12}\frac{\left(F_{nm}^{\sigma _1}F_+\right)\left(\rho \sigma _2\right)\left(F_{nm}^{\sigma _1}F_+\right)}{\mathrm{tr}_{123}\left(F_{nm}^{\sigma _1}F_+\right)\left(\rho \sigma _2\right)\left(F_{nm}^{\sigma _1}F_+\right)}$$
where
$$F_{nm}^{\sigma _1}:=\left(B_nU_m\mathrm{\Gamma }(T^{})\right)\sigma _1\left(B_nU_m\mathrm{\Gamma }(T^{})\right)^{}$$
In this paper we have considered the situation where $`H_1`$ is spanned by the ONS
$$\left(|\mathrm{exp}(aK_1g_k)\mathrm{exp}(0)\right)_{k=1}^N$$
and $`H_2`$ is spanned by the ONS
$$\left(|\mathrm{exp}(aK_2g_k)\mathrm{exp}(0)\right)_{k=1}^N$$
Further the model discussed in section 2.2 corresponds to the special case $`\sigma _1=\sigma _2=\sigma `$, e. g.
$$\mathrm{\Lambda }_{nm}=\mathrm{\Omega }_{nm}^{\sigma \sigma }(n,m=1,\mathrm{},N)$$
(perfect in the sense of conditions (E1) and (E2)).
The model discussed in section 2.2 corresponds to the special case $`\sigma _1=\sigma \sigma _2=\stackrel{~}{\sigma }`$, e. g.
$$\mathrm{\Theta }_{nm}=\mathrm{\Omega }_{nm}^{\sigma \stackrel{~}{\sigma }}$$
(perfect in the sense of (E1), and only asymptotically perfect in the sense of (E2)).
Finally the model from this section corresponds to the special case $`\sigma _1=\sigma _2=\stackrel{~}{\sigma }`$, e. g.
$$\stackrel{~}{\mathrm{\Theta }}_{nm}=\mathrm{\Omega }_{nm}^{\stackrel{~}{\sigma }\stackrel{~}{\sigma }}$$
(non-perfect, neither (E2) nor (E1) hold, but asymptotically perfect in the sense of both conditions)
## 4 Proof of Theorem 3.1
From (14) we get
$$\mathrm{exp}\left(aK_sg_j\right)\mathrm{exp}(0)^2=e^{\frac{a^2}{2}}1(s=1,2;j=1,\mathrm{},N)$$
(60)
$$\mathrm{exp}\left(aK_sg_j\right)=e^{\frac{a^2}{2}}(s=1,2;j=1,\mathrm{},N)$$
(61)
Using (46), (60) and (61) one easily checks
$`V(|\mathrm{exp}\left(aK_1g_j\right))\mathrm{exp}(0)|\mathrm{exp}\left(aK_1g_k\right)`$ (62)
$`=\left((e^{\frac{a^2}{2}}1)e^{\frac{a^2}{2}}\right)^{\frac{1}{2}}[\mathrm{exp}\left({\displaystyle \frac{a}{\sqrt{2}}}K_1(g_jg_k)\right)\mathrm{exp}\left({\displaystyle \frac{a}{\sqrt{2}}}K_1(g_j+g_k)\right)`$
$`\mathrm{exp}({\displaystyle \frac{a}{\sqrt{2}}}K_1\left(g_k\right))\mathrm{exp}\left({\displaystyle \frac{a}{\sqrt{2}}}K_1g_k\right)]`$
$`(k,j=1,\mathrm{},N)`$
###### LEMMA 4.1
Put for $`j,k=1,\mathrm{},N`$
$$\alpha _{jk}:=|\mathrm{exp}(0)|\stackrel{~}{\eta },V\left(|\mathrm{exp}\left(aK_1g_j\right)\mathrm{exp}(0)|\mathrm{exp}\left(aK_1g_k\right)\right)$$
Then it holds for $`j,k=1,\mathrm{},N`$
$$\alpha _{jk}=\left(\left(1e^{\frac{a^2}{2}}\right)e^{a^2}\right)^{\frac{1}{2}}\frac{\gamma }{\sqrt{N}}(kj)$$
(63)
$$\alpha _{jj}=\left(1e^{\frac{a^2}{2}}\right)^{\frac{1}{2}}\frac{\gamma }{\sqrt{N}}$$
(64)
Proof: We have
$$\mathrm{exp}(0),\mathrm{exp}(f)=1\left(fL_2(G)\right)$$
(65)
Using (62), (65), and (45) we get for $`j,k=1,\mathrm{},N`$
$`\alpha _{jk}=\left((e^{\frac{a^2}{2}}1)e^{\frac{a^2}{2}}\right)^{\frac{1}{2}}{\displaystyle \frac{\gamma }{\sqrt{N}}}{\displaystyle \underset{s=1}{\overset{N}{}}}[|\mathrm{exp}(\sqrt{2}aK_1g_s),\mathrm{exp}\left({\displaystyle \frac{a}{\sqrt{2}}}K_1(g_j+g_k)\right)`$
$`|\mathrm{exp}\left(\sqrt{2}aK_1g_s\right),\mathrm{exp}\left({\displaystyle \frac{a}{\sqrt{2}}}K_1\left(g_k\right)\right)]`$ (66)
We have
$$\mathrm{exp}(f_1),\mathrm{exp}(f_2)=e^{f_1,f_2}\left(f_1,f_2L_2(G)\right)$$
(67)
Using (13) and (67) we obtain
$`0=\mathrm{exp}\left(\sqrt{2}aK_1g_s\right),\mathrm{exp}\left({\displaystyle \frac{a}{\sqrt{2}}}K_1\left(g_j+g_k\right)\right)(sj)`$
$`\mathrm{exp}(\sqrt{2}aK_1,g_s,\mathrm{exp}\left({\displaystyle \frac{a}{\sqrt{2}}}K_1g_k\right))`$ (68)
From (61), (66), (67), and (68) it follows
$$\alpha _{jk}=\left(\left(e^{\frac{a^2}{2}}1\right)e^{\frac{a^2}{2}}e^{a^2}\right)^{\frac{1}{2}}\frac{\gamma }{\sqrt{n}}\left(e^{a^2K_1g_j,K_1(g_j+g_k)}e^{a^2K_1g_j,K_1g_k}\right)$$
(69)
Now (13) and (14) implies
$$K_1g_j,K_1g_k=\frac{1}{2}\delta _{jk}$$
(70)
For that reason (63), and (64) follow from (69). $`\mathrm{}`$
In the following we fix a pair $`n,m\{1,\mathrm{},N\}`$.
###### REMARK 4.2
Without loss of generality we can assume
$$B_n=\mathrm{𝟏}$$
(71)
That we can explain as follows:
Using (57,58), (59), and (54) we obtain in the case (71)
$`\stackrel{~}{\mathrm{\Theta }}_{km}(\rho )`$ $`=`$ $`\stackrel{~}{\mathrm{\Theta }}_{nm}\left(B_k^{}\rho B_k\right)(k=1,\mathrm{},N)`$
$`\stackrel{~}{p}_{km}(\rho )`$ $`=`$ $`\stackrel{~}{p}_{nm}\left(B_k^{}\rho B_k\right)(k=1,\mathrm{},N)`$
On the other hand from theorem 2.1 follows that in the case (71 ) for all states $`\rho `$ with (16) and (17) it holds
$$\mathrm{\Lambda }_{km}(\rho )=\mathrm{\Lambda }_{nm}\left(B_k^{}\rho B_k\right)(k=1,\mathrm{},N)$$
Finally it is easy to show that $`B_k^{}\rho B_k`$ fulfills (16), and (17) if the state $`\rho `$ fulfills (16) and (17).
For that reasons theorem 3.1 would be proved if we could prove (40), and (41) on the assumption that we have (71).
Now from (30), (49), and (37) we get
$$(U_m^{}\mathrm{\Gamma }(T))V(|\stackrel{~}{\eta }|\mathrm{exp}(0))=\frac{\gamma }{\sqrt{N}}\underset{k=1}{\overset{N}{}}|\mathrm{exp}\left(aK_1g_k\right)|\mathrm{exp}(aK_2g_{km})$$
(72)
###### LEMMA 4.3
Put for $`s=1,\mathrm{},N`$
$`\beta _s:=((|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|)V\mathrm{𝟏})(\mathrm{𝟏}U_m^{}\mathrm{\Gamma }(T))(\mathrm{𝟏}V)|\mathrm{\Psi }_s|\stackrel{~}{\eta }`$
$`|\mathrm{exp}(0)\mathrm{exp}(0)|(s=1,\mathrm{},N)`$ (73)
Then it holds
$`\beta _s={\displaystyle \frac{\gamma ^2}{N}}(1e^{\frac{a^2}{2}})^{\frac{1}{2}}|\mathrm{exp}(0)|\stackrel{~}{\eta }((1e^{\frac{a^2}{2}}){\displaystyle \underset{j=1}{\overset{N}{}}}c_{sj}|\mathrm{exp}\left(aK_2g_{jm}\right)`$
$`+e^{\frac{a^2}{2}}{\displaystyle \underset{j=1}{\overset{N}{}}}c_{sj}{\displaystyle \underset{k=1}{\overset{N}{}}}|\mathrm{exp}\left(aK_2g_k\right))`$ (74)
Proof: From (17), (72), and (73) we get
$`\beta _s={\displaystyle \underset{j=1}{\overset{N}{}}}c_{sj}{\displaystyle \frac{\gamma }{\sqrt{N}}}{\displaystyle \underset{k=1}{\overset{N}{}}}[(|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|)V(|\mathrm{exp}\left(aK_1g_j\right)\mathrm{exp}(0)`$
$`|\mathrm{exp}\left(aK_1g_k\right))]|\mathrm{exp}\left(aK_2g_{km}\right)`$ (75)
Further we have
$`(|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|)V(|\mathrm{exp}\left(aK_1g_j\right))\mathrm{exp}(0)|\mathrm{exp}\left(aK_1g_k\right)`$
$`=|\mathrm{exp}(0)|\stackrel{~}{\eta }|\mathrm{exp}(0)|\stackrel{~}{\eta },V\left(|\mathrm{exp}\left(aK_1g_j\right)\mathrm{exp}(0)|\mathrm{exp}\left(aK_1g_k\right)\right)`$
$`(j,k=1,\mathrm{},N)`$ (76)
Using Lemma 4.1, (75), and (76) we obtain
$`\beta _s={\displaystyle \frac{\gamma ^2}{N}}\left(1e^{\frac{a^2}{2}}\right)^{\frac{1}{2}}|\mathrm{exp}(0)|\stackrel{~}{\eta }{\displaystyle \underset{j=1}{\overset{N}{}}}c_{sj}|\mathrm{exp}\left(aK_2g_{jm}\right)`$
$`+{\displaystyle \frac{\gamma ^2}{N}}\left((1e^{\frac{a^2}{2}})e^{a^2}\right)^{\frac{1}{2}}|\mathrm{exp}(0)|\stackrel{~}{\eta }({\displaystyle \underset{j}{}}{\displaystyle \underset{kj}{}}c_{sj}|\mathrm{exp}\left(aK_2g_{km}\right))`$
That implies (74). $`\mathrm{}`$
Now we put
$$|\mathrm{\Psi }_0:=\frac{1}{\sqrt{N}}\underset{j=1}{\overset{N}{}}|\mathrm{exp}\left(aK_1g_j\right)\mathrm{exp}(0)$$
(77)
Since
$$F_+=\mathrm{𝟏}|\mathrm{exp}(0)\mathrm{exp}(0)|,$$
one easily checks
$`F_+|\mathrm{exp}\left(aK_rg_k\right)=\left(1e^{\frac{a^2}{2}}\right)^{\frac{1}{2}}|\mathrm{exp}\left(aK_rg_k\right)\mathrm{exp}(0)`$
$`(r=1,2;k=1,\mathrm{},m)`$ (78)
Using (77), and (78) we obtain
$`F_+\left({\displaystyle \underset{k=1}{\overset{N}{}}}|\mathrm{exp}\left(aK_2g_k\right)\right)`$ $`=`$ $`\left(1e^{\frac{a^2}{2}}\right)^{\frac{1}{2}}\sqrt{N}U_m\mathrm{\Gamma }(T)|\mathrm{\Psi }_0`$
$`=`$ $`\left(1e^{\frac{a^2}{2}}\right)^{\frac{1}{2}}\sqrt{N}\mathrm{\Gamma }(T)U_m|\mathrm{\Psi }_0`$
Using the same arguments we get
$`F_+\left({\displaystyle \underset{j=1}{\overset{N}{}}}c_{sj}|\mathrm{exp}\left(aK_2g_{jm}\right)\right)`$ $`=`$ $`\left(1e^{\frac{a^2}{2}}\right)^{\frac{1}{2}}U_m\mathrm{\Gamma }(T)|\mathrm{\Psi }_s(s=1,\mathrm{},N)`$ (80)
$`=`$ $`\left(1e^{\frac{a^2}{2}}\right)^{\frac{1}{2}}\mathrm{\Gamma }(T)U_m|\mathrm{\Psi }_s`$ (81)
Finally we have
$`\left(|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|F_+\right)(V\mathrm{𝟏})`$ (82)
$`=`$ $`\left(\mathrm{𝟏}\mathrm{𝟏}F_+\right)\left(|\mathrm{exp}(0)\mathrm{exp}(0)||\stackrel{~}{\eta }\stackrel{~}{\eta }|\right)V\mathrm{𝟏}`$
Using (54), (71), (4), (80), and Lemma 4.3 one easily checks the following equality.
$`W_{nm}\left(|\mathrm{\Psi }_s|\stackrel{~}{\eta }|\mathrm{exp}(0)\right)(s=1,\mathrm{},N)`$
$`=`$ $`{\displaystyle \frac{\gamma ^2}{N}}(1e^{\frac{a^2}{2}})(|\mathrm{exp}(0)|\stackrel{~}{\eta }\mathrm{\Gamma }(T)U_mB_n^{})`$
$`\left(\left(1e^{\frac{a^2}{2}}\right)|\mathrm{\Psi }_s+e^{\frac{a^2}{2}}\left({\displaystyle \underset{j}{}}c_{sj}\right)\sqrt{N}|\mathrm{\Psi }_0\right)`$
For that reason we have the following Lemma
###### LEMMA 4.4
For each bounded operator $`A`$ on $``$ and $`s=1,\mathrm{},N`$ it holds
$`\vartheta _s(A)`$ $`:`$ $`=W_{nm}(|\mathrm{\Psi }_s|\stackrel{~}{\eta }|\mathrm{exp}(0)),(\mathrm{𝟏}\mathrm{𝟏}A)W_{nm}(|\mathrm{\Psi }_s|\stackrel{~}{\eta }|\mathrm{exp}(0))`$
$`=`$ $`\left({\displaystyle \frac{\gamma ^2}{N}}\right)^2(1e^{\frac{a^2}{2}})^2[(1e^{\frac{a^2}{2}})^2\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_s,A\mathrm{\Gamma }(T)U_mB_N^{}|\mathrm{\Psi }_s`$
$`+e^{\frac{a^2}{2}}\left(1e^{\frac{a^2}{2}}\right)\left({\displaystyle \underset{j}{}}c_{sj}\right)\sqrt{N}\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_s,A\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_0`$
$`+e^{\frac{a^2}{2}}\left(1e^{\frac{a^2}{2}}\right)\left(\overline{{\displaystyle \underset{j}{}}c_{sj}}\right)\sqrt{N}\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_0,A\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_s`$
$`+e^{a^2}|{\displaystyle \underset{j}{}}c_{sj}|^2N\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_0,A\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_0`$
Now from (16) we get
$`\rho |\stackrel{~}{\eta }\stackrel{~}{\eta }||\mathrm{exp}(0)\mathrm{exp}(0)|`$ (84)
$`=`$ $`{\displaystyle \underset{s=1}{\overset{N}{}}}\lambda _s|\mathrm{\Psi }_s\stackrel{~}{\eta }\mathrm{exp}(0)\mathrm{\Psi }_s\stackrel{~}{\eta }\mathrm{exp}(0)|`$
On the other hand $`\left(|\mathrm{\Psi }_s\stackrel{~}{\eta }\mathrm{exp}(0)\right)_{s=1}^N`$ is an ONS because $`\left(\mathrm{\Psi }_s\right)_{s=1}^N`$ is an ONS. for that reason from (57,58), (84), and Lemma 4.4 with $`A=\mathrm{𝟏}`$ it follows
$`\stackrel{~}{p}_{nm}(\rho )=\left({\displaystyle \frac{\gamma ^2}{N}}\right)^2(1e^{\frac{a^2}{2}})^2[(1e^{\frac{a^2}{2}})^2+Ne^{a^2}{\displaystyle \underset{s=1}{\overset{N}{}}}\lambda _s|{\displaystyle \underset{s=1}{}}c_{sj}|^2`$
$`+\sqrt{N}e^{\frac{a^2}{2}}(1e^{\frac{a^2}{2}}){\displaystyle \underset{s=1}{\overset{N}{}}}\lambda _s{\displaystyle \underset{j=1}{\overset{N}{}}}(c_{sj}\mathrm{\Psi }_s,\mathrm{\Psi }_0+\overline{c_{sj}\mathrm{\Psi }_s,\mathrm{\Psi }_0})]`$ (85)
As $`\left(|\mathrm{exp}(aK_1g_j)\mathrm{exp}(0)\right)_{j=1}^N`$ is an ONS we can calculate easily
$$|\mathrm{\Psi }_s,|\mathrm{\Psi }_0=\frac{1}{\sqrt{N}}\underset{k=1}{\overset{N}{}}c_{sk}$$
For that reason from (85) follows
$`\stackrel{~}{p}_{nm}(\rho )=\left({\displaystyle \frac{\gamma ^2}{N}}\right)^2(1e^{\frac{a^2}{2}})^2[(1e^{\frac{a^2}{2}})^2+{\displaystyle \underset{s=1}{\overset{N}{}}}\lambda _s|{\displaystyle \underset{j=1}{\overset{N}{}}}c_{sj}|^2`$
$`(Ne^{a^2}+2\sqrt{N}e^{\frac{a^2}{2}}(1e^{\frac{a^2}{2}}))]`$ (86)
Further we have $`\underset{s}{}\lambda _s=1`$ and
$$\left|\underset{j=1}{\overset{N}{}}c_{sj}\right|^2\underset{j}{}\underset{k}{}\left|c_{sj}\right|\left|c_{sk}\right|\underset{j}{}\underset{k}{}\left(\frac{|c_{sj}|^2}{2}+\frac{|c_{sk}|^2}{2}\right)N$$
(87)
Using (86), (87) and the definition of $`\gamma `$ (cf. (29 )) we can estimate
$$\begin{array}{cc}& \left|\stackrel{~}{p}_{nm}(\rho )\frac{1}{N^2}\right|\hfill \\ & =\left(\frac{\gamma ^2}{N}\right)^2\left|\left(1e^{\frac{a^2}{2}}\right)^2\left[\left(1e^{\frac{a^2}{2}}\right)^2+\underset{s}{}\lambda _s\left|\underset{j}{}c_{sj}\right|^2\left(Ne^{a^2}+2\sqrt{N}e^{\frac{a^2}{2}}\left(1e^{\frac{a^2}{2}}\right)\right)\right]\frac{1}{\gamma ^4}\right|\hfill \\ & \frac{1}{N^2}|(1e^{\frac{a^2}{2}})^4+(1e^{\frac{a^2}{2}})^2\underset{s}{}\lambda _s|\underset{j}{}c_{sj}|^2(Ne^{a^2}+2\sqrt{N}e^{\frac{a^2}{2}}(1e^{\frac{a^2}{2}}))\hfill \\ & (1+(N1)e^{a^2})^2|\hfill \\ & \frac{1}{N^2}\left(\left|\left(1e^{\frac{a^2}{2}}\right)^4\left(1+(N1)e^{a^2}\right)^2\right|+e^{\frac{a^2}{2}}N\left(N+2\sqrt{N}\right)\right)\hfill \\ & \frac{1}{N^2}\left(e^{\frac{a^2}{2}}\left(14+N^2\right)+e^{\frac{a^2}{2}}N(N+2\sqrt{N})\right)\hfill \end{array}$$
That implies (41).
###### LEMMA 4.5
We use the notation $`\vartheta _s(A)`$ from Lemma 4.4. Then for each bounded operator $`A`$ on $``$ and $`s=1,\mathrm{},N`$ it holds
$`Z_s(A)`$ $`:`$ $`=\left|{\displaystyle \frac{\vartheta _s(A)}{\stackrel{~}{p}_{nm}(\rho )}}\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_s,A\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_s\right|`$
$``$ $`{\displaystyle \frac{2e^{\frac{a^2}{2}}}{\left(1e^{\frac{a^2}{2}}\right)^2}}\left(N^2+N\sqrt{N}+N\right)`$
Proof: Using Lemma 4.4 and the estimation
$$|\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_k,A\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_r|A(k,r=0,\mathrm{},N)$$
We get
$$\begin{array}{cc}\hfill Z_s(A)& A\left|\left(\frac{\gamma ^2}{N}\right)^2\left(1e^{\frac{a^2}{2}}\right)^2\left(\stackrel{~}{p}_{nm}(\rho )\right)^11\right|\hfill \\ & +A\left(\frac{\gamma ^2}{N}\right)^2(1e^{\frac{a^2}{2}})^4\left(\stackrel{~}{p}_{nm}(\rho )\right)^1\left[2e^{\frac{a^2}{2}}(1e^{\frac{a^2}{2}})\sqrt{N}\right|\underset{j}{}c_{sj}|\hfill \\ & +e^{a^2}N|\underset{j}{}c_{sj}|^2]\hfill \end{array}$$
Because of (86) it follows
$$\begin{array}{cc}\hfill Z& A\left[\right|\frac{\left(1e^{\frac{a^2}{2}}\right)^2}{\left(1e^{\frac{a^2}{2}}\right)^2+\underset{s}{}\lambda _s|\underset{j}{}c_{sj}|^2\left(Ne^{a^2}+2\sqrt{N}e^{\frac{a^2}{2}}\left(1e^{\frac{a^2}{2}}\right)\right)}1|\hfill \\ & +\frac{2e^{\frac{a^2}{2}}\left(1e^{\frac{a^2}{2}}\right)\sqrt{N}|\underset{j}{}c_{sj}|+e^{a^2}N|\underset{j}{}c_{sj}|^2}{\left(1e^{\frac{a^2}{2}}\right)^2+\underset{s}{}\lambda _s|\underset{j}{}c_{sj}|^2\left(Ne^{a^2}+2\sqrt{N}e^{\frac{a^2}{2}}\left(1e^{\frac{a^2}{2}}\right)\right)}]\hfill \end{array}$$
Using (87) we get
$$\begin{array}{cc}& \left|\frac{\left(1e^{\frac{a^2}{2}}\right)^2}{\left(1e^{\frac{a^2}{2}}\right)^2+\underset{s}{}\lambda _s|\underset{j}{}c_{sj}|^2\left(Ne^{a^2}+2\sqrt{N}e^{\frac{a^2}{2}}\left(1e^{\frac{a^2}{2}}\right)\right)}1\right|\hfill \\ & \frac{e^{\frac{a^2}{2}}}{\left(1e^{\frac{a^2}{2}}\right)^2}\left(N^2+2N\sqrt{N}\right)\hfill \end{array}$$
and
$$\begin{array}{cc}& \frac{2e^{\frac{a^2}{2}}\left(1e^{\frac{a^2}{2}}\right)\sqrt{N}|\underset{j}{}c_{sj}|+e^{a^2}N|\underset{j}{}c_{sj}|^2}{\left(1e^{\frac{a^2}{2}}\right)^2+\underset{s}{}\lambda _s|\underset{j}{}c_{sj}|^2\left(Ne^{a^2}+2\sqrt{N}e^{\frac{a^2}{2}}\left(1e^{\frac{a^2}{2}}\right)\right)}\hfill \\ & \frac{e^{\frac{a^2}{2}}}{\left(1e^{\frac{a^2}{2}}\right)^2}\left(2N+N^2\right)\hfill \end{array}$$
That proves Lemma 4.5. $`\mathrm{}`$
We have the representation (84) of $`\rho |\stackrel{~}{\eta }\stackrel{~}{\eta }||\mathrm{exp}(0)\mathrm{exp}(0)|`$ as a mixture of orthogonal projections. Thus from (56) and (57,58) we get with the notation $`\vartheta _s(A)`$ from Lemma 4.4
$$\mathrm{tr}\left(\stackrel{~}{\mathrm{\Theta }}_{nm}(\rho )A\right)=\underset{s}{}\lambda _s\vartheta _s(A)\left(\stackrel{~}{p}_{nm}(\rho )\right)^1$$
On the other hand from Theorem 2.1 follows
$$\mathrm{tr}\left(\mathrm{\Lambda }_{nm}(\rho )A\right)=\underset{s}{}\lambda _s\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_s,A\mathrm{\Gamma }(T)U_mB_n^{}|\mathrm{\Psi }_s$$
Consequently we have with notation $`Z_s(A)`$ from the Lemma 4.5
$$|\mathrm{tr}\left(\stackrel{~}{\mathrm{\Theta }}_{nm}(\rho )A\right)\mathrm{tr}\left(\mathrm{\Lambda }_{nm}(\rho )A\right)|\underset{s}{}\lambda _sZ_s(A)$$
For that reason (40) follows from Lemma 4.5, and $`\underset{s}{}\lambda _s=1`$.
That completes the proof of Theorem 3.1.
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# Slow holes in the triangular Ising antiferromagnet
## I Introduction
The problem of mobile holes in a magnetic background is an especially interesting example of charge motion in a strongly correlated electron system. In the cases that we have in mind, the magnetism involves local moments which reflect an underlying strong electron-electron interaction that can produce new and interesting effects when doping creates charge carriers. Much of the recent work on this problem has been inspired by the cuprate superconductors, whose physics is widely believed to be centrally connected to their genesis as doped Mott insulators. In this context, the central observation is that hole motion in an ordered antiferromagnetic background is frustrated, leading to the expectation that doping will result in a rearrangement of the magnetic backgound in a manner conducive to relieving this frustration. Suggestions for what this entails include RVB (resonating valence bond) theory, stripe formation, and, in models with purely short-ranged interactions, phase separation. A second theme in this setting is the large spin degeneracy of the extreme Mott insulator (e.g. the $`U=\mathrm{}`$ Hubbard model) at half filling and the role of doping in lifting it. Perhaps the most celebrated result along these lines is the Nagaoka theorem, which established that a single hole would completely polarize the spin background. While the situation at finite dopings remains unsettled, the Nagaoka result does show that doping can lift the degeneracy in striking ways.
In this paper we report some results on dilute holes introduced into a frustrated magnetic system—the particular system studied is the canonical example of this class, the triangular lattice Ising antiferromagnet first studied by Wannier and Houtappel. This system realises aspects of both themes touched on above. It has local antiferromagnetic order which leads to frustration of the hole motion. In addition however, as this is an Ising system, the magnetic frustration leads to a finite zero-point entropy per site, $`𝒮`$. This feature is reminiscent of the Mott insulator cited above, but we should note that in this case the degeneracy is generated as a cooperative effect: the Ising model on the triangular lattice has a large number of ground states, $`𝒩_{GS}`$, because of the geometrically frustrated nature of its magnetic interactions. In each triangle, at least one pair of spins has to be aligned, and hence one bond frustrated. Any state with exactly one bond per triangle is hence a ground state, and $`𝒩_{GS}`$ is found to scale exponentially with the number of spins, $`N`$, with $`𝒮=k_B\mathrm{ln}(𝒩_{GS})/N0.323k_B`$.
In the following, we present a set of results on the question of if, and how, this degeneracy is lifted upon dilute doping. To make progress, we will assume the hole motion is slow—their kinetic energy being taken to be much smaller than the magnetic exchange. While Ising magnets exist, slow holes might be harder to find. Our interest in this limit then is that a) it poses the question of how doping interacts with a frustrated magnetic background most cleanly, which is of theoretical interest, and b) that the ordering patterns we find in this limit could well persist when more enterprising holes are considered. While more work is needed to investigate the validity of (b) we should note concerning (a) that the response of highly frustrated magnets to perturbations is more generally interesting. On account of the large degeneracy, these systems are unstable in a large number of directions, promising a wide range of unexpected and unusual physical phenomena. Along these lines, in a recent study along with P. Chandra, we have explored the phase structure produced by switching on the quantum dynamics of a transverse magnetic field instead.
We now turn to the Hamiltonian, $`H`$, we study. It is
$`H`$ $`=`$ $`H_t+H_J+H_\eta `$ (1)
$`=`$ $`t{\displaystyle \underset{ij,\sigma }{}}P\left(c_{i\sigma }^{}c_{j\sigma }+c_{j\sigma }^{}c_{i\sigma }\right)P+J^z{\displaystyle \underset{ij}{}}S_i^zS_j^z`$ (3)
$`{\displaystyle \frac{\eta J^z}{4}}{\displaystyle \underset{ij}{}}n_in_j,`$
where $`J^z>0`$ is the strength of the antiferromagnetic Ising exchange, $`t`$ is the hopping integral, and $`c_{i\sigma }(c_{i\sigma }^{})`$ annihilates (creates) an electron at site $`i`$ with spin $`\sigma `$. $`S_i^z=\frac{1}{2}(c_i^{}c_ic_i^{}c_i)`$ is the $`z`$-component of the spin-1/2 operator at site $`i`$, $`n_i=_\sigma c_{i\sigma }^{}c_{i\sigma }`$ is the total density operator and $`P`$ is a projector that ensures that all sites are at most singly occupied. Finally, the sum $`ij`$ runs over nearest neighbor bonds and the sum $`\sigma `$ over up and down spin orientations. As noted earlier, we study the range of parameters $`tJ^z`$, where the hopping acts as a perturbation to the exchange coupling. We consider strengths of the nearest neighbour density-density coupling of $`0\eta 1`$. In the derivation of the rotationally invariant $`tJ`$ model for $`S=1/2`$ from the Hubbard model, such a coupling is generated with $`\eta =1`$. For Ising systems, possibly with spins greater than $`1/2`$, there is no such precise fixing of its magnitude and for that reason, as well as for more general theoretical interest, we consider a range of values for it.
In the rest of the paper, we first prove a “frustrated Nagaoka theorem”, namely that the introduction of a single hole into the ground state manifold causes the system to order into a state that is spin polarized and breaks translational symmetry with a $`\sqrt{3}\times \sqrt{3}`$-unit cell. In this state, the frustrated bonds form a hexagonal (honeycomb) lattice on which the hole hops. Next we consider generalizing this physics to a finite but dilute density of holes, which turns out to depend strongly on $`\eta `$. For $`\eta >1/3`$, phase separation obtains. The case $`\eta <1/3`$ is much more interesting. Here, we find a spinless Fermi liquid state that lives on the “hexagonal backbone” of the single hole problem. Further, the surrounding sites mediate an attractive interaction that implies that the Fermi liquid is unstable to a superconducting state in a non-zero odd angular momentum channel. The resulting state then breaks three symmetries of the parent insulating state at once: Ising, translation and the charge $`U(1)`$. Finally, we turn to non-zero temperatures $`J^zTt`$, and show using Grassman functional integral techniques, that there is an algebraically decaying, entropic interaction between two holes which is oscillatory as a function of distance and which indicates a tendency for the holes to segregate on the hexagonal backbone.
## II A frustrated Nagaoka theorem
We show in this section that the ground state of our Hamiltonian (Eq. 3) with a single hole present is macroscopically ordered. We will begin by establishing this in the limit where we take $`J^z=\mathrm{}`$, i.e. in the degenerate perturbation theory problem among a set of single hole states that we identify below. We will comment on why the conclusion is unaltered when $`J^zt`$ but not infinite, at the end of the proof of this first part.
The proof will consist of first establishing a lower bound for the energy of a one-hole state on arbitrary finite lattices and then demonstrating that this can be saturated only on a subset of finite lattices and that on the latter there is, up to global symmetry operations, just one state that does so. For specificity, we will assume that our lattices come with periodic boundary conditions.
We first identify the one-hole states that minimise the exchange energy, noting that the term $`H_\eta `$ in the Hamiltonian (Eq. 3) is the same for any one-hole state. Starting from any ground state, we can only remove an electron which experiences zero net exchange field, i.e., which has three frustrated and three satisfied bonds. There are exponentially many such one-hole states.
The hopping Hamiltonian connects those one-hole states which only differ by the exchange of a hole and a neighbouring spin. Fig. 1a shows a spin configuration which allows the hole to hop to three of its neighbouring sites. Up to symmetry operations, this is the only such (local) spin configuration. Since the hopping integral is $`t`$, this means that the energy which can be gained from $`H_t`$ is at most $`3t`$. We note that this bound is quite general and holds for all finite lattices with arbitrary boundary conditions. Indeed, starting with any given one-hole state, we can think of the all the states reached by hops as inducing a graph, whose topology is in general far more complicated than that of the parent lattice. The statement is that regardless of the complexity, the graph has a maximum local coordination number of 3 and that the nearest neighbor hopping problem on it has at best an energy of $`3t`$.
This energy gain can only be realised if any state reached after one hop again permits hoppping in three directions. Demanding that this be the case completely constrains the parent spin configuration to the one depicted in Fig. 1b; the other five states that satisfy this condition are obtained using translations or Ising reversal. In this configuration, the hole can occupy any site with a frustrated bond. With periodic boundary conditions, it is clear that not all finite lattices accomodate this pattern, but that there is a subset of arbitrarily large size which does—these are the lattices with linear sizes that are integer multiples of 3. In the following we construct our desired one-hole eigenstate $`|h`$ with energy $`3t`$ on this subset.
Define $`|G=_ic_{i,\sigma (i)}^{}|0`$, where $`\sigma (i)`$ is the spin at site $`i`$ in the hexagonal state depicted in Fig. 1b, and $`|0`$ is the empty lattice; the ordering of the operators is immaterial in what follows. Next, define the state with a hole on site $`n`$ as $`|n=(1)^Sc_{n,}|G`$. Here, $`S(n)=0(1)`$ if the hole is created on sublattice A (B) of the bipartite hexagonal lattice defined by the frustrated bonds. The matrix elements $`m|H_t|n`$ vanish unless $`m`$ and $`n`$ refer to neighbouring sites on the hexagonal lattice, in which case
$`m|H_t|n`$ $`=`$ $`m\left|t{\displaystyle \underset{ij,\sigma }{}}c_{i\sigma }^{}c_{j\sigma }\right|n`$ (4)
$`=`$ $`tm\left|c_n^{}c_m\right|n`$ (5)
$`=`$ $`t(1)^{(0+1)}G\left|c_m^{}c_n^{}c_mc_n\right|G`$ (6)
$`=`$ $`tG\left|c_n^{}c_nc_m^{}c_m\right|G=t.`$ (7)
Now, let $`|h=_{i=1}^{2N/3}|n`$. For this state, we have $`H_t|h=3t|h`$ since $`H_t`$ connects, with amplitude $`t`$, each single-hole component of the wavefunction to the three components with the hole on a neighbouring site. Note that this is, in essence, the ground state of a spinless hole hopping on a hexagonal lattice.
Unlike the ensemble of undiluted ground states, which has neither long-range order nor a net magnetisation, $`|h`$ incorporates long-range spin order. There is a three sublattice ($`\sqrt{3}\times \sqrt{3}`$) ordering pattern with the spin on two sublattices pointing up and down on the third. The magnetisation induced by adding the single hole is therefore $`N/3`$. It is interesting to note that this ordering pattern would also have been generated by an explicitly symmetry-breaking magnetic field pointing along the Ising axis and that the $`\sqrt{3}\times \sqrt{3}`$ structure is also picked out by a transverse field.
Finally, we ask whether the conclusion at $`J^z=\mathrm{}`$ is altered when $`J^z`$ is made finite. While we do not have a proof of this, we believe that this does not happen. Conceptually, it is useful to imagine generating an effective Hamiltonian in the space of all one-hole states by including fluctuations outside this space perturbatively in powers of $`t/J^z`$. Such an expansion will not invalidate the conclusion reached inside the hexagonal backbone states—it will only modify the state $`|h`$ perturbatively, preserving its symmetry characteristics at small $`t/J^z`$. The real danger is that a state in a different sector that has an energy $`ϵ=3t+\theta (L)`$ for a system of linear dimension $`L`$, when $`J^z=\mathrm{}`$, with $`lim_L\mathrm{}\theta (L)=0^+`$, will acquire a perturbative gain of $`O\left([t/J^z]^n\right)`$ for some $`n`$, in excess of the state $`|h`$, thereby invalidating our conclusion at sufficiently large $`L`$.
To gain insight into why this is unlikely to happen, consider the energetics at $`O\left([t/J^z]^2\right)`$. At this order, the hole can hop to a site where the net exchange field does not vanish and, depending on the spin configuration, either be forced to return to its starting site or hop to a different one. In terms of the induced graph mentioned above, the former corresponds to an on-site potential energy term, while the latter generates a new, or modifies the strength of an existing, edge (hopping amplitude). In the hexagonal backbone states, only the first option is available, yielding an effective potential energy of $`3t^2/J^z`$ at each site. If we consider altering the state near a particular site to generate more favourable potential or hopping matrix elements to $`O\left([t/J^z]^2\right)`$, this invariably requires sacrificing a hopping element of $`O(t)`$. We believe that this points to a certain rigidity of the hexagonal backbone which will operate to exclude similar subtleties at higher orders as well. Needless to say, a proof of this assertion would be desirable.
## III Small density of holes and leading $`t/J^z`$ corrections
We begin with the case $`\eta =0`$. First, we need to find the allowed many hole states at low dilution (and when $`t/J^z0`$). As before, we need to remove those electrons experiencing zero net exchange field since this ensures that the leading term of the Hamiltonian, $`H_J`$, continues to be optimised. This implies that one cannot remove two neighouring electrons: this would leave a triangle with only one occupied site and hence no contribution to the magnetic energy, whereas the optimal states are those in which the three bonds of each triangle have total magnetic energy $`J/4`$.
We showed in the last section that the kinetic energy selects a subset of the one-hole states in which the hole is maximally delocalized. We believe that at sufficiently low dopings, this logic continues to operate and the low-lying eigenstates of $`H`$ will be constructed out of the many-hole states obtained by removing electrons from the hexagonal backbone of the state depicted in Fig 1b. (However, we are no longer able to prove this.) This state continues to be favourable as it allows the holes to be maximally mobile by having coordination three everywhere. The only restriction being placed on the holes is not to sit immediately next to each other, which is a demand present for all possible background states and one that should not be too onerous at low dopings.
With these restrictions, the problem reduces to that of spinless fermions hopping on a hexagonal lattice with amplitude $`t`$ and with a nearest neighbor hard core repulsion. At low densities, this is a well studied problem and known to be a Fermi liquid. In this fashion, we find that the ground state of the dilute hole problem is magnetized, breaks translational symmetry and is a two dimensional Fermi liquid.
Now this is not quite true, for the Fermi liquid suffers from the well-known Kohn-Luttinger superconducting instability and so strictly at $`T=0`$ the system will break the charge $`U(1)`$ symmetry as well—thereby breaking three symmetries of the parent insulating state at once! While this happens as $`t/J^z0`$, we show next that for $`t/J^z`$ large but not infinite, the magnetic background of the hexagonal conducting lattice mediates an attractive interaction that also drives a superconducting instability—at low densities this is a much larger effect than the purely “internal” Kohn-Luttinger effect.
To this end we consider the effective Hamiltonian, $`H_{eff}`$, between holes on the hexagonal backbone to next-leading order in $`t/J^z`$, induced by excursions to the surrounding sites in which the system leaves the ground state manifold of $`H_J`$. The matrix elements of this effective Hamiltonian are determined by perturbation theory as follows. First, we define a $`k`$-hole state with holes at sites $`m_1,m_2,\mathrm{}m_k`$ (with some ordering of the sites, and for concreteness with $`m_i<m_j`$ for $`i<j`$) as
$`|m_1m_2\mathrm{}m_n={\displaystyle \underset{i=1}{\overset{k}{}}}c_i|G.`$ (8)
In the two-hole sector ($`k=2`$) we thus obtain:
$`m_1m_2\left|H_{eff}\right|m_1m_2={\displaystyle \underset{p_1,p_2}{}}{\displaystyle \frac{\left|p_1p_2\left|H_t\right|m_1m_2\right|^2}{E_{m_1,m_2}^0E_{p_1,p_2}^0}};`$ (9)
here $`E^0`$ are the energies of the two-hole states for $`t=0`$. Note that the overall sign of the many-hole states is immaterial at this order as only the squares of matrix elements of $`H_t`$ are used in determining $`H_{eff}`$.
To this order, more general, non-diagonal matrix elements are absent since holes hopping of the backbone have to hop back to their original site in order to reconstitute a leading-order ground state. When the two holes are not on the same hexagon, $`H_{eff}`$ always takes the same value which is just twice the contribution of an isolated hole and hence has the interpretation of an effective one-body potential as discussed in the last section. Subtracting this leaves a term that depends on the joint presence of the two holes and can thus be viewed as an effective two-body potential, $`V(R)`$, for the holes. This has the form, restoring $`\eta `$ for future use:
$`V(R=1)`$ $`=`$ $`\mathrm{}`$ (10)
$`V(R=\sqrt{3})`$ $`=`$ $`{\displaystyle \frac{t^2}{J}}{\displaystyle \frac{2(134\eta \eta ^2)}{34\eta +\eta ^2}}`$ (11)
$`V(R=2)`$ $`=`$ $`{\displaystyle \frac{t^2}{J}}{\displaystyle \frac{2(1+\eta )}{3\eta }}`$ (12)
$`V(R>2)`$ $`=`$ $`0,`$ (13)
where $`R`$ is measured in lattice constants of the triangular lattice.
This effective interaction is dominantly repulsive on account of the hard core piece considered earlier. Hence the additional, weak, attraction can only induce a superconducting instability of the parent Fermi liquid and we do not need to worry about the possiblity of phase separation at $`\eta =0`$. In the continuum, it is easy to see that the two body problem with a hard core and a weak attractive tail posesses attractive phase shifts in sufficiently high angular momentum channels—at low densities this is sufficient to lead to pairing. While we have not carried out a detailed analysis on the hexagonal lattice, the general conclusion will hold.
We end this section by considering the effect of switching on an $`\eta >0`$. For $`\eta <1/3`$, the above arguments retain their validity, with the main consequence of $`\eta `$ being a strengthening of the attractive part of the effective interaction. For $`\eta >1/3`$, the allowed many-hole states change entirely in character. The density-density attraction $`H_\eta `$ then is strong enough to overcome the exchange term $`H_J`$ and produce phase separation. This happens because phase-separation costs an energy of $`J(3\eta +1)/4`$ per hole, whereas the optimal states described above achieve a cost of $`3J\eta /2`$.
## IV Effective interaction for classical holes
Thus far we have focussed on the physics precisely at $`T=0`$. In this limit the degeneracy of the parent magnet is “selected away”. At non-zero temperatures, but still below $`J^z`$, the ground-state entropy of the parent magnet will again play a role. As a first step towards an understanding of finite temperatures we consider $`tTJ`$, equivalently, we discuss the behaviour of classical annealed holes on the triangular antiferromagnet. Our central result in this part is an effective interaction, induced entropically, between two holes. We close with some remarks on the finite density problem.
### A Two holes
Two holes on the triangular lattice experience an entropic interaction because the number of ground-state spin configurations, $`Z(r)`$, minimising $`H_J`$ depends on their separation vector, $`r`$. From a knowledge of $`Z(r)`$, upon subtracting the “creation entropy” of the two holes, $`\mathrm{ln}Z(\mathrm{})`$, we can obtain an effective interaction potential $`\beta v(r)`$, which vanishes as $`r\mathrm{}`$:
$`\beta v(r)`$ $`=`$ $`\mathrm{ln}Z(r)+\mathrm{ln}Z(\mathrm{})`$ (14)
$`=`$ $`\mathrm{ln}(Z(r)/Z)+\mathrm{ln}(Z(\mathrm{})/Z);`$ (15)
here $`Z`$ is the partition function of the undiluted system.
Defining $`Z(\mathrm{})/Z\mathrm{{\rm Y}}^2`$ and $`Z(r)/Z\mathrm{{\rm Y}}^2\zeta (r)`$, we obtain
$`\beta v(r)=\mathrm{ln}\left[1{\displaystyle \frac{\zeta (r)}{\mathrm{{\rm Y}}^2}}\right]{\displaystyle \frac{\zeta (r)}{\mathrm{{\rm Y}}^2}},`$ (16)
the last step being justified at large $`|r|`$ by the smallness of $`\zeta (r)/\mathrm{{\rm Y}}^2=0`$, as we will show below.
To determine $`\zeta (r)`$ and $`\mathrm{{\rm Y}}`$, it is convenient to use a standard representation of the spin problem as a classical dimer model on the dual hexagonal lattice. The presence of a hole is encoded by a certain dimer configuration around it. The latter is soluble by the Pfaffian techniques introduced by Kasteleyn although we will find it convenient to use the language of Grassmanian functional integrals introduced into such problems by Samuel. Thus, the function $`\zeta (r)`$ is obtained by evaluating a twelve-fermion correlation function.
In detail, this calculation proceeds as follows. First, we map ground state spin configurations on the triangular lattice to dimer configurations on its dual hexagonal lattice by drawing a dimer through each frustrated bond, placing its endpoints at the centres of the two triangles sharing the frustrated bond. Noting that each triangle has one and only one frustrated bond in the ground state, the dimers thus provide a hard-core covering of the lattice dual to the triangular lattice, namely the hexagonal lattice. Conversely, up to the global Ising symmetry, each dimer covering corresponds to a ground state of the Ising model.
The statistical mechanics of classical hard core dimer coverings on the hexagonal lattice has been studied by Kastelyn and some correlation functions have been computed by Yokoi, Nagle and Salinas.
We extend this method to include the presence of diluted sites on the original triangular lattice. As outlined above, a hole can only occupy the site of a spin experiencing a net exchange field of zero. In dimer language, this means that only a spin at the centre of a hexagonal plaquette occupied by exactly three dimers can be replaced by a hole, as depicted in Fig. 2. To determine $`Z(r)`$, we are thus only allowed to count those dimer configurations in which such plaquettes are encountered at the location of the holes.
To calculate $`Z`$ in the first instance, one defines Grassman variables, $`\psi _i`$, on each site of the $`s`$ sites on the hexagonal lattice. One then writes
$`Z={\displaystyle \left(\underset{l=s}{\overset{1}{}}d\psi _l\right)\mathrm{exp}(\psi _iA_{ij}\psi _j)}=\mathrm{Pf}(A)=\sqrt{detA},`$ (17)
where Pf$`(A)`$ is the Pfaffian of the (antisymmetric) matrix $`A`$.
The problem of determining $`A`$ for a general class of two-dimensional lattices has been solved by Kasteleyn. For the hexagonal lattice, which we are interested in, it turns out that one has to double the unit cell of the lattice to contain 4 spins, with the unit cell labelled as in Fig. 2, so that $`A_{ij}(r)`$ equals
$`\left(\begin{array}{cccc}0& \delta _{r,0}\delta _{r,\xi }& \delta _{r,0}& 0\\ \delta _{r,0}+\delta _{r,\xi }& 0& 0& \delta _{r,\upsilon }\\ \delta _{r,0}& 0& 0& \delta _{r,0}\delta _{r,\xi }\\ 0& \delta _{r,\upsilon }& \delta _{r,0}+\delta _{r,\xi }& 0\end{array}\right).`$ (18)
Here, the vectors $`\xi `$ and $`\upsilon `$ denote the translation vectors of the lattice with the doubled unit cell (Fig. 2). The holes are located in cells $`r_1`$, site $`i`$, and $`r_2`$, site $`j`$, with $`r=r_2r_1`$.
This choice of $`A`$ accomplishes the following. If the exponential in the integrand for the partition function is expanded, the only terms that result in a nonzero contribution to the integral are those in which each $`\psi _i`$ appears exactly once. Each such term represents a covering of dimers with sites $`i`$ and $`j`$ being connected by a dimer if $`A_{ij}`$ is one of the prefactors of the term. The signs of the entries of $`A`$ are chosen such that every such term in fact integrates to 1 so that all dimer coverings are allocated equal weight. One thus obtains $`𝒮/k_B=\mathrm{ln}Z/N0.323`$.
We first caclulate the probability that a spin can be replaced with a hole, which we require as this will turn out to be $`\mathrm{{\rm Y}}`$, defined above. To ensure the plaquette surrounding a spin (the sites of which we label $`1\mathrm{}6`$) is covered by three dimers, we insert a prefactor $`A_{12}A_{34}A_{56}\psi _1\psi _2\psi _3\psi _4\psi _5\psi _6`$ into the integrand, which forces the three dimers into the required position. We thus need to calculate
$`\mathrm{{\rm Y}}`$ $`=`$ $`{\displaystyle \frac{1}{Z}}{\displaystyle }\left({\displaystyle \underset{l=s}{\overset{1}{}}}d\psi _l\right)\times `$ (20)
$`A_{12}A_{34}A_{56}\psi _1\psi _2\psi _3\psi _4\psi _5\psi _6\mathrm{exp}(\psi _iA_{ij}\psi _j).`$
Since the Fermionic action is quadratic, this becomes by Wick’s theorem the sum over all $`6!`$ contractions of pairs of $`\psi `$s:
$`\mathrm{{\rm Y}}`$ $`=`$ $`{\displaystyle \underset{P}{}}\mathrm{sign}(P)G_{i_1i_2}G_{i_3i_4}G_{i_5i_6},`$ (21)
where the sum runs over all permutations $`P=\{i_l|l=1\mathrm{}6\}`$ of $`\{1\mathrm{}6\}`$. These Green functions have been worked out by Yokoi et al. together with identities expressing all $`G_{ij}`$ in terms of $`G_{24}`$, which for $`y0`$ is given by
$`G_{24}(x,y)`$ $`=`$ $`{\displaystyle \frac{2}{\pi }}(1)^x{\displaystyle _0^{\pi /3}}𝑑\varphi {\displaystyle \frac{\mathrm{cos}(2x\varphi )}{\left|1+\mathrm{exp}(2i\varphi )\right|^{2(y+1)}}},`$ (22)
with $`G_{24}`$ given by the negative of this expression for $`y<0`$.
In fact, the sum for $`\mathrm{{\rm Y}}`$ only contains $`3!`$ nonzero terms as Green functions arising from contractions of sites belonging to the same sublattice of the hexagonal lattice vanish. We thus obtain
$`\mathrm{{\rm Y}}={\displaystyle \frac{2}{27}}{\displaystyle \frac{3\sqrt{3}}{8\pi ^3}}+{\displaystyle \frac{1}{2\sqrt{3}\pi }}0.14501.`$ (23)
To confirm the calculation so far, we remark that Monte Carlo simulations we have carried out in a different context (frustrated transverse field Ising models) give $`\mathrm{{\rm Y}}=0.147\pm 0.003`$.
The evaluation of $`Z(r)`$ thus requires twelve-fermion correlation functions as two holes have to be introduced. There are two technical difficulties in this calculation. Firstly, the distance-dependence of the Green functions is not known in closed form, and secondly, we now require $`6!=720`$ terms each containing at least six Green functions even after using the vanishing of same-sublattice Green functions.
As the reader can see, this calculation is straightforward in principle but somewhat involved in practice—it will turn out though that the final answer has a transparent rationale. Before proceeding further it is useful to discuss the same general phenonmenon in a toy model which is exactly solvable and where the combinatorics is more transparent. Consider a one-dimensional dimer model in which the dimers cover a two-leg ladder (see Fig 3). The entropy per rung of this dimer model is the golden mean, $`G=(\sqrt{5}+1)/21.62`$. The dimer model on this ladder arises in the study of a fully frustrated three-leg Ising ladder, and removing a spin from this Ising model corresponds to forcing the square surrounding this spin to be covered by two parallel dimers.
If two holes are adjacent to one another (Fig. 3b, top), there are no rungs separating the plaquettes surrounding them. Moving the hole on the right two units further away (Fig. 3b, bottom) costs an entropy of $`\mathrm{ln}G^2\mathrm{ln}2.62`$ in the semi-infinite region outside to the right of the hole pair and only gains an entropy of $`\mathrm{ln}2`$ as the dimer pair between the holes can only have two states. For even separations, it is thus advantageous for the holes to move close together.
Separating the holes by one unit costs entropy $`\mathrm{ln}G`$ outside while gaining nothing inside as a single dimer can only exist in the vertical state (Fig. 3b, middle). Moving the hole further away by two units, however, yields a total gain of $`\mathrm{ln}3\mathrm{ln}G^2>0`$ in entropy as the three dimers between the hole plaquettes can have three states. Therefore, for odd separations, the holes want to move further apart. In fact, the alternating two-hole potential for this model can be calculated both by Pfaffian and transfer matrix methods, and takes the form
$`\beta v_{\mathrm{ladder}}(r){\displaystyle \frac{\zeta _{\mathrm{ladder}}(r)}{\mathrm{{\rm Y}}_{\mathrm{ladder}}^2}}={\displaystyle \frac{4}{5}}{\displaystyle \frac{G1}{G^3\mathrm{{\rm Y}}_{\mathrm{ladder}}^2}}(1)^{r+1}G^{2r},`$ (24)
with $`\mathrm{{\rm Y}}_{\mathrm{ladder}}=1/\left(G\left(2G1\right)\right)`$.
Returning to the triangular lattice problem, the large $`|r|`$ asymptotics of the Green functions are calculated in Ref. :
$`G_{24}(x,y)\left(1\right)^x{\displaystyle \frac{\sqrt{3}y\mathrm{cos}(\frac{2\pi x}{3})+x\mathrm{sin}(\frac{2\pi x}{3})}{\pi x^2+3\pi y^2}}.`$ (25)
As these Green functions decay algebraically with $`|r|`$, the terms can be arranged according to how many contractions between sites belonging to different holes they contain. The leading terms are the ones in which all contractions are between sites of one or the other hole, and these terms are clearly independent of $`r`$. In fact, it is easy to see that the contractions around a given hole are simply those which yielded $`\mathrm{{\rm Y}}`$ in the single-hole calculation, and thus the leading, $`r`$-independent term is $`Z(\mathrm{})/Z=\mathrm{{\rm Y}}^2`$, as promised above.
The leading non-trivial $`r`$-dependence arises from terms which contain two Green functions connecting sites around the two holes, an odd number of connections being impossible. There are 324 such terms. With the assistance of Mathematica, we were able to analyze their asymptotics and obtain the exact long distance form,
$`\zeta (x,y)\mathrm{\Xi }{\displaystyle \frac{\mathrm{cos}(4\pi x/3)}{x^2+3y^2}},`$ (26)
where
$`\mathrm{\Xi }`$ $`=`$ $`{\displaystyle \frac{729324\sqrt{3}\pi 324\pi ^2+96\sqrt{3}\pi ^3+64\pi ^4}{288\pi ^6}}`$ (27)
$``$ $`0.025853.`$ (28)
Translating the co-ordinates back to the original triangular lattice, our final asymptotic expression for the interaction potential is
$`\beta v(r){\displaystyle \frac{\mathrm{\Xi }}{\mathrm{{\rm Y}}^2}}{\displaystyle \frac{\mathrm{cos}(4\pi x/3)}{x^2+y^2}}.`$ (29)
The physics of this expression is quite transparent. The wavevector $`𝐪=\frac{4\pi }{3}\widehat{x}`$ on the triangular lattice is characteristic of the hexagonal frustrated bond pattern sketched in Fig 1b. It has long been recognized (already by Wannier ) that this pattern represents the dominant ordering tendency in the ground state average of the classical problem. In a more modern formulation in terms of height representations the pattern is the reference flat surface about which the system exhibits critical fluctuations. In accordance with this we find from Eq. 29 that the holes attract when they share a hexagonal backbone of frustrated bonds, and repel when they do not. The algebraic decay could have been anticipated in the dimer language from our observation that the leading piece at large distance comes from a sum of dimer-dimer correlation functions each of which individually decays as the inverse-square of the distance. However, and this is what motivated the computation, it is not evident in advance that the various dimer correlators (which enter with different signs) do not sum to zero and hence the acutal result does not just yield a non-trivial constant but also assures us that it is not zero.
### B Finite density of classical holes: an incompressible state
As the two hole interaction derived previously has considerable spatial structure and a long range, the reader may wonder whether it leads to a condensation of a finite density of holes on the hexagonal backbone. This does not happen. To see this note that $`\beta v(r)`$ is independent of $`T`$, which is to say that there are flucutations even as $`T0`$—as indeed there are on account of the degeneracy. In the dilute limit, in which the two hole calculation is reliable, the effective interaction strength is proportional to the inverse square of the hole separation, i.e. to the (low) density and hence the hole gas remains in a liquid phase. Nevertheless, the long range of the interactions indicates that it exhibits substantial correlations coming from the entropic potential, and—on account of the form of the interaction—we suspect that such correlations will be algebraic as well. As we have already argued that the zero temperature state involves hole condensation on the hexagonal backbone driven by the hole kinetic energy, it appears that this phenomenon arises for both energetic and entropic reasons, and is thus quite robust so that it will characterize the entire temperature range $`TJ^z`$. (In all of this we are again referring to $`\eta <1/3`$, else the energetic attraction ignored in the entropic derivation will again dominate and drive phase separation.)
Finally we report a curiosity regarding classical holes for $`\eta <1/3`$. As the density of holes is increased, it is not clear whether the interaction finally becomes strong enough to produce phase separation into regions either fully occupied or maximally diluted (but respecting the hard core repulsion of the holes). However, even if it does not, there exists a density of holes at which long-range order is present.
The closest packing compatible with the ground-state constraint is obtained by placing the holes such that they cover one of two sublattices of the hexagaonal backbone of up spins depicted in Fig. 1b. At a hole density of 1/3, this sublattice is completely occupied by holes, at zero energy cost at $`O(J)`$ compared to the undiluted system.
Adding further holes above this density costs a magnetic energy of order $`J`$, and thus the chemical potential experiences a discontinuity at a hole density of 1/3. The state at this density, which has perfect unfrustrated antiferromagnetic order on the bipartite hexagonal lattice occupied by the remaining spins, is thus an incompressible state.
As the density of holes is increased further, phase separation occurs for certain. Completely depleted regions are created by removing one hexagonal ring after another as the system optimises its energy by removing those spins with the smallest number of neighbours. This is the same mechanism that would operate in any unfrustrated magnet. As mentioned above, for $`\eta >1/3`$, phase separation occurs for all hole densities.
## V Summary
We have considered the problem of dilute, slow holes in the triangular Ising antiferromagnet. We find that doping has two very different outcomes depending on the dimensionless strength $`\eta `$ of the density-density interaction. For $`\eta <1/3`$ we find quite generally that at all temperatures $`TJ^z`$ the holes tend to condense on a hexagonal backbone of frustrated bonds and that at $`T=0`$ they form a superconducting state that coexists with magnetic order and the breaking of translational symmetry in a “magnetic supersolid”. As a special case of this more general phenomenon, we were able to prove a “Nagaoka” theorem for a single hole at $`J^z=\mathrm{}`$. For $`\eta >1/3`$ in our short ranged model, we find phase separation. Whether the inclusion of the long range piece of the Coulomb interaction will change that outcome and stablize the magnetic supersolid and how much of this structure will persist to larger values of $`t/J^z`$ remain topics for future work.
## Acknowledgements
We are grateful to D. Priour and P. Chandra, who also kindly commented on the manuscript, for collaboration on closely related work. We would also like to acknowledge A. Yodh for a stimulating talk on entropic forces that led us to examine the same question in the context of this paper. This work was supported in part by grants from the Deutsche Forschungsgemeinschaft, NSF grant No. DMR-9978074, the A. P. Sloan Foundation and the David and Lucille Packard Foundation.
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# A CUPRONICKEL ROTATING BAND PION PRODUCTION TARGET FOR MUON COLLIDERS1
## 1 Introduction
High power pion production targets are required in current scenarios for muon colliders. The pion secondaries from protons on the target are captured in a solenoidal magnetic channel and decay into the muon bunches needed for cooling, acceleration and injection into the collider ring. Bunched proton beams of several megawatts will be needed for the currently specified muon currents : approximately $`6\times 10^{20}`$ muons of each sign at repetition rates of 15 Hertz and in bunches of up to $`4\times 10^{12}`$ muons per bunch. This is an extrapolation from today’s high power targets in rate of target heating, shock stresses and integrated radiation damage to the target.
Because of the high beam power, liquid metal jet targets have been the subject of much recent study and form the bulk of a proposed experimental R & D program of targetry studies that has recently been submitted to the BNL AGS Division. More conventional solid targets have several challenges. Along with concerns about shock heating stresses and radiation damage, it is challenging to design a cooling scenario consistent with both the large beam power and the small target cross sections that are needed for high pion yields.
This paper introduces the idea of a solid target in the form of a band that addresses this cooling issue by rotating the band to carry heat away from the targetry region and through a cooling channel.
Figures 1 and 2 give schematic views of the targetry setup we are considering and figure 3 shows the trajectory of the proton beam into the target band. It must be emphasized that details such as the rollers and cooling setup are only shown schematically and no effort has been put into their design.
The target band is enclosed in a 20 Tesla solenoidal magnetic pion capture magnet whose general design has previously been studied by the Muon Collider Collaboration (MCC). The major design modification specific to this particular geometry concerns the provision of entry and exit ports for the target band.
The high-power bunched proton beam strikes the target band at a glancing angle and travels along inside the target material for two nuclear interaction lengths before the curvature of the band brings it again to an exit point at the outer edge of the band. The beam is tilted at 150 milliradians to the longitudinal axis of the solenoidal magnet; MARS simulations described below show that this gives a larger pion yield than a beam parallel to the solenoid.
## 2 MARS Simulations of Target Heating and Pion Yield
Full MARS tracking and showering Monte Carlo simulations were conducted for a 16 GeV proton beam of 1$`\times `$10<sup>14</sup> ppp with a repetition rate of 15 Hz on Ni band (R=250 cm, 6 cm height, 0.6 cm thickness) in a 20 T solenoid of R<sub>a</sub>=7.5 cm half-aperture. Both untilted targets and targets tilted by $`\alpha `$=150 mrad were studied and detailed 3-dimensional maps of energy deposition densities were generated for input to the ANSYS stress analyses.
The yield per proton at 90 cm downstream from the central intersection of the beam with the target was determined for pions plus muons in the momentum range 0.05$`<`$p$`<`$0.8 GeV/c. The yields of positive and negative pions were, respectively, $`Y_+`$ = 0.491 and $`Y_{}`$ = 0.498 at $`\alpha `$=0 and $`Y_+`$ = 0.622 and $`Y_{}`$ = 0.612 at $`\alpha `$=150 mrad. Figure 5 shows the momentum spectra for all hadrons and figure 6 gives more detailed information for the pions. Figure 7 shows the time distribution for when these pions are formed and figure 8 shows several scatter plots to illustrate their distribution in phase space. These pion yields and densities in phase space are comparably good to the predictions for the best of the liquid jet targets under consideration.
The peak energy deposition density was found to be 68.6 J/g per pulse, corresponding to a temperature rise of $`\mathrm{\Delta }T`$=150.5C. This corresponded to a total power dissipation in the target of 0.324 MW at $`\alpha `$=150 mrad. Contributions to the deposited energy come from dE/dx from hadrons and muons (44%), electromagnetic showering (46%) and from absorbtion of sub-threshold particles (10%). Power dissipation in inner layer of tungsten shielding (7.5$`<`$r$`<`$15 cm) was also determined, and was found to be 0.624 MW at $`\alpha `$=0 and 0.766 MW at $`\alpha `$=150 mrad.
## 3 ANSYS Stress Simulations
The survivability of solid targets in the face of repeated shock heating is probably the most challenging problem faced in these scenarios for pion production for muon colliders.
To investigate this, we are beginning to conduct finite element computer simulations of the shock heating stresses using ANSYS, a commercial package that is very widely used for stress and thermal calculations. These studies are still at an early stage.
It is encouraging that the instantaneous energy deposition predicted by MARS of approximately 70 J/g per proton pulse is much less than the 500-600 J/g depositions in microsecond timescales that the Fermilab pbar source nickel target routinely operates at <sup>1</sup><sup>1</sup>1Editting note: this sentence has been modified from the original version, which erroneously referred to temperature rather than energy deposition.. Further, if the predicted stresses turn out to be higher than, say, 50% of the target’s tensile strength then possibilities exist for redimensioning the target and the proton spot size to reduce the stress.
## 4 Conclusions
In conclusion, initial studies indicate that cupronickel rotating band targets may well be a viable and attractive option to satisfy the difficult high power targetry requirements of muon colliders.
## 5 references
The Muon Collider Collaboration, “Status of Muon Collider Research and Development and Future Plans”, to be submitted to Phys. Rev. E.
Alessi et al., “An R & D Program for Targetry and Capture at a Muon-Collider Source - A Proposal to the BNL AGS Division”. Spokesperson Kirk T. McDonald, email: mcdonald@puphep.princeton.edu .
N.V Mokhov, “The MARS Code System User’s Guide, Version 13 (98)”, FERMILAB-FN-628 (Feb. 1998).
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# TURKU-FL-P35-00 Q-ball collisions in the MSSM: gauge-mediated supersymmetry breaking
## 1 Introduction
Various field theories can support stable non-topological solitons , Q-balls . A Q-ball is a coherent scalar condensate that carries a conserved charge, typically a $`U(1)`$ charge. Due to charge conservation, the Q-ball configuration is the ground state in the sector of fixed charge. Q-balls may have physical importance because the supersymmetric extensions of the Standard Model have scalar potentials that are suitable for Q-balls to exist in the theory. In particular, lepton or baryon number carrying Q-balls are present in the Minimal Supersymmetric Standard Model (MSSM) due to the existence of flat directions in the scalar sector of the theory .
If supersymmetry (SUSY) is broken at low energy scales by a gauge mediated mechanism, the scalar potential is completely flat for large enough field values. Therefore the energy per unit charge decreases like $`B^{1/4}`$ ($`B`$ is the baryon number) and for large enough $`B`$ one can have completely stable B-balls since there are no light enough baryon number carrying particles the Q-ball can decay into . If SUSY breaking is due to a hidden supergravity sector in the theory, the potential is not flat. Radiative corrections allow for $`Q`$-balls to exists but they are unstable and decay typically to quarks and nucleons . The decay progresses by evaporation from the surface of the Q-ball .
Q-balls can be cosmologically significant in various ways. Stable (or long living) Q-balls are natural candidates for dark matter . Their decay offers a way to understand the baryon to dark matter ratio and the baryon asymmetry of the universe . Q-balls can also protect the baryons from electroweak sphalerons and may be an important factor in considering the stability of neutron stars .
For Q-balls to be cosmologically significant one needs to have a mechanism that creates them in early stages of the evolution of the universe. Q-balls can be created in the early universe from an Affleck-Dine (AD) condensate . This process has been studied recently by numerical simulations where both the gauge- and gravity-mediated SUSY breaking scenarios were considered. In both cases Q-balls with various charges were seen to form from the condensate.
Collisions of Q-balls have been considered previously in various potentials -. In collisions were simulated on a one dimensional lattice in the gravity-mediated case and it was found that Q-balls typically merge, exchange charge or pass through each other . To our knowledge, collisions have not been studied in the gauge-mediated case previously and the gravity-mediated case has been studied in more than one dimension only in . In a recent paper Q-ball collisions were studied numerically in one, two and three dimensions in a polynomial potential . The main qualitative features of the collision processes were similar in all the three cases.
Since the charge of Q-ball can change in a collision process, they may play an important role in the evolution of the Q-ball charge distribution after their formation. On the other hand, the Q-ball charge distribution affects the cosmological role that Q-balls may have in the evolution of the universe. The effect of collisions can therefore be an important factor in evaluating the cosmological role of Q-balls.
In this paper we have studied Q-ball collisions in the gauge-mediated scenario on a two dimensional lattice. The gravity-mediated scenario has been analyzed in a previous paper .
## 2 Q-ball solutions
Consider a field theory with a U(1) symmetric scalar potential, $`U(\varphi )`$, with a global minimum at $`\varphi =0`$. The complex scalar field $`\varphi `$ carries a unit quantum number with respect to the $`U(1)`$-symmetry. The charge and energy of a field configuration $`\varphi `$ in D dimensions are
$$Q=\frac{1}{i}(\varphi ^{}_t\varphi \varphi _t\varphi ^{})d^Dx$$
(1)
and
$$E=[|\dot{\varphi }|^2+|\varphi |^2+U(\varphi ^{}\varphi )]d^Dx.$$
(2)
The single Q-ball solution is the minimum energy configuration in the sector of fixed charge. The Q-ball will be stable against radiative decays into $`\varphi `$-scalars if condition
$$E<mQ,$$
(3)
where $`m`$ is the mass of the $`\varphi `$-scalar, holds. It is then energetically favourable to store charge in a Q-ball rather than in form of free scalars.
Finding the minimum energy is straightforward using Lagrange multipliers. The Q-ball can be shown to be of the form
$$\varphi (x,t)=e^{i\omega t}\varphi (r),$$
(4)
where $`\varphi (x)`$ is now time independent and real, $`\omega `$ is the Q-ball frequency, $`|\omega |[0,m]`$ and $`\varphi `$ is spherically symmetric.
The charge of a Q-ball with spherical symmetry in D-dimensions is given by
$$Q=2\omega \varphi (r)^2d^Dr$$
(5)
and the equation of motion at a fixed $`\omega `$ is
$$\frac{d^2\varphi }{dr^2}+\frac{D1}{r}\frac{d\varphi }{dr}=\varphi \frac{U(\varphi ^2)}{\varphi ^2}\omega ^2\varphi .$$
(6)
To find the Q-ball solution we must solve (6) with the boundary conditions $`\varphi ^{}(0)=0,\varphi (\mathrm{})=0`$.
In the present paper we consider a potential of the form
$$U(\varphi )=m^4(1+\mathrm{log}(\frac{\varphi ^2}{m^2}))+\frac{\lambda ^2}{M^2}\varphi ^6,$$
(7)
with parameter values $`m=10^4`$ GeV, $`\lambda =0.5`$ and $`M=2.4\times 10^{18}`$ GeV. This corresponds to a potential along a flat direction that has been lifted by soft supersymmetry-breaking terms. Supersymmetry is broken here by a gauge-mediated mechanism as opposed to the gravity-mediated case analyzed previously .
We have calculated the charge and energy of Q-balls for different values of $`\omega `$. Energy vs. charge curves are shown in Figure 1(a). The axis scales are chosen differently for two and three dimensions; for two dimensions, $`Q_0=240(m/\mathrm{GeV})^2`$, $`E_0=mQ_0\mathrm{GeV}`$ and for three dimensions, $`Q_0=4.0(m/\mathrm{GeV})^2`$, $`E_0=mQ_0\mathrm{GeV}`$. The dashed line is the stability line, $`E=mQ`$, that indicates that the Q-balls considered here are stable with respect to scalar decays.
As from Fig. 1(a) can be seen the two and three dimensional energy vs. charge curves follow each other very closely. Also the Q-ball profiles in two and three dimensions have very similar shapes, Fig. 1(b). The similarities in the energy vs. charge curves and Q-ball profiles gives an indication that collisions studied in two and three dimensions are likely to give similar results. In it was also found that the two and three dimensional cases possess similar qualitative features. In contrast, the one dimensional case is fundamentally different from the higher dimensional cases. As from (6) can be seen, there is no dissipation term in one spatial dimension. Processes can then have properties that are not seen in higher dimensions. In the one dimensional simulations we have studied this also seemed to be the case.
Even though the general features of collision processes are similar in different potentials, the exact form of the potential is still important. From the basis of our simulations it appears that the details of collisions are dependent on the choice of potential.
One should also point out that in three dimensions new, ring-like intermediate states were seen in collisions at high velocities . The corresponding process in two dimensions is a right angle scattering -process that we did not observe in our simulations. However, at high velocities, $`𝒪(0.1)`$, there can be new types of processes that we have not seen in our low velocity simulations.
## 3 Collisions
We have studied collisions of Q-balls with equal charges in the potential (7). The range of $`\omega `$ for which numerical simulations have been done is $`\omega /m=0.15,0.20,0.30,0.60`$. In terms of charges these values of $`\omega `$ in the two dimensional case correspond to $`5.6,2.1,0.47,0.032`$ (in units of $`(\frac{m}{\mathrm{GeV}})^2`$). The initial velocities of the balls are allowed to have two values, $`v=10^3`$ or $`v=10^2`$.
Collisions are studied over different relative phases, $`0\mathrm{\Delta }\varphi \pi `$. Here we have defined the relative phase, $`\mathrm{\Delta }\varphi `$, to be the difference in individual Q-ball phases at the point where the distance between them would be at a minimum assuming there were no interaction between them. If the balls are of equal size, the point at which the relative phase difference is defined is irrelevant. In general, however, one needs to define the phase difference so that it is independent of the initial positions of the balls.
The position of a Q-ball is defined by the location of its maximum amplitude. We have varied the impact parameter to study the scattering cross-sections. As from Fig. 1 can be seen, the Q-balls studied here have thick walls so that one needs a definition for their size. In the gravity mediated case we defined the size of a ball by a Gaussian fit. Here, however, we have found that the Gaussian is a much poorer fit than in the gravity-mediated case. Instead we use a kink solution, $`\varphi =A+B\mathrm{tanh}(Cr+D)`$, and find that it approximates the numerical profiles well. The radius of the ball is defined as $`R=\frac{1}{C}(\mathrm{tanh}^1(\frac{A}{B})D)`$. This definition has the advantage that as the profile of the ball approaches the thin-wall solution, radius becomes defined in a natural way. It is worth noting that even though the profiles visually appear to be similar to the Q-ball profiles in the gravity-mediated scenario, they are fundamentally different from them and approach a purely thin-walled profile in the large charge limit.
A two dimensional, typically a $`300\times 300`$, lattice with continuous boundary conditions was used in all calculations.
### 3.1 Numerical Results
As in , the collision processes that we have observed can be roughly divided into three categories; fusion, charge exchange and elastic scattering. Fusion is defined as a process where most of the initial charge is in a single Q-ball after the collision and the rest of the charge is lost either as radiation or as small Q-balls. By charge exchange we mean a process where Q-balls exchange some of their charge while the total amount of charge carried by the two balls is essentially conserved. An elastic scattering is defined to be a scattering process where less than $`1\%`$ of the total charge is exchanged.
The type of a collision process is mainly dependent on the relative phase difference, $`\mathrm{\Delta }\varphi `$, between the colliding Q-balls. When $`\mathrm{\Delta }\varphi `$ is small the Q-balls fuse and form a larger ball. Excess charge is lost in form of radiation and small lumps of charge. As $`\mathrm{\Delta }\varphi `$ increases the Q-balls no longer fuse and start to scatter while exchanging a significant amount of their charge. The amount of charge that is exchanged in a collision decreases with increasing phase difference until the Q-balls scatter elastically. The change from a fusion process to an elastic scattering is typically rapid so that a significant amount of charge is exchanged only at a very narrow range of $`\mathrm{\Delta }\varphi `$. This is a different behavior from the gravity-mediated case where charge exchange was generally a much more dominant process . The qualitative effects of changing the relative phase difference are the same for the whole range of $`\omega `$:s and initial velocities that we have studied.
The effect of the impact parameter is much less pronounced than the phase difference on the type of collision. Typically if the balls fuse at zero impact parameter, they continue to fuse with an increasing impact parameter until at some point the balls start to scatter elastically or while exchanging little of their charge. The interaction probability, averaged over the different phases, also has a clear cut-off with respect to the impact parameter.
From the simulations we can now calculate the total, fusion and geometric cross-sections averaged over the relative phases. The charge exchange cross-section is typically much smaller than the other cross-sections and is not quoted here. The total cross-section includes all the different types of processes. The quoted cross-sections are the three dimensional cross-sections with the interaction radius taken from the two dimensional simulations.
The geometric, total and fusion cross-sections for $`v=10^3`$ are plotted in Fig. 2 and in Fig. 3 for $`v=10^2`$. In each case we have fitted a curve $`\sigma =A(\omega /m)^B`$ through the points. From the figures it is clear that the cross-sections increase with decreasing $`\omega `$. This is naturally due to the increasing size of the Q-balls with decreasing $`\omega `$. More importantly, one can also see that the ratio of the total cross-section to the geometrical cross-section decreases with decreasing $`\omega `$ i.e. increasing charge. This reflects the fact that as the profile of the Q-ball approaches the thin-wall type, the geometrical cross-section is an increasingly better approximation to the total cross-section. This is intuitively clear; as the balls become more and more thin-walled, the effects of the boundary diminish and the geometrical cross-section dominates the total cross-section. The ratio of the fusion cross-section to the geometrical cross-section is increasing, but only very slightly, with increasing charge. Extrapolating to large, thin-walled balls one can therefore conclude that the total cross-section of large Q-balls is well approximated by their geometrical cross-sections. The fusion cross-section of such balls can then be bounded by a fraction of the geometrical cross-section, $`\sigma _F/\sigma _G>\mathrm{\hspace{0.33em}0.6}`$. The geometrical cross-sections of large Q-balls can be estimated analytically, $`R\frac{1}{\sqrt{2}}m^1Q^{1/4}`$ so that $`\sigma _G2\pi m^2Q^{1/2}`$.
The effect of the increased velocity can also be seen from the figures. As the initial velocity is increased tenfold to $`10^2`$, the total and fusion cross-sections both decrease. This is expected since faster balls have less time to interact with each other and are more likely to pass without interacting.
Compared to the gravity-mediated case , the results presented here show both some similarities and some differences. In the gravity-mediated case the radius of the balls, and hence the geometrical cross-section, was approximately constant whereas here the radius of the ball varies greatly which obviously has an effect on the different cross-sections. Also the probabilities of different types of processes are different in the two scenarios: in the gravity-mediated case a charge exchange-process is much more likely to occur than a fusion process, here the charge exchange cross-section is much smaller. The fusion processes also can be distinguished in the two cases, in the gravity-mediated case charge was typically lost as small Q-balls whereas here most of the charge is lost as radiation and lumps of charge. This demonstrates clearly that the exact form of the potential is significant even though the general qualitative features are alike. One can also spot similarities between the gravity- and gauge-mediated scenarios; in both cases the fusion cross-sections increase with increasing balls and an increase in velocity decreases the total cross-sections.
## 4 Conclusions
In this paper we have studied Q-ball collisions in the MSSM with SUSY broken by a gauge mediated mechanism. It was found that Q-balls may fuse, exchange charge or scatter elastically in a collision depending on the relative phase difference between them. The probability of each process is dependent not only on the relative phase difference but also on the size of the balls. Larger balls are more probable to fuse in a collision whereas smaller balls are more likely to scatter either elastically or while exchanging some of their charge. Collisions can therefore alter the charge distribution of Q-balls quite significantly, provided that collisions are frequent enough.
Our simulations give an indication that the total cross-section, $`\sigma _{\mathrm{tot}}`$, approaches the geometrical cross-section, $`\sigma _G`$, as the Q-ball size grows so that in the thin-wall limit the total and geometrical cross-sections are equal. The fusion cross-section also grows with the ball size and on the basis of our results we can give an estimate for the fusion cross-section of large balls, $`\sigma _F>\mathrm{\hspace{0.33em}0.6}\sigma _G`$.
If collisions are to play a significant role in cosmology, the interaction rate must be large enough at some time in the evolution of the universe. If the average interaction rate is smaller than the Hubble rate after the Q-balls are formed, the distribution will ’freeze out’ and collisions will not alter the charge distribution. If, on the other hand, the collision rate is initially larger than the Hubble rate, collisions can affect the charge distribution until the distribution freezes out due to the expansion of the universe. Furthermore, if fusion processes dominate, the number density of Q-balls may decrease rapidly which can also freeze the Q-ball charge distribution.
On the basis of our results, the size and the relative phase differences are important in determining the evolution of the Q-ball charge distribution in the gauge mediated scenario. If the Q-balls are initially in the same phase, they typically fuse in a collision. The average size of a Q-ball then increases while the number density decreases. As our results show, the fusion probability increase with increasing Q-ball size so that larger balls are more likely to increase their size even further by collisions. The Q-ball interaction rate is dependent on several factors; clearly the initial Q-ball charge distribution, the interaction cross-section, number density and the velocity distribution affects the average interaction rate.
Q-balls may play a significant role in the past and present stages of the universe. To be able to address the question of their significance more decisively one needs to consider not only the initial charge distribution but also its evolution. Collisions may be an important factor and to study their effect in more detail poses an interesting question that motivates further study.
Acknowledgements. We thank K. Enqvist for discussions and the Center for Scientific Computing for computation resources. This work has been supported by the Finnish Graduate School in Nuclear and Particle Physics and by the Academy of Finland under the project no. 40677.
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# Molecular gas in blue compact dwarf galaxies
## 1 Introduction
A particular class of dwarf galaxies named Blue Compact Dwarf Galaxies (BCDGs, Sargent & Searle, \[Sargent & Searle, 1970\]) has seen increasing interest among astrophysicists because of their extreme current star forming activity which is in contrast to their apparent “youth” in terms of chemical evolution. BCDGs represent about 5% of all dwarfs \[Salzer 1989\], \[Sage et al., 1992\] and are among the smallest star forming galactic systems known.
One of their outstanding properties is that their optical spectra are dominated by lines characteristic of HII regions, which is the reason why they are frequently termed “HII galaxies”. From optical spectroscopy we know that many HII galaxies have low heavy element abundances, typically down by a factor of three up to more than twenty compared to the solar neighbourhood \[Kunth & Östlin, 2000\].
It quickly became clear that these objects must form stars in what is called a burst, otherwise the observed star formation rate would be in conflict with their total gas masses as derived from HI observations \[Thuan & Martin, 1981\]. Such a burst may last some $`10^8`$ yrs, with a time span between bursts of the order of 10<sup>9</sup> yrs. It has been suggested that interaction with companions might trigger their star formation \[Brinks, 1990\], but Taylor et al. \[Taylor et al., 1995\] found that only about 60% of HII galaxies have companions, often with masses about $`\frac{1}{10}`$ of the main galaxy.
One of the most interesting and important issues which has not been settled so far is the molecular gas content of these galaxies. Molecular hydrogen is believed to be the preponderant seed for star formation, so it is a natural assumption that large amounts of H<sub>2</sub> should be present in BCDGs. Yet the results have been anything but conclusive so far. Following early attempts to detect the CO line in BCDGs \[Tacconi & Young, 1984\], there have been a number of observations to confirm or reject those inconclusive measurements (e.g. Sage et al. \[Sage et al., 1992\], hereafter SSLH; Gondhalekar et al. \[Gondhalekar et al., 1998\]; Taylor, Kobulnicky & Skillman \[Taylor et al., 1998\], hereafter TKS). Surprisingly, the results remained partially contradictory, as for instance in the case of II Zw 40 (Arnault et al. \[Arnault et al., 1988\]; SSLH), although improved instrumentation had been involved.
Prompted by the difficulty to detect the CO line – relied upon as a good tracer of molecular hydrogen content – in BCDGs, some of the pertinent publications prematurely concluded that molecular gas is deficient in these systems. However, part of the difficulties to detect CO might have been due to beam filling and sensivitity problems. Taking e.g. 30 Dor in the LMC as a template giant star-forming complex, it is clear that, if placed at some larger distance and covered by the beam, it could have escaped detection, as CO brightness is rather low there, due to strong photo-dissociation in the high interstellar radiation field \[Cohen et al., 1988\]. The same might be true for BCDGs. In this case, high-sensitivity mapping might reveal previously undetected CO emission.
We have therefore conducted a search for CO in gas-rich (based on HI) HII galaxies, using the IRAM 30 m telescope. In contrast to previous projects (e.g. SSLH, Gondhalekar et al. \[Gondhalekar et al., 1998\]), our observations not only consisted of single pointings towards the brightest position in the galaxies, but involved mapping a number of positions in them, to detect possible gas concentrations away from photodissociation regions. Obtaining sensitive upper limits to the CO luminosity in these systems is as much a goal of this study as detecting emission.
In Section 2 we present details of our observations. In Section 3 we present our data, and compare it with previous results. In the subsequent section (Sec. 4) we discuss possible causes for the detections and non-detections. This section is divided into three subsections: Section 4.1 deals with the physical conditions of the gas derived from a LVG model; in Section 4.2 we analyze the relationship between metallicity and CO luminosity; finally, in Section 4.3 we discuss the $`X_{\mathrm{CO}}`$ factor problem, which has been heatedly debated in the past years and has not yet found a clear resolution. Our conclusions are presented in the last section (Sec. 5).
## 2 Observations
The CO observations were carried out in October 1996 with the IRAM 30 m Telescope on Pico Veleta in Spain. Our target galaxies – except Haro 2 – have been selected from the sample examined by \[Taylor et al., 1995\] in the 21 cm line of atomic hydrogen. Pointing positions were chosen from peak HI column densities, as seen in their high angular resolution VLA observations. The basic properties of the galaxies are listed in Table 1. In the $`J=10`$ transition (115 GHz) the HPBW is 22<sup>′′</sup>, whilst in the $`J=21`$ transition (230 GHz) it is $`12.^{\prime \prime }`$5. At the distance of the galaxies, 14.7 Mpc to 23.3 Mpc, the beam size at 115 GHz corresponds to 1.6 kpc to 2.5 kpc.
Two independent SIS receivers have been used simultaneously at each frequency. More than 90% of our observations had $`T_{sys}`$ 600 K both at 230 and at 115 GHz. An autocorrelator with a spectral resolution of 0.625 MHz at 115 MHz and 1.25 MHz at 230 MHz and a filter spectrometer with 1 MHz resolution were used. A baseline of zeroth or first order was always subtracted, and the spectra were summed up to improve the signal-to-noise ratio; the spectra were finally smoothed to roughly the same velocity resolution (5.2 km s<sup>-1</sup> for the $`J=10`$ and 4.8 km s<sup>-1</sup> for the $`J=21`$ transition). Spectra were obtained with a wobbling secondary mirror with a wobbler throw of $`\pm 4^{}`$ in azimuth. All temperatures throughout this article refer to main beam brightness with $`T_{mb}`$ derived using $`\eta _{mb}(115)=0.74`$ and $`\eta _{mb}(230)=0.45`$.
## 3 Results
Table 1 gives an overview of the observed galaxies and some of their known properties, and lists results of the galaxy-averaged CO spectra. We mapped most of the galaxies in our sample in order to cover most of the area where emission (perhaps in “hot spots”) might be present, and obtained low rms noise levels. Nevertheless, we detected CO only in Haro 2 and in UM 465. For all the other galaxies we obtained low upper limits. It is interesting to note that the SMC, with an H<sub>2</sub> mass of $`310^7`$ M would be just detectable with our sensitivity if placed at a distance of 15 Mpc.
In Table 2, all positions observed in all galaxies are listed, with the final rms obtained for single positions. In Fig. 1 and 3, the $`J=10`$ maps of Haro 2 and UM 465 are displayed.
In Table 1 we also give a summary of our results. For the detections, we give the parameters of a Gaussian fit which we obtained for the lines averaged over all positions. In all other cases, upper limits to the CO intensity were calculated, again based on the average of all positions. These upper limits were derived as $`I_{\mathrm{CO}}\sigma \mathrm{\Delta }v_{\text{ch}}\sqrt{N}`$, where $`\sigma `$ is the rms noise level obtained in the baseline range, $`\mathrm{\Delta }v_{\text{ch}}`$ is the velocity width of each channel and $`N`$ is the number of channels involved. We always assumed a total line width of 70 km s<sup>-1</sup>, as this is the average velocity width one would expect from the CO detections in the literature (see for example SSLH and Gondhalekar et al. \[Gondhalekar et al., 1998\]). In the case of detections, $`I_{\mathrm{CO}}`$ has been calculated based on Gaussian fits to the spectra. Based on the $`I_{\mathrm{CO}}`$ values, CO luminosities ($`L_{\mathrm{CO}}`$) were calculated. Those for UM 465 and Haro 2 are lower limits, because we did not completely map the CO gas. The other values are upper limits for the galaxies.
In the following sections we sumarize briefly our main results for the individual galaxies in comparison to previously published data where applicable. We also give a short description of the galaxies as they appear on the Digitized Sky Survey (DSS).
### 3.1 Haro 2
Haro 2 is a relatively well studied BCDG. Its metallicity is about $`\frac{1}{3}`$ solar. It has the shape and the brightness profiles of an elliptical galaxy, but possesses a brilliant blue nucleus which shows intense star formation. The absolute blue magnitude is M$`{}_{B}{}^{}=18.^\mathrm{m}4`$ \[Loose & Thuan, 1986\]. A comparison between the UV, optical and FIR spectra of Haro 2 with evolutionary population synthesis models has allowed to estimate the age of the youngest star formation episode to be 4 million years, followed by two older bursts, the younger of which was over 20 million years ago.
Our CO observations confirm the previous detection of CO in both the $`J=10`$ and $`J=21`$ transitions by SSLH and by Knapp & Rupen \[Knapp & Rupen, 1996\] in the latter transition. The emission is clearly extended, as seen in Fig. 1, with significant lines in all positions (see Tab. 3). Moreover, Fig. 2 suggests that the line in Haro 2 has two components which are seen at the same velocities in both transitions. Recently, observations with the IRAM interferometer fully confirmed this finding (Fritz et al., in prep.). The line ratio of the ($`21`$) to the ($`10`$) line seems to be independent of the velocity (see Fig. 2). Following the path of SSLH, we calculate two extreme line ratios, one assuming the source to be point-like, and one assuming a uniformly filled beam. In the latter case, since a beam filling factor of 1 is assumed, the ratio of the lines is effectively the ratio of integrated lines. This is also true if the filling factor is $`<1`$, but equal for both transitions (i.e. in the presence of large scale clumping). For a point source the other extreme is considered: one assumes that the (2-1)/(1-0) intensity ratio is overestimated by exactly the ratio of the two beam areas, so that the maximum line ratio must be divided by the ratio of the squares of the two beam widths ($`3.1`$). The results are shown in Table 4 and range from 0.5 to 1.5. The differences between our ratios and those of SSLH might be partly due to different main beam efficiencies used for the two transitions (not specified by SSLH).
### 3.2 UM 422
The dominant HII region of this dwarf galaxy is embedded in an extended faint irregular stellar body. We have obtained 5 independent spectra; none of them shows significant CO emission. UM 422 has also been observed with a single pointing by TKS using the NRAO 12 m telescope. We confirm their non-detection with a significantly lower upper limit.
### 3.3 UM 439
In the optical, UM 439 has a slightly elongated compact appearance with one prominent HII region south of the center. High resolution HI observations by van Zee, Skillman & Salzer \[Van Zee et al., 1998\] reveal that the star formation is taking place in the peak of the extended gas distribution. It was observed in CO by SSLH and by Gondhalekar et al. \[Gondhalekar et al., 1998\]. We confirm their non-detections. Our upper limit to $`I_{\mathrm{CO}}`$ is a factor of 2.4 lower than that of SSLH obtained with the NRAO 12 m telescope and a factor of 2.5 lower than that obtained with the Onsala 20 m telescope. Because we additionally observed 4 independent positions with higher angular resolution than those just mentioned, our upper limit to $`L_{\mathrm{CO}}`$ is significantly lower than those previously published.
### 3.4 UM 446
The stellar component of this galaxy, as detected in optical imaging, is very faint. We have observed it only in the central position. Our upper limit for CO is the first one ever published.
### 3.5 UM 452
We have observed this galaxy towards one position only, where the optical emission is strongest. The optical extent of the galaxy looks much smaller as compared to the HI. The mass of HI is quite low, $`M_{HI}=5\times 10^7`$ M \[Martin, 1999\].
### 3.6 UM 456
The star forming regions in this galaxy are confined to the center of an extended and distorted stellar component. Taylor et al. \[Taylor et al., 1995\] have detected two companions of UM 456. UM 456 A seems to be a pure “HI cloud” with no optical counterpart, whereas UM 456 B is seen both in HI and on optical images. Both companions seem to be gravitationally bound to UM 456. None of them shows CO; with our better rms we do not confirm the “marginal detection” of UM 456 by SSLH; our upper limit is a factor of 6 lower than their claimed detection.
### 3.7 UM 462
The two BCDGs UM 461 and UM 462 seem to form a bound system with a linear separation of about 70 kpc at a distance of 13.9 Mpc. Two centrally located knots of star formation dominate the optical image of UM 462. They are associated with the peak of the HI column density \[Van Zee et al., 1998\]. The claimed detection of CO in UM 462 by SSLH could not be confirmed; our upper limit is a factor of 2 lower. This galaxy had also been observed by Gondhalekar et al. \[Gondhalekar et al., 1998\] with a higher upper limit obtained at lower angular resolution.
### 3.8 UM 465
The optical appearance of this dwarf galaxy is circular in shape with an exponential law brightness distribution \[Doublier et al, 1997\]. The nuclear starburst and extended dust lanes and patches are well resolved by HST imaging \[Malkan et al., 1998\]. A faint nearby object was not detected in HI \[Taylor et al., 1995\] but confirmed as a physical neighbour \[Doublier et al, 1997\] using 6 m telescope spectroscopy. The HST images of this companion reveal an irregular structure. While SSLH reported a marginal detection in UM 465 of CO, the present work delivers a clear one. The CO emission is extended in this galaxy (see Fig. 3), but with a lower intensity than Haro 2. Only in one of the three positions no $`J=21`$ line was detected. The calculated line ratios are listed in Table 4, and for this galaxy they range between 0.4 and 1.3.
## 4 Discussion
Our deep observations corroborate the difficulty to detect CO in BCDGs. Only two of the galaxies observed show significant CO emission; in both cases it is extended in both transitions. We note here that these two sources are those with the highest metallicity in our sample. In the other sources, even with observations towards several positions, we were unable to find CO emission. One could be tempted to say that these galaxies are void of molecular gas, but this conclusion is premature because the relationship of CO emission and H<sub>2</sub> content in a galaxy depends on many factors (Maloney & Black \[Maloney & Black, 1988\], Israel \[Israel et al., 1997\]), some of which are not fully understood. Therefore, one can only conclude from the detection of CO that H<sub>2</sub> is present, whereas the absence of CO does not necessarily imply a lack of H<sub>2</sub>.
### 4.1 Physical conditions of the gas
As a first step in our analysis, we try to infer some information on the physical properties of the CO-emitting gas in BCDGs. We make use of a large velocity gradient (LVG) model to predict ratios for the lowest CO lines. The basic LVG assumption is that of a systematic monotonic velocity gradient, which allows to treat the molecular excitation as a local problem (see de Jong, Chu & Dalgarno, \[De Jong et al., 1975\] and White \[White, 1977\] for details). This is certainly an idealization of extra-galactic cloud complexes; however, an LVG code does not require detailed knowledge of the velocity field. As a first rough estimate this model further assumes constant density and kinetic temperature within the molecular cloud. In Fig. 4, we show an example of the dependence of line ratios on $`n`$(H<sub>2</sub>) and $`T_{kin}`$. For this figure, $`\frac{[\mathrm{CO}/\mathrm{H}_2]}{|\stackrel{}{v}|}=210^5`$ ( km s<sup>-1</sup>pc$`{}_{}{}^{1})^1`$ has been assumed. This corresponds to a velocity gradient of 1 km s<sup>-1</sup>pc<sup>-1</sup> and an abundance of \[CO/H<sub>2</sub>\]$`=210^5`$ or, due to how the LVG code is constructed, to a velocity gradient of 5 km s<sup>-1</sup>pc<sup>-1</sup> and \[CO/H<sub>2</sub>\]=10<sup>-4</sup>. The figure shows the ratios of the intensities of the $`21`$ to $`10`$ transitions, $``$<sub>2/1</sub>, and that of $`32`$ and the $`21`$ lines, $``$<sub>3/2</sub>.
As discussed above, for Haro 2 and UM 465 the line ratios are $`0.49_{2/1}1.51`$ and $`0.42_{2/1}1.30`$, respectively, depending on the filling of the sources in our beam. From our limited mapping of the two galaxies, we know that the sources are extended and the “point source” limit can be firmly excluded. On the other hand, observations of Haro 2 obtained with the Plateu de Bure Interferometer (Fritz et al., in prep.) show that the galaxy does not fill the beam of the IRAM 30 m telescope uniformly. It is then reasonable to expect a value for the line intensity ratios in between those listed in Table 4, and may thus be close to unity.
Adopting a ratio $`_{2/1}`$ between 0.8 and 1.0, we can derive from Fig. 4 that the gas is either at high temperatures with medium densities ($``$ a few hundred cm<sup>-3</sup>) or at high densities ($`2000`$ cm<sup>-3</sup>) and low temperatures ($`T_{kin}50`$ K). More stringent limits on volume density and kinetic temperatures require observations of higher <sup>12</sup>CO transitions and/or <sup>13</sup>CO transitions. These transitions are expected to be very weak though, and thus difficult and time-consuming.
### 4.2 Dependence of the CO luminosity on metallicity and absolute blue magnitude
We subsequently examine the relation of the CO emission and metallicity. Because only the two galaxies in our sample with the highest metallicities are detected in CO, one could expect that CO luminosity depends on metallicity. Therefore, we plot our CO luminosities ($`L_{\mathrm{CO}}`$) – listed in Table 1 – vs. the metallicities of the galaxies of our sample. Since not all of the galaxies have known metallicities, we are left with only 6 galaxies. These are shown in Fig. 5. In this figure we also include the data points given by TKS because they represent the most comprehensive sample of CO observations of dwarf galaxies with metallicity determinations. We only selected those dwarf galaxies from the sample with metallicities better determined than 0.1 dex; these galaxies are, however, not necessarily classical BCDGs.
We have chosen the CO luminosity because it is largely independent of distance, although the error in the distance determination enters as the square in the CO luminosity. Furthermore, the CO luminosity is directly proportional to the H<sub>2</sub> mass corrected for helium ($`M(\mathrm{H}_2)=2.2X_{\mathrm{CO}}L_{\mathrm{CO}}`$), with $`M`$ in M and $`L_{\mathrm{CO}}`$ in K km s<sup>-1</sup> pc<sup>2</sup>. $`X_{\mathrm{CO}}`$ (in units of $`10^{20}`$ molecules cm<sup>-2</sup> (K km s<sup>-1</sup>)<sup>-1</sup>) is the well-known but poorly determined $`X_{\mathrm{CO}}`$ factor which relates the molecular hydrogen column density to the observed integrated CO line intensity ($`X_{\mathrm{CO}}N(\mathrm{H}_2)/I_{\mathrm{CO}}`$).
Fig. 5 shows a general trend that galaxies with higher metallicities have higher CO luminosities, although no functional correlation is visible. TKS have recently proposed that galaxies with metallicities below 7.9 are basically undetectable in CO. Our data do not contradict this finding. For galaxies close to that limit we were only able to derive upper limits. We note, however, that even above that limit there are galaxies not detected in CO.
Although it is qualitatively expected that a higher metallicity leads to a higher $`L_{\mathrm{CO}}`$ because of the availability of the building blocks of the CO molecule, Fig. 5 shows that the oxygen abundance cannot be the only factor influencing the CO luminosity. Because BCDGs are actively star-forming galaxies, the UV radiation field may be locally high. This plays two conflicting roles for CO: on the one hand, it heats the gas, so that the excitation temperature is higher and CO is brighter; on the other hand, if hard enough, it destroys CO via photodissociation, thus the CO emission becomes weaker. Pak et al. \[Pak et al., 1998\] and Bolatto, Jackson & Ingalls \[Bolatto et al., 1999\] have studied these effects and found that the CO emitting regions are effectively smaller in low-metallicity environments and most of the carbon is present in atomic form. The net effect in galaxies with low metallicities is that, due to low beam filling of the clouds, the molecular gas becomes invisible in the CO lines and might be better traced in the fine structure lines of CI and CII. This different beam filling might be the cause for the dependence of the integrated CO line intensity $`I_{\mathrm{CO}}`$ on the metallicity found by TKS.
To study the effect of the radiation field on the CO luminosity, we labelled the values in the CO luminosity-metallicity plane with the absolute blue magnitudes of the galaxies (right side of Fig. 5). Again there is no simple relation between any two of these quantities. However, there is evidence for an influence of both metallicity and absolute blue magnitude on the CO luminosity, meaning that higher blue magnitudes lead to higher CO luminosity for a given metallicity. It also appears that at lower metallicities a higher absolute blue magnitude is necessary to reach a certain CO luminosity. Clearly, more well-observed galaxies are necessary to study the relation between the three quantities.
### 4.3 Molecular gas masses and the $`X_{\mathrm{CO}}`$ factor
Directly linked to the question of the dependence of the CO luminosity on metallicity and radiation field is the question of which $`X_{\mathrm{CO}}`$ factor is applicable to low-metallicity galaxies. A number of studies have examined possible correlations between $`X_{\mathrm{CO}}`$ and the metallicity (e.g. Dettmar & Heithausen \[Dettmar & Heithausen, 1989\], Wilson \[Wilson, 1995\], Verter & Hodge \[Verter & Hodge, 1995\], Arimoto, Sofue & Tsujimoto \[Arimoto et al., 1996\]). Klein \[Klein, 1999\] has proposed an additional dependence of $`X_{\mathrm{CO}}`$ on the cosmic ray flux as judged from the radio continuum brightness.
An independent determination of the $`X_{\mathrm{CO}}`$ factor for our BCDG sample would be useful to determine their molecular masses; this is, however, beyond the scope of this paper. A reliable value for $`X_{\mathrm{CO}}`$ has been established for the disk of the Milky Way, $`X_{\mathrm{MW}}`$. The currently best accepted value is $`X_{\mathrm{MW}}=1.610^{20}`$ molecules cm<sup>-2</sup> K<sup>-1</sup> km s<sup>-1</sup> \[Hunter et al., 1997\]. $`X_{\mathrm{CO}}`$ factors determined for galaxies with lower metallicities are usually significantly higher (e.g. Cohen et al. \[Cohen et al., 1988\], Dettmar & Heithausen \[Dettmar & Heithausen, 1989\]).
In the following, we assume that the correlation between $`X_{\mathrm{CO}}`$, derived from a virialization analysis of several galaxies, and the oxygen abundance obtained by Arimoto et al. \[Arimoto et al., 1996\] gives $`X_{\mathrm{CO}}`$ factors applicable to our galaxy sample. H<sub>2</sub> masses derived under this assumption are listed in Table 5. Also given are the HI masses as derived by \[Taylor et al., 1995\] for the UM galaxies and by SSLH for Haro 2. One remarkable result from this calculation is that those galaxies undetected in CO have upper limits on the molecular gas mass significantly below the HI mass, whereas in the two galaxies where CO is detected, HI and H<sub>2</sub> masses are about the same.
The $`X_{\mathrm{CO}}`$ factors used here to derive molecular gas masses are global factors. Studying the CO emission of the Magellanic Clouds with different angular resolutions Rubio, Lequeux & Boulanger \[Rubio et al., 1993\] noted that the derived $`X_{\mathrm{CO}}`$ factor depends on the linear resolution, implying that the local $`X_{\mathrm{CO}}`$ value is lower than the global one.
In order to calculate the local $`X_{\mathrm{CO}}`$ we use an LVG approximation assuming that H<sub>2</sub> and CO share the same volume. The LVG approximation can then be used to derive the H<sub>2</sub> column density, $`N(\mathrm{H}_2)`$, from $`N(\mathrm{H}_2)=n(\mathrm{H}_2)\mathrm{\Delta }v/|\stackrel{}{v}|`$, where $`n(\mathrm{H}_2)`$ is the H<sub>2</sub> volume density, $`\mathrm{\Delta }v`$ is the line width and $`|\stackrel{}{v}|`$ is the velocity gradient. Since $`I_{\mathrm{CO}}T_{mb}\mathrm{\Delta }v`$, one can derive $`X_{\mathrm{CO}}`$ from $`X_{\mathrm{CO}}=N(\mathrm{H}_2)/I_{\mathrm{CO}}n(\mathrm{H}_2)/|\stackrel{}{v}|T_{mb}`$.
The dependence of the local $`X_{\mathrm{CO}}`$ on varying volume density and kinetic temperature is shown in Fig. 6. The velocity gradient is fixed to 2 km s<sup>-1</sup>pc<sup>-1</sup>. In the top panel, the abundance is \[CO/H<sub>2</sub>\]=10<sup>-4</sup>; in the bottom panel \[CO/H<sub>2</sub>\]= 10<sup>-5</sup>. These are extreme values for the metallicities of BCDGs. $`X_{\mathrm{MW}}`$ is indicated by the thick solid line. We note that, once the density becomes high enough, $`X_{\mathrm{CO}}`$ will not change significantly any more, due to the growth of the optical depth, and $`X_{\mathrm{CO}}`$ is independent of the abundance. On the other hand, at lower densities and kinetic temperatures, $`X_{\mathrm{CO}}`$ changes significantly depending on $`n`$(H<sub>2</sub>) and $`T_{kin}`$ for a given metallicity.
These simple calculations indicate that $`X_{\mathrm{CO}}`$ does not only depend on the metallicity of a galaxy. Physical parameters, such as average volume density and kinetic temperature, play important roles, if the density is below $`10^4`$ cm<sup>-3</sup> and/or the kinetic temperature is below 50 K, values found in many molecular clouds in the Milky Way. The $`X_{\mathrm{CO}}`$-factor, that we find for a standard density of $`10^3`$ cm<sup>-3</sup> and a low-metallicity environment, is low, i.e. close to Galactic. In contrast, $`X_{\mathrm{CO}}`$ derived from the formula of Arimoto et al. \[Arimoto et al., 1996\] is higher by an order of magnitude. This supports the concept of large amounts of hidden gas - either atomic or molecular.
Interferometric observations of BCDGs might help to resolve this issue. Such a study of the nearby post-starburst dwarf NGC 1569, a galaxy that may be considered as a nearby BCDG in a post-starburst phase, yielded a rather high value of $`X_{\mathrm{CO}}=6.6X_{\mathrm{MW}}`$, based on virial masses of resolved GMCs \[Taylor et al., 1999\]. This indicates that this method, which is also the basis of Arimoto’s formula, is sensitive to the ‘hidden’ H<sub>2</sub> and tends to yield global values for $`X_{\mathrm{CO}}`$. In contrast, we expect to find local values if line ratio studies encompassing several CO isotopomers become available, since these studies directly probe the physics of the gas from which the CO emission arises.
## 5 Conclusions
We have searched for emission from the <sup>12</sup>CO ($`J=10`$ and $`J=21`$) transitions in 10 dwarf galaxies, 8 of which are BCDGs and 2 are the companions of one of these. We detected CO in 2 of them (Haro 2 and UM 465) and found it to be extended in both galaxies. Although we mapped part of the other galaxies, we were unable to detected CO. We obtained very stringent upper limits. We could not confirm the “marginal detection” of CO in UM 456 and UM 462 previously reported by SSLH.
The observed line ratios of the $`21`$ to $`10`$ transitions are not very sensitive to changes in the kinetic temperature. Modelling the ratio with a simple LVG code helps only to exclude low densities. Higher CO transitions and/or observations of CO isotopomers will help to get more stringent limits on these physical parameters.
We could not find any simple relation between metallicity and CO luminosity. Molecular gas masses for the galaxies are derived assuming the relation between $`X_{\mathrm{CO}}`$ and metallicity given by \[Arimoto et al., 1996\]. We find that for those galaxies detected in the CO lines the molecular gas mass is comparable to the HI mass, whereas for those galaxies undetected in CO the HI mass is significantly larger than the limits on the molecular gas mass.
Even in the sources where CO has not been detected, we do not argue against the presence of H<sub>2</sub>. While it is certainly possible that in the extreme environment of a BCDG not just CO but also H<sub>2</sub> is destroyed, at least in regions close to young massive stellar clusters, a picture in which a large amount of H<sub>2</sub> exists without CO is attractive. Sensitive observations of CI and CII in these galaxies would thus be desirable in the future to shed light on this issue.
## Acknowledgements
L.T.B. would like to thank Prof. Loretta Gregorini and the Socrates/Erasmus project which made this exchange possible and financed it, the Faculty of Science of the University of Bologna and the C.N.A.A. (Consorzio Nazionale per l’Astronomia e l’Astrofisica) for grants supporting this work. This project was supported by the Deutsche Forschungsgemeinschaft via the Graduiertenkolleg “The Magellanic Clouds and other Dwarf Galaxies”.
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# 1 Introduction
## 1 Introduction
Recently, it has been realized that renormalizations in a softly broken SUSY theory are not independent but follow from those of an unbroken SUSY theory. According to the approach advocated in Refs. , one can perform the renormalization of a softly broken SUSY theory in the following straightforward way:
One takes renormalization constants of a rigid theory, calculated in some massless scheme, substitutes for the rigid couplings (gauge and Yukawa) their modified expressions, that depend on a Grassmannian variable, and expands over this variable.
This gives renormalization constants for the soft terms. Differentiating them with respect to a scale, one can find corresponding renormalization-group equations.
Thus, the soft-term renormalizations are not independent but can be calculated from the known renormalizations of a rigid theory with the help of the differential operators. Explicit form of these operators has been found in a general case and in some particular models like SUSY GUTs or the MSSM . The same expressions have been obtained also in a somewhat different approach in Ref. .
There is, however, some minor difference. The authors of have used the component approach, while in , the superfield formalism is exploited. This creates the usual difference in gauge-fixing and ghost field terms and in the renormalization scheme. The latter is related to the choice of regularization. In , the dimensional reduction (DRED) regularization is used. In this case, one is bounded to introduce the so-called $`ϵ`$-scalars to compensate the lack of bosonic degrees of freedom in 4-2$`ϵ`$ dimensions. These $`ϵ`$-scalars in due course of renormalization acquire a soft mass that enters into the RG equations for soft masses of physical scalar particles. This problem has been discussed in . If one gets rid of the $`ϵ`$-scalar mass by changing the renormalization scheme, $`\mathrm{DRED}\mathrm{DRED}^{}`$, there appears an additional term in RG equations for the soft scalar masses called X . This term is absent in RG equations in Refs. .
We have to admit that, indeed in our approach, though the $`ϵ`$-scalars in the superfield formalism are absent, that term appears in higher orders and is related to the soft masses of other unphysical particles, the auxiliary gauge fields. We show below how it emerges in the superfield formalism and coincides with that of Ref. . Thus, the two approaches finally merge.
## 2 Massive Auxiliary Fields
Consider an arbitrary $`N=1`$ SUSY gauge theory with unbroken SUSY within the superfield formalism. The Lagrangian of a rigid theory is given by
$`_{rigid}`$ $`=`$ $`{\displaystyle d^2\theta \frac{1}{4g^2}\mathrm{Tr}W^\alpha W_\alpha }+{\displaystyle d^2\overline{\theta }\frac{1}{4g^2}\mathrm{Tr}\overline{W}_{\dot{\alpha }}\overline{W}^{\dot{\alpha }}}`$
$`+`$ $`{\displaystyle d^2\theta d^2\overline{\theta }\overline{\mathrm{\Phi }}^i(e^V)_i^j\mathrm{\Phi }_j}+{\displaystyle d^2\theta 𝒲}+{\displaystyle d^2\overline{\theta }\overline{𝒲}},`$
where
$$W_\alpha =\frac{1}{4}\overline{D}^2e^VD_\alpha e^V,\overline{W}_{\dot{\alpha }}=\frac{1}{4}D^2e^V\overline{D}_{\dot{\alpha }}e^V,$$
are the gauge field strength tensors and the superpotential $`𝒲`$ has the form
$$𝒲=\frac{1}{6}y^{ijk}\mathrm{\Phi }_i\mathrm{\Phi }_j\mathrm{\Phi }_k+\frac{1}{2}M^{ij}\mathrm{\Phi }_i\mathrm{\Phi }_j.$$
(2)
To fix the gauge, the usual gauge-fixing term can be introduced. It is useful to choose it in the form
$$_{gaugefixing}=\frac{1}{16}d^2\theta d^2\overline{\theta }\mathrm{Tr}\left(\overline{f}f+f\overline{f}\right)$$
(3)
where the gauge fixing condition is taken as
$`f=\overline{D}^2{\displaystyle \frac{V}{\sqrt{\xi g^2}}},\overline{f}=D^2{\displaystyle \frac{V}{\sqrt{\xi g^2}}}.`$ (4)
Then, the corresponding ghost term is
$$_{ghost}=id^2\theta \frac{1}{4}\mathrm{Tr}b\delta _cfid^2\overline{\theta }\frac{1}{4}\mathrm{Tr}\overline{b}\delta _{\overline{c}}\overline{f},$$
(5)
where $`c`$ and $`b`$ are the Faddeev–Popov ghost and antighost chiral superfields, respectively, and $`\delta _c`$ is the gauge transformation with the replacement of gauge superfield parameters $`\mathrm{\Lambda }(\overline{\mathrm{\Lambda }})`$ by chiral (antichiral) ghost fields $`c(\overline{c})`$.
For our choice of the gauge-fixing condition, the gauge transformation of $`f`$ looks like
$$\delta _\mathrm{\Lambda }f=\overline{D}^2\delta _\mathrm{\Lambda }\frac{V}{\sqrt{\xi g^2}}=i\overline{D}^2\frac{1}{\sqrt{\xi g^2}}_{V/2}[\mathrm{\Lambda }+\overline{\mathrm{\Lambda }}+\mathrm{coth}(_{V/2})(\mathrm{\Lambda }\overline{\mathrm{\Lambda }})],$$
(6)
where $`_XY[X,Y]`$. Equation (5) then takes the form
$`_{ghost}`$ $`=`$ $`{\displaystyle d^2\theta \frac{1}{4}\mathrm{Tr}b\overline{D}^2\frac{1}{\sqrt{\xi g^2}}_{V/2}[c+\overline{c}+\mathrm{coth}(_{V/2})(c\overline{c})]}+h.c.`$
$`=`$ $`{\displaystyle d^2\theta d^2\overline{\theta }\mathrm{Tr}\left(\frac{b+\overline{b}}{\sqrt{\xi g^2}}\right)_{V/2}[c+\overline{c}+\mathrm{coth}(_{V/2})(c\overline{c})]}`$
$`={\displaystyle d^2\theta d^2\overline{\theta }\mathrm{Tr}\left(\frac{b+\overline{b}}{\sqrt{\xi g^2}}\right)\left(\left(c\overline{c}\right)+\frac{1}{2}[V,\left(c+\overline{c}\right)]+\frac{1}{12}[V,[V,\left(c\overline{c}\right)]]+\mathrm{}\right)}.`$
The resulting Lagrangian together with the gauge-fixing and the ghost terms are invariant under the BRST transformations. For a rigid theory in our normalization of the fields, they have the form
$`\delta V=ϵ_{V/2}[c+\overline{c}+\mathrm{coth}(_{V/2})(c\overline{c})],`$
$`\delta c^a={\displaystyle \frac{i}{2}}ϵf^{abc}c^bc^c,\delta \overline{c}^a={\displaystyle \frac{i}{2}}ϵf^{abc}\overline{c}^b\overline{c}^c,`$
$`\delta b^a={\displaystyle \frac{1}{8}}ϵ\overline{D}^2\overline{f}^a,\delta \overline{b}^a={\displaystyle \frac{1}{8}}ϵD^2f^a.`$ (8)
To perform the SUSY breaking, that satisfies the requirement of ”softness”, one can introduce a gaugino mass term as well as cubic and quadratic interactions of scalar superpartners of the matter fields
$`_{softbreaking}`$ $`=`$ $`[{\displaystyle \frac{M}{2}}\lambda \lambda +{\displaystyle \frac{1}{6}}A^{ijk}\varphi _i\varphi _j\varphi _k+{\displaystyle \frac{1}{2}}B^{ij}\varphi _i\varphi _j+h.c.]+(m^2)_j^i\varphi _i^{}\varphi ^j,`$ (9)
where $`\lambda `$ is the gaugino field, and $`\varphi _i`$ is the lowest component of the chiral matter superfield.
One can rewrite the Lagrangian (9) in terms of N=1 superfields introducing the external spurion superfields $`\eta =\theta ^2`$ and $`\overline{\eta }=\overline{\theta }^2`$, where $`\theta `$ and $`\overline{\theta }`$ are Grassmannian parameters, as
$`_{soft}`$ $`=`$ $`{\displaystyle d^2\theta \frac{1}{4g^2}(12M\theta ^2)\mathrm{Tr}W^\alpha W_\alpha }+{\displaystyle d^2\overline{\theta }\frac{1}{4g^2}(12\overline{M}\overline{\theta }^2)\mathrm{Tr}\overline{W}^{\dot{\alpha }}\overline{W}_{\dot{\alpha }}}`$
$`+{\displaystyle d^2\theta d^2\overline{\theta }\overline{\mathrm{\Phi }}^i(\delta _i^k(m^2)_i^k\eta \overline{\eta })(e^V)_k^j\mathrm{\Phi }_j}`$
$`+{\displaystyle d^2\theta \left[\frac{1}{6}(y^{ijk}A^{ijk}\eta )\mathrm{\Phi }_i\mathrm{\Phi }_j\mathrm{\Phi }_k+\frac{1}{2}(M^{ij}B^{ij}\eta )\mathrm{\Phi }_i\mathrm{\Phi }_j\right]}+h.c.`$
Thus, one can interpret the soft terms as the modification of the couplings of a rigid theory. The couplings become external superfields depending on Grassmannian parameters $`\theta `$ and $`\overline{\theta }`$. To get the explicit expression for the modified couplings, consider eqs.(2). The first two terms give
$$\frac{1}{g^2}\frac{1}{\stackrel{~}{g}^2}=\frac{1M\theta ^2\overline{M}\overline{\theta }^2}{g^2}.$$
(11)
Since the gauge field strength tensors $`W_\alpha `$ ($`\overline{W}_\alpha `$) are chiral (antichiral) superfields, they enter into the chiral (antichiral) integrands in eq.(2), respectively. Correspondingly, the $`M\theta ^2`$ term of eq.(11) contributes to the chiral integral, while the $`\overline{M}\overline{\theta }^2`$ term contributes to the antichiral one. There is no $`\theta ^2\overline{\theta }^2`$ term in eq.(11), since it is neither chiral, no antichiral and gives no contribution to eq.(2).
However, modifying the gauge coupling in the gauge part of the Lagrangian, one has to do the same in the gauge-fixing (4) and ghost (2) parts in order to preserve the BRST invariance. Here one has the integral over the whole superspace rather than the chiral one. This means that if one adds to eq.(11) a term proportional to $`\theta ^2\overline{\theta }^2`$, it gives a nonzero contribution. Moreover, even if this term is not added, it reappears as a result of renormalization.
We suggest the following modification of eq.(11)
$$\frac{1}{g^2}\frac{1}{\stackrel{~}{g}^2}=\frac{1M\theta ^2\overline{M}\overline{\theta }^2\mathrm{\Delta }\theta ^2\overline{\theta }^2}{g^2},$$
(12)
which gives the final expression for the soft gauge coupling
$`\stackrel{~}{g}^2=g^2\left(1+M\theta ^2+\overline{M}\overline{\theta }^2+2M\overline{M}\theta ^2\overline{\theta }^2+\mathrm{\Delta }\theta ^2\overline{\theta }^2\right).`$ (13)
In our previous papers , this $`\mathrm{\Delta }`$ term was absent. It will be clear below that it is self-consistent to put $`\mathrm{\Delta }=0`$ in the lowest order of perturbation theory, but it appears in higher orders due to renormalizations.
One has to take into account, however, that, since the gauge-fixing parameter $`\xi `$ may be considered as an additional coupling, it also becomes an external superfield and has to be modified. The soft expression can be written as
$$\stackrel{~}{\xi }=\xi \left(1+x\theta ^2+\overline{x}\overline{\theta }^2+(x\overline{x}+z)\theta ^2\overline{\theta }^2\right),$$
(14)
where parameters $`x`$ and $`z`$ can be obtained by solving the corresponding RG equation (see Appendix A).
Having this in mind, we perform the following modification of the gauge fixing condition (4) first used in
$$f\overline{D}^2\frac{V}{\sqrt{\stackrel{~}{\xi }\stackrel{~}{g}^2}},\overline{f}D^2\frac{V}{\sqrt{\stackrel{~}{\xi }\stackrel{~}{g}^2}}.$$
(15)
Then, the gauge-fixing term (3) becomes
$$_{gaugefixing}=\frac{1}{8}d^2\theta d^2\overline{\theta }\mathrm{Tr}\left(\overline{D}^2\frac{V}{\sqrt{\stackrel{~}{\xi }\stackrel{~}{g}^2}}D^2\frac{V}{\sqrt{\stackrel{~}{\xi }\stackrel{~}{g}^2}}\right),$$
(16)
This leads to the corresponding modification of the associated ghost term (5)
$`_{ghost}`$ $`=`$ $`{\displaystyle d^2\theta d^2\overline{\theta }\mathrm{Tr}\frac{1}{\sqrt{\stackrel{~}{\xi }\stackrel{~}{g}^2}}\left(b+\overline{b}\right)_{V/2}[c+\overline{c}+\mathrm{coth}(_{V/2})(c\overline{c})]}.`$ (17)
To understand the meaning of the $`\mathrm{\Delta }`$ term, consider the quadratic part of the ghost Lagrangian (17)
$`_{ghost}^{(2)}`$ $`=`$ $`{\displaystyle d^2\theta d^2\overline{\theta }\mathrm{Tr}\frac{1}{\sqrt{\xi g^2}}\left(1\frac{1}{2}M\xi \theta ^2\frac{1}{2}\overline{M}\xi \overline{\theta }^2\frac{1}{2}\mathrm{\Delta }\xi \theta ^2\overline{\theta }^2\right)(b+\overline{b})\left(c\overline{c}\right)}`$
$`=`$ $`{\displaystyle d^2\theta d^2\overline{\theta }\mathrm{Tr}\frac{1}{\sqrt{\xi g^2}}\left(1\frac{1}{2}\mathrm{\Delta }\xi \theta ^2\overline{\theta }^2\right)(b+\overline{b})\left(c\overline{c}\right)}`$
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle d^2\theta \mathrm{Tr}\frac{1}{\sqrt{\xi g^2}}M\xi bc}+{\displaystyle \frac{1}{2}}{\displaystyle d^2\overline{\theta }\mathrm{Tr}\frac{1}{\sqrt{\xi g^2}}\overline{M}\xi \overline{b}\overline{c}},`$
where we have used the explicit form of $`\stackrel{~}{\xi }`$ given in Appendix A.
If one compares this expression with the usual Lagrangian for the matter fields (2), one finds an obvious identification of the second line with the soft scalar mass term and the third line with the mass term in a superpotential. Thus, $`\mathrm{\Delta }`$ plays the role of a soft mass providing the splitting in the ghost supermultiplet.
The other place where the $`\mathrm{\Delta }`$-term appears is the gauge-fixing term (16). Here it manifests itself as a soft mass of the auxiliary gauge field, one of the scalar components of the gauge superfield $`V`$.
To see this, consider the gauge-fixing term (16) in more detail. Expanding the vector superfield $`V(x,\theta ,\overline{\theta })`$ in components
$`V(x,\theta ,\overline{\theta })`$ $`=`$ $`(x)+i\theta \chi (x)i\overline{\theta }\overline{\chi }(x)+{\displaystyle \frac{i}{2}}\theta \theta N(x){\displaystyle \frac{i}{2}}\overline{\theta }\overline{\theta }\overline{N}(x)\theta \sigma ^\mu \overline{\theta }v_\mu (x)`$
$`+i\theta \theta \overline{\theta }[\overline{\lambda }(x)+{\displaystyle \frac{i}{2}}\overline{\sigma }^\mu _\mu \chi (x)]i\overline{\theta }\overline{\theta }\theta [\lambda +{\displaystyle \frac{i}{2}}\sigma ^\mu _\mu \overline{\chi }(x)]+{\displaystyle \frac{1}{2}}\theta \theta \overline{\theta }\overline{\theta }[D(x){\displaystyle \frac{1}{2}}\mathrm{}(x)].`$
and substituting it into eq.(16) one finds
$`_{gaugefixing}`$ $`=`$ $`{\displaystyle \frac{1}{2\xi g^2}}[(D\mathrm{}\mathrm{\Delta }\xi +{\displaystyle \frac{i}{2}}M\xi \overline{N}{\displaystyle \frac{i}{2}}\overline{M}\xi N)^2(^\mu v_\mu )^2`$ (20)
$`+`$ $`(\overline{N}i\overline{M}\xi )\mathrm{}(N+iM\xi )+2i(\lambda +{\displaystyle \frac{1}{2}}\overline{M}\xi \chi )\sigma ^\mu _\mu (\overline{\lambda }+{\displaystyle \frac{1}{2}}M\xi \overline{\chi })`$
$``$ $`2(\lambda +{\displaystyle \frac{1}{2}}\overline{M}\xi \chi )\mathrm{}\chi 2(\overline{\lambda }+{\displaystyle \frac{1}{2}}M\xi \overline{\chi })\mathrm{}\overline{\chi }2i\mathrm{}\chi \sigma ^\mu _\mu \overline{\chi }].`$
One can see from eq.(20) that the parameter $`M`$, besides being the gaugino soft mass, plays the role of a mass of the auxiliary field $`\chi `$, while $`\mathrm{\Delta }`$ is the soft mass of the auxiliary fields $`N`$ and $``$. All these fields are unphysical degrees of freedom of the gauge superfield. They are absent in the Wess-Zumino gauge, however, when the gauge-fixing condition is chosen in supersymmetric form (3), this gauge is no longer possible, and the auxiliary fields $`\chi `$, $`N`$, and $``$ survive. Thus, the extra $`\mathrm{\Delta }`$ term is associated with unphysical, ghost, degrees of freedom, just like in the component approach, one has the mass of unphysical $`ϵ`$-scalars. When we go down with energy, all massive fields decouple, and we get the usual nonsupersymmetric Yang-Mills theory.
The $`\mathrm{\Delta }`$-term is renormalized and obeys its own RG equation which can be obtained from the corresponding expression for the gauge coupling via Grassmannian expansion. In due course of renormalization, this term is mixing with the soft masses of scalar superpartners and gives an additional term in RG equations for the latter ($`X`$ term of Ref. mentioned above).
At the end of this section, we would like to comment on the BRST invariance in a softly broken SUSY theory. The BRST transformations (8) due to our choice of normalization of the gauge and ghost fields do not depend on the gauge coupling. Hence, in a softly broken theory they remain unchanged. One can easily check that, despite the substitution $`g^2\stackrel{~}{g}^2`$ and $`\xi \stackrel{~}{\xi }`$, the softly broken SUSY theory remains BRST invariant .
## 3 RG Equations for the Soft Parameters.
Thus, following the procedure described in Refs , to get the RG equations for the soft terms, one has to modify the gauge ($`g_i^2`$) and Yukawa ($`y_{ijk}`$) couplings replacing them by external superfields :
$`\stackrel{~}{g}_i^2`$ $`=`$ $`g_i^2(1+M_i\eta +\overline{M}_i\overline{\eta }+(2M_i\overline{M}_i+\mathrm{\Delta }_i)\eta \overline{\eta }),`$ (21)
$`\stackrel{~}{y}^{ijk}`$ $`=`$ $`y^{ijk}A^{ijk}\eta +{\displaystyle \frac{1}{2}}(y^{njk}(m^2)_n^i+y^{ink}(m^2)_n^j+y^{ijn}(m^2)_n^k)\eta \overline{\eta },`$ (22)
$`\stackrel{~}{\overline{y}}_{ijk}`$ $`=`$ $`\overline{y}_{ijk}\overline{A}_{ijk}\overline{\eta }+{\displaystyle \frac{1}{2}}(y_{njk}(m^2)_i^n+y_{ink}(m^2)_j^n+y_{ijn}(m^2)_k^n)\eta \overline{\eta }.`$ (23)
Then, the $`\beta `$ functions of RG equations for the soft masses of scalar superpartners of the matter fields and for the mass of the auxiliary gauge field are given by
$`[\beta _{m^2}]_j^i=D_2\gamma _j^i,`$ (24)
$`\beta _{\mathrm{\Sigma }_{\alpha _i}}=D_2\gamma _{\alpha _i},`$ (25)
where $`\gamma _j^i`$ and $`\gamma _{\alpha _i}=\beta _{\alpha _i}/\alpha _i`$ are the anomalous dimensions of the matter fields and of the gauge coupling, respectively, and we have introduced the notation
$$\mathrm{\Sigma }_{\alpha _i}=M_i\overline{M}_i+\mathrm{\Delta }_i.$$
The modified expression for the operator $`D_2`$ is
$`D_2`$ $`=`$ $`\overline{D}_1D_1+\mathrm{\Sigma }_{\alpha _i}\alpha _i{\displaystyle \frac{}{\alpha _i}}+{\displaystyle \frac{1}{2}}(m^2)_n^a(y^{nbc}{\displaystyle \frac{}{y^{abc}}}+y^{bnc}{\displaystyle \frac{}{y^{bac}}}+y^{bcn}{\displaystyle \frac{}{y^{bca}}}`$ (26)
$`+y_{abc}{\displaystyle \frac{}{y_{nbc}}}+y_{bac}{\displaystyle \frac{}{y_{bnc}}}+y_{bca}{\displaystyle \frac{}{y_{bcn}}}).`$
It coincides now with that of Ref. with $`X_i=\mathrm{\Delta }_i`$.
To find $`\mathrm{\Sigma }_{\alpha _i}`$, one can use equation (25). In particular, using the expression for the anomalous dimension $`\gamma _\alpha `$ in case of a single non-abelian gauge group calculated up to three loops
$`\gamma _\alpha `$ $`=`$ $`\alpha Q+2\alpha ^2QC(G){\displaystyle \frac{2}{r}}\alpha \gamma _{j}^{i}{}_{}{}^{(1)}C(R)_i^j\alpha ^3Q^2C(G)+4\alpha ^3QC^2(G)`$ (27)
$`{\displaystyle \frac{6}{r}}\alpha ^3QC(R)_j^iC(R)_i^j{\displaystyle \frac{4}{r}}\alpha ^2C(G)\gamma _{j}^{i}{}_{}{}^{(1)}C(R)_i^j+{\displaystyle \frac{3}{r}}\alpha y^{ikm}y_{jkn}\gamma _{m}^{n}{}_{}{}^{(1)}C(R)_i^j`$
$`+{\displaystyle \frac{1}{r}}\alpha \gamma _{j}^{i}{}_{}{}^{(1)}\gamma _{p}^{j}{}_{}{}^{(1)}C(R)_i^p+{\displaystyle \frac{6}{r}}\alpha ^2\gamma _{j}^{i}{}_{}{}^{(1)}C(R)_p^jC(R)_i^p,`$
and the anomalous dimension of the matter field calculated up to two loops
$`\gamma _j^i`$ $`=`$ $`{\displaystyle \frac{1}{2}}y^{ikl}y_{jkl}2\alpha C(R)_j^i`$
$``$ $`(y^{imp}y_{jmn}+2\alpha C(R)_j^p\delta _n^i)({\displaystyle \frac{1}{2}}y^{nkl}y_{pkl}2\alpha C(R)_p^n)+2\alpha ^2QC(R)_j^i,`$
one can get the solution
$`\mathrm{\Sigma }_{\alpha }^{}{}_{}{}^{(1)}`$ $`=`$ $`M^2,`$ (29)
$`\mathrm{\Sigma }_{\alpha }^{}{}_{}{}^{(2)}`$ $`=`$ $`\mathrm{\Delta }_\alpha ^{(2)}=2\alpha [{\displaystyle \frac{1}{r}}(m^2)_j^iC(R)_i^jM^2C(G)],`$ (30)
$`\mathrm{\Sigma }_{\alpha }^{}{}_{}{}^{(3)}`$ $`=`$ $`\mathrm{\Delta }_\alpha ^{(3)}={\displaystyle \frac{\alpha }{2r}}[{\displaystyle \frac{1}{2}}(m^2)_n^iy^{nkl}y_{jkl}+{\displaystyle \frac{1}{2}}(m^2)_j^ny^{ikl}y_{nkl}+2(m^2)_n^my^{ikn}y_{jkm}`$ (31)
$`+A^{ikl}A_{jkl}8\alpha M^2C(R)_j^i]C(R)_i^j2\alpha ^2QC(G)M^2`$
$`4\alpha ^2C(G)[{\displaystyle \frac{1}{r}}(m^2)_j^iC(R)_i^jM^2C(G)].`$
These expressions for $`\mathrm{\Delta }_\alpha `$ (30,31) coincide with those obtained in Ref. for the mass of the $`ϵ`$-scalars.
The nonzero $`\mathrm{\Delta }`$-term modifies the expression for the $`\beta `$ function of the soft scalar mass starting from the second loop. Substituting eq.(30) into the expression for the differential operator $`D_2`$ gives in two loops
$`\left[\beta _{m^2}\right]_j^{i(2)}`$ $`=`$ $`(A^{ikp}A_{jkn}+{\displaystyle \frac{1}{2}}(m^2)_l^iy^{lkp}y_{jkn}+{\displaystyle \frac{1}{2}}y^{ikp}y_{lkn}(m^2)_j^l+{\displaystyle \frac{2}{2}}y^{ilp}(m^2)_l^sy_{jsn}`$ (32)
$`+`$ $`{\displaystyle \frac{1}{2}}y^{iks}(m^2)_s^py_{jkn}+{\displaystyle \frac{1}{2}}y^{ikp}(m^2)_n^sy_{jks}+4\alpha m_A^2C(R)_j^p\delta _n^i)({\displaystyle \frac{1}{2}}y^{nst}y_{pst}2\alpha C(R)_p^n)`$
$``$ $`(y^{ikp}y_{jkn}+2\alpha C(R)_j^p\delta _n^i)({\displaystyle \frac{1}{2}}A^{nst}A_{pst}+{\displaystyle \frac{1}{4}}(m^2)_l^ny^{lst}y_{pst}+{\displaystyle \frac{1}{4}}y^{nst}y_{lst}(m^2)_p^l`$
$`+`$ $`{\displaystyle \frac{4}{4}}y^{nlt}(m^2)_l^sy_{pst}4\alpha m_A^2C(R)_p^n)+12\alpha ^2m_A^2QC(R)^i_j`$
$``$ $`(A^{ikp}y_{jkn}2\alpha m_AC(R)_j^p\delta _n^i)({\displaystyle \frac{1}{2}}y^{nst}A_{pst}+2\alpha m_AC(R)_p^n)`$
$``$ $`(y^{ikp}A_{jkn}2\alpha m_AC(R)_j^p\delta _n^i)({\displaystyle \frac{1}{2}}A^{nst}y_{pst}+2\alpha m_AC(R)_p^n)`$
$`+`$ $`4\alpha ^2C(R)_j^i[{\displaystyle \frac{1}{r}}(m^2)_l^kC(R)_k^lM^2C(G)],`$
where the last term is an extra contribution due to nonzero $`\mathrm{\Delta }_\alpha `$.
To argue that a solution for $`\mathrm{\Delta }_\alpha `$ exists in all orders of PT, one can consider the so-called NSVZ-scheme where the anomalous dimension $`\gamma _\alpha `$ is equal to
$$\gamma _\alpha ^{NSVZ}=\alpha \frac{Q2r^1\mathrm{Tr}[\gamma C(R)]}{12C(G)\alpha }.$$
(33)
Then the solution for $`\mathrm{\Delta }_\alpha `$ to all orders is
$$\mathrm{\Delta }_\alpha ^{NSVZ}=2\alpha \frac{r^1\mathrm{Tr}[m^2C(R)]M^2C(G)}{12C(G)\alpha }.$$
(34)
It coincides with $`X`$ of Ref. .
This problem has been also addressed in Ref. , where originally the additional contribution to the soft term $`\beta `$ function was absent. In a comment to paper it is suggested that the discrepancy can be eliminated by introducing the term proportional to the mass of the $`ϵ`$-scalar in the superfield formalism
$`{\displaystyle \frac{\stackrel{~}{m}_A^2}{2}}V_\mu ^AV_\nu ^A\widehat{g}^{\mu \nu }={\displaystyle \frac{\stackrel{~}{m}_A^2}{2}}{\displaystyle d^4\theta \overline{\eta }\eta \frac{1}{16g^2}\overline{\sigma }_\mu ^{\dot{\alpha }\alpha }\overline{D}_{\dot{\alpha }}(e^{2gV}D_\alpha e^{2gV})\overline{\sigma }_\nu ^{\dot{\beta }\beta }\overline{D}_{\dot{\beta }}(e^{2gV}D_\beta e^{2gV})\widehat{g}^{\mu \nu }},`$ (35)
where $`\widehat{g}^{\mu \nu }`$ is a $`2ϵ`$-dimensional metric tensor.
Similar things were done in Ref. , where the appearance an extra term in RGE for the soft scalar masses is due to additional ”evanescent” operator in DRED scheme as in eq.(35). It leads to additional contribution in higher loops.
However, whenever it is true, technically, it is complicated. We propose here the other solution of this problem.
## 4 Illustration
As an illustration of the described procedure, we consider the case of the MSSM. Here instead of one there are three gauge couplings, and though the recipe is still the same, one faces some problem of the general nature. We obtain below the explicit solutions for the $`\mathrm{\Sigma }_{\alpha _i}`$ terms that can be of interest for the applications in higher loops.
In the MSSM we have three gauge and three Yukawa couplings and, to simplify the formulas, we use the following notation
$$\alpha _i\frac{g_i^2}{16\pi ^2},i=1,2,3;Y_k\frac{y_k^2}{16\pi ^2},k=t,b,\tau .$$
Then, the modified couplings (21-23) take the form
$`\stackrel{~}{\alpha }_i`$ $`=`$ $`\alpha _i\left(1+M_i\eta +\overline{M}_i\overline{\eta }+(M_i\overline{M}_i+\mathrm{\Sigma }_{\alpha _i})\eta \overline{\eta }\right),`$ (36)
$`\stackrel{~}{Y}_k`$ $`=`$ $`Y_k\left(1A_k\eta \overline{A}_k\overline{\eta }+(A_k\overline{A}_k+\mathrm{\Sigma }_k)\eta \overline{\eta }\right),`$ (37)
where $`\mathrm{\Sigma }_k`$ is the sum of the soft masses squared corresponding to a given Yukawa vertex
$$\mathrm{\Sigma }_t=\stackrel{~}{m}_Q^2+\stackrel{~}{m}_U^2+m_{H_2}^2,\mathrm{\Sigma }_b=\stackrel{~}{m}_Q^2+\stackrel{~}{m}_D^2+m_{H_1}^2,\mathrm{\Sigma }_\tau =\stackrel{~}{m}_L^2+\stackrel{~}{m}_E^2+m_{H_1}^2.$$
Now the RG equations for a rigid theory can be written in a universal form
$$\dot{a}_i=a_i\gamma _i(a),a_i=\{\alpha _i,Y_k\},$$
(38)
where $`\gamma _i(a)`$ stands for a sum of corresponding anomalous dimensions. In the same notation, the soft terms (36,37) take the form
$$\stackrel{~}{a}_i=a_i(1+m_i\eta +\overline{m}_i\overline{\eta }+S_i\eta \overline{\eta }),$$
(39)
where $`m_i=\{M_i,A_k\}`$ and $`S_i=\{M_i\overline{M}_i+\mathrm{\Sigma }_{\alpha _i},A_k\overline{A}_k+\mathrm{\Sigma }_k\}`$.
Substituting eq.(39) into eq.(38) and expanding over $`\eta `$ and $`\overline{\eta }`$, one can get the RG equations for the soft terms
$$\dot{\stackrel{~}{a}}_i=\stackrel{~}{a}_i\gamma _i(\stackrel{~}{a}).$$
(40)
Consider first the F-terms. Expanding over $`\eta `$, one has
$$\dot{m}_i=\gamma _i(\stackrel{~}{a})|_F=\underset{j}{}a_j\frac{\gamma _i}{a_j}m_j.$$
(41)
This is just the RG equation for the soft terms $`M_i`$ and $`A_k`$ . Proceeding in the same way for the D-terms, one gets after some algebra
$$\dot{S}_i=2m_i\underset{j}{}a_j\frac{\gamma _i}{a_j}m_j+\underset{j}{}a_j\frac{\gamma _i}{a_j}S_j+\underset{j,k}{}a_ja_k\frac{^2\gamma _i}{a_ja_k}m_jm_k.$$
(42)
Substituting $`S_i=m_i\overline{m}_i+\mathrm{\Sigma }_i`$, one has the following RG equation for the mass terms
$$\dot{\mathrm{\Sigma }}_i=\gamma _i(\stackrel{~}{a})|_D=\underset{j}{}a_j\frac{\gamma _i}{a_j}(m_jm_j+\mathrm{\Sigma }_j)+\underset{j,k}{}a_ja_k\frac{^2\gamma _i}{a_ja_k}m_jm_k.$$
(43)
Using the explicit form of anomalous dimensions calculated up to some order, one can reproduce the desired RG equations for the soft terms. In case of squark and slepton masses, they contain the contributions from unphysical masses $`\mathrm{\Sigma }_{\alpha _i}`$. To eliminate them, one has to solve the equation for $`\mathrm{\Sigma }_{\alpha _i}`$. In the case of the MSSM up to three loops, the solutions are
$`\mathrm{\Sigma }_{\alpha _1}`$ $`=`$ $`M_1^2\alpha _1\sigma _1{\displaystyle \frac{199}{25}}\alpha _1^2M_1^2{\displaystyle \frac{27}{5}}\alpha _1\alpha _2M_2^2{\displaystyle \frac{88}{5}}\alpha _1\alpha _3M_3^2`$ (44)
$`+`$ $`{\displaystyle \frac{13}{5}}\alpha _1Y_t(\mathrm{\Sigma }_t+A_t^2)+{\displaystyle \frac{7}{5}}\alpha _1Y_b(\mathrm{\Sigma }_b+A_b^2)+{\displaystyle \frac{9}{5}}\alpha _1Y_\tau (\mathrm{\Sigma }_\tau +A_\tau ^2),`$
$`\mathrm{\Sigma }_{\alpha _2}`$ $`=`$ $`M_2^2\alpha _2(\sigma _24M_2^2)\alpha _2^2(4\sigma _2+9M_2^2){\displaystyle \frac{9}{5}}\alpha _2\alpha _1M_1^224\alpha _2\alpha _3M_3^2`$ (45)
$`+`$ $`3\alpha _2Y_t(\mathrm{\Sigma }_t+A_t^2)+3\alpha _2Y_b(\mathrm{\Sigma }_b+A_b^2)+\alpha _2Y_\tau (\mathrm{\Sigma }_\tau +A_\tau ^2),`$
$`\mathrm{\Sigma }_{\alpha _3}`$ $`=`$ $`M_3^2\alpha _3(\sigma _36M_3^2)\alpha _3^2(6\sigma _322M_3^2){\displaystyle \frac{11}{5}}\alpha _3\alpha _1M_1^29\alpha _3\alpha _2M_2^2`$ (46)
$`+`$ $`2\alpha _3Y_t(\mathrm{\Sigma }_t+A_t^2)+2\alpha _3Y_b(\mathrm{\Sigma }_b+A_b^2),`$
where we have used the combinations
$`\sigma _1`$ $`=`$ $`{\displaystyle \frac{1}{5}}\left[3(m_{H_1}^2+m_{H_2}^2)+3(\stackrel{~}{m}_Q^2+3\stackrel{~}{m}_L^2+8\stackrel{~}{m}_U^2+2\stackrel{~}{m}_D^2+6\stackrel{~}{m}_E^2)\right],`$ (47)
$`\sigma _2`$ $`=`$ $`m_{H_1}^2+m_{H_2}^2+3(3\stackrel{~}{m}_Q^2+\stackrel{~}{m}_L^2),`$ (48)
$`\sigma _3`$ $`=`$ $`3(2\stackrel{~}{m}_Q^2+\stackrel{~}{m}_U^2+\stackrel{~}{m}_D^2).`$ (49)
Notice, however, that the solutions (47-49) correspond to particular boundary conditions, while, in general, one can use arbitrary ones. Here we encounter the general problem that the solutions for physical masses depend on the unphysical parameter ($`ϵ`$-scalar mass in the component approach in the DRED scheme and an auxiliary field mass $`\mathrm{\Delta }`$ in the superfield approach).
The solution to this paradox, mentioned also in Ref. , follows from the observation that the running soft masses that obey the RG equations are not, strictly speaking, the observables and are scheme-dependent. More appropriate are the pole masses, that are scheme-independent. The authors of Ref. proposed the solution of the paradox by passing to the DRED scheme via the shift of the running soft mass, which allows one to get rid of the unwanted $`ϵ`$-scalar mass and does not influence the pole mass. In one-loop order, the shift is
$$(m^2)_i^j|_{\overline{DR}^{^{}}}=(m^2)_i^j|_{\overline{DR}}\frac{2g_A^2C_A(i)}{(4\pi )^2}\delta _i^j\stackrel{~}{m}_ϵ^2,$$
(50)
where $`\stackrel{~}{m}_ϵ`$ is the $`ϵ`$-scalar mass. The procedure can be continued in the same way in higher loops.
One can easily see how a similar trick works in our approach in case of one gauge coupling (and, consequently, one $`\mathrm{\Delta }`$ term). Indeed, consider eq.(43). It is a linear inhomogeneous differential equation. Hence, to any given solution of this equation one can add an arbitrary solution of a homogeneous equation. In our case, the solution of a homogeneous equation is
$$\mathrm{\Sigma }_i=𝒞\gamma _i,i=\alpha _1,\alpha _2,\alpha _3,t,b,\tau ,$$
(51)
where $`𝒞`$ is an overall constant.
Hence, if one has the only gauge coupling one can choose the constant $`𝒞`$ so that one can get any desirable boundary condition for $`\mathrm{\Sigma }_\alpha `$. The price for this is extra terms in the other $`\mathrm{\Sigma }`$’s (and soft masses) proportional to the corresponding anomalous dimensions. However, the shift of the running mass by a term proportional to the anomalous dimension does not change the pole mass, since it can be absorbed into the scale redefinition. This is due to the fact that the coefficient of the $`\mathrm{log}\mu ^2`$ term is just the anomalous dimension of the field.
Thus, the arbitrariness in the unphysical mass boundary condition does not influence the physical masses.
However, one has only one overall constant $`𝒞`$, and the above argument clearly works when one has only one gauge coupling. In case of many couplings, it is more tricky, and we have not found an obvious explanation.
## 5 Conclusion
Summarizing, we would like to stress once again that soft breaking of supersymmetry can be realized via interaction with an external superfield that develops nonzero v.e.v.’s for its $`F`$ and $`D`$ components. In the superfield notation, it can be reformulated as a modification of the rigid couplings that become external superfields. The same is true for the gauge-fixing parameter that can also be considered as a rigid coupling. The soft masses of scalar particles obtain their contribution from the D-components of external superfields. The latter also lead to nonzero masses for unphysical degrees of freedom, ghost and gauge auxiliary fields. These unphysical masses enter into the RG equations for the physical scalars and have to be eliminated. This creates an ambiguity in the running scalar masses; it can be resolved by passing to the pole masses.
Acknowledgments
The authors are grateful to A.Bakulev, S.Mikhailov, G.Moultaka, A.Onishchenko and A.Slavnov for useful discussions. Financial support from RFBR (grants # 99-02-16650 and 00-15-96691) is kindly acknowledged.
## Appendix A
The RG equation for the parameter $`\xi `$ in a rigid theory is
$$\dot{\xi }=\gamma _V\xi ,$$
(A.1)
where $`\gamma _V`$ is the anomalous dimension of the gauge superfield. To find the soft terms $`x,\overline{x}`$ and $`z`$, one should solve the modified equation
$$\dot{\stackrel{~}{\xi }}=\gamma _V(\stackrel{~}{\alpha },\stackrel{~}{y},\stackrel{~}{\xi })\stackrel{~}{\xi }.$$
(A.2)
In one-loop order $`\gamma _V=(b_1+b_2\xi )\alpha `$, where $`b_1+b_2=Q`$, and the solutions are
$`x`$ $`=`$ $`(M+x_0){\displaystyle \frac{b_1+b_2\xi }{Q}},\overline{x}=(\overline{M}+\overline{x}_0){\displaystyle \frac{b_1+b_2\xi }{Q}},`$ (A.3)
$`z`$ $`=`$ $`(\mathrm{\Sigma }_\alpha +z_0){\displaystyle \frac{b_1+b_2\xi }{Q}}+{\displaystyle \frac{b_2\xi }{Q}}(M+x_0)(\overline{M}+\overline{x}_0){\displaystyle \frac{b_1+b_2\xi }{Q}},`$ (A.4)
where $`x_0,\overline{x}_0`$, and $`z_0`$ are arbitrary constants. In the Abelian case when $`b_1=Q,b_2=0`$, the solutions are simplified and can be chosen as
$$x=M(1\xi ),\overline{x}=\overline{M}(1\xi ),z=\mathrm{\Sigma }_\alpha (1\xi )M\overline{M}\xi (1\xi ).$$
Together with the expression for $`\stackrel{~}{\alpha }`$ (13) it gives eq.(2) above.
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# Incidence relations among the Schubert cells of equivariant Hilbert schemes
## 1 Introduction
The problem (general wording).
Fixing an algebraically closed field $`k`$, the Hilbert scheme $``$ parametrizing the zero dimensional subschemes of the affine plane $`Speck[x,y]`$ admits a natural action of the two dimensional torus $`k^{}\times k^{}`$, induced by the linear action $`(t_1,t_2).x^\alpha y^\beta =(t_1.x)^\alpha (t_2.y)^\beta `$ of $`k^{}\times k^{}`$ on $`k[x,y]`$. If $`T_{ab}=\{(t^b,t^a),tk^{}\}`$ is a one dimensional sub-torus, the closed subscheme $`_{ab}`$ which parametrizes by definition the zero dimensional subschemes invariant under the action of $`T_{ab}`$ is a disjoint union of irreducible subschemes $`_{ab}(H)`$. Each $`_{ab}(H)`$ can be embedded in a product of Grassmannians and inherits the stratification whose cells are the inverse images of the products of the Schubert cells. The problem is to understand the geometry of this stratification, and in particular to describe the incidences between the cells.
Motivations
Before explaining the results and the techniques, let’s explain the motivations. Roughly speaking, Grassmannians are easier to study than Hilbert schemes because they are stratified by Schubert cells. Those stratifications enable for instance to compute the intersection rings of Grassmannians, the Hilbert function of their natural embeddings in a projective space (\[GH\],\[M\]). On the other hand, Hilbert schemes admit a stratification furnished by the theory of Grobner bases (or standard basis) which is fussy to describe (\[Gr\]). The equivariant Hilbert scheme $`_{ab}(H)`$ is an object which both looks like a Hilbert scheme and a Grassmannian, a sort of interface between them. It has two stratifications, the stratification in Schubert cells introduced above and the Gröbner stratification obtained by restricting the stratification of $``$ to $`_{ab}(H)`$. Moreover these two stratifications coincide. So the philosophy is to try to describe the geometry of $`_{ab}=_{ab}(H)`$, and then to lift the information on $``$. This is our motivation to study $`_{ab}(H)`$. To see some precise examples where the geometry of $`_{ab}`$ can be used to describe the geometry of $``$, see $`[I]`$, where structures of bundles are highlighted, \[ES\] where the Betti numbers of the connected components of $``$ are computed via the study of the equivariant inclusion $`_{ab}`$ for $`(a,b)`$ general enough, or \[B\] which describes links between the equivariant cohomology of $`_{ab}`$ and the equivariant cohomology of $``$. Beyond this motivation, note that $`_{ab}`$ is a classical object of study (\[Gö\],\[I\],\[Y\] for instance).
Known results and precise formulation of the problem.
For some special Hilbert functions $`H`$, the equivariant Hilbert scheme $`_{ab}(H)`$ is a Grassmannian and the cells of $`_{ab}(H)`$ are the Schubert cells (for this reason, we will say a “Schubert cell” to designate a cell on $`_{ab}(H)`$ ). In this simple case, the closure of a cell is a union of cells and there is an explicit numerical condition to check incidence. This is no longer so simple when $`_{ab}(H)`$ is not a Grassmannian: Yameogo has given in \[Y\] an example where the closure of a cell is not a union of cells. So the problem splits into two different pieces: the “weak incidence” problem which consists in deciding whether the closure of a cell meets another cell, and the “strong incidence” problem which consists in deciding whether the closure of a cell contains another cell. In \[Y\], Yameogo faced the weak incidence problem and gave a necessary but not sufficient condition to have a relation $`\overline{C}C^{}\mathrm{}`$. In this paper, we address the weak incidence problem too.
The results.
Our main answer will be a necessary condition (theorem 8) for weak incidence. No counter-example to sufficiency is known to the author. The other propositions try to analyse the pertinence of the condition, by comparing it to the previous results (prop. 21, 22) and by testing it on small Hilbert functions (section 5).
We now formulate our results with more details. We associate with any Schubert cell $`C`$ a combinatorial data and we express the relation $`\overline{C}C^{}\mathrm{}`$ by a combinatorial property linking $`C`$ and $`C^{}`$, as follows.
Each Schubert cell is determined by a staircase $`E`$, i.e. a subset of $`^2`$ whose complementary is stable under the addition of $`^2`$: the ideal $`I^E`$ generated by the monomials whose exponent is not in $`E`$ is the unique monomial ideal of the cell. If $`n`$ is the cardinal of $`E`$, $`E`$ is included in a box of size $`n\times n`$ and we denote by $`E^\nu `$ the dual staircase defined by the complementary of $`E`$ in the box (see figure). A system of arrows $`S`$ on $`E`$ is the data for each element $`p`$ of $`E`$ of an arrow $`(p,f_p)`$ with origin $`p`$, the set of arrows being compelled to compatibility conditions.
Replacing each element $`p`$ of $`E`$ by $`f_p`$ gives a new subset of $`^2`$ which turns out to be a staircase thanks to the compatibility conditions. Thinking of $`S`$ as an operator on $`E`$, we denote by $`S(E)`$ the staircase obtained from $`E`$ and $`S`$ by the above procedure. The theorem says that two cells $`C(E)`$ and $`C(F)`$ corresponding respectively to staircases $`E`$ and $`F`$ verify $`\overline{C(E)}C(F)\mathrm{}`$ only if
1. there exists a system of arrows $`S`$ on $`E`$ such that $`S(E)=F`$
2. there exists a system of arrows $`S^\nu `$ on $`E^\nu `$ such that $`S^\nu (E^\nu )=F^\nu `$
We then compare this condition to the known conditions and we test it on an example. In the particular case of Grassmannians, 1) and 2) are equivalent and reduce to the well known necessary and sufficient incidence condition (prop. 21). In general, they are not equivalent (ex. 10) but any of them imply the condition of \[Y\] (proposition 22). Finally, we study an example (section 5) which illustrates that when the Hilbert function $`H`$ is small, we can solve the weak incidence problem because our condition can be shown to be necessary and sufficient.
The methods.
The first step consists in putting the problem in the context of Gröbner basis theory and to describe a cell $`C(E)`$ as the locus formed by the ideals $`I`$ having $`I^E`$ as initial ideal. Then we interpret the incidence relation $`\overline{C(E)}C(F)\mathrm{}`$ as the existence of a rational morphism from $`𝔸^1=Speck[t]`$ to $`C(E)`$ which can be extended in $`\mathrm{}`$ by putting $`f(\mathrm{})=X_FC(F)`$. The universal ideal on $`_{ab}(H)`$ restricted on the affine curve defines an ideal $`I(t)k[x,y,t]`$. The key consists in using ideas coming from Grobner basis theory to exhibit within $`I(t)`$ a set of elements $`P_1,P_2,\mathrm{}`$ from which we read the system of arrows on $`E^\nu `$. Then, using an argument of duality corresponding to the geometric notion of linkage of two zero dimensional schemes in a complete intersection, we get the system of arrows on $`E`$.
I thank Michel Brion who told me about the problems that motivated this work and who answered several questions.
## 2 The stratification on $`_{ab}(H)`$.
### 2.1 Notations
First, we keep the notations of the introduction: $`k`$ is an algebraically closed field, $`a`$ and $`b`$ are two relatively prime integers, and $`T_{a,b}=\{(t^b,t^a),tk^{}\}`$ is a one dimensional subtorus of $`k^{}\times k^{}=:T`$. The torus $`T`$ acts on the Hilbert scheme $``$ parametrizing the zero-dimensional subschemes of $`Speck[x,y]`$, and $`_{ab}`$ is the subscheme of $``$ parametrizing the subschemes invariant under the action of $`T_{ab}T`$. Alternatively, the subschemes of $`_{ab}`$ can be characterized by their ideals using degrees. The degrees $`d,d_x,d_y`$ are defined on monomials by $`d(x^\alpha y^\beta )=b\alpha +a\beta ,d_x(x^\alpha y^\beta )=\alpha ,d_y(x^\alpha y^\beta )=\beta `$. If $`I`$ is an ideal of $`k[x,y]`$, we let $`I_n:=Ik[x,y]_n`$, where $`k[x,y]_n`$ denotes the vector space generated by the monomials $`m`$ of degree $`d(m)=n`$. A subscheme $`Z`$ is in $`_{ab}`$, if and only if its ideal is quasi-homogeneous with respect to $`d`$, ie. $`I(Z)=_{n0}I(Z)_n`$.
We order the monomials of $`k[x,y]`$ by the rule $`m_1<m_2`$ if $`d(m_1)<d(m_2)`$ or ($`d(m_1)=d(m_2)`$ and $`d_y(m_1)<d_y(m_2)`$).
The scheme $`_{ab}`$ is not connected but the connected components are determined by a Hilbert function. By semi-continuity, if a subscheme $`Z^{}_{ab}`$ is a specialization of $`Z_{ab}`$, then the codimensions of $`I_n(Z)`$ and $`I_n(Z^{})`$ in $`k[x,y]_n`$ verify $`codimI_n(Z)codimI_n(Z^{})`$. But $`Z`$ and $`Z^{}`$ have the same length $`l=_{n0}codimI_n(Z)=_{n0}codimI_n(Z^{})`$. It follows that the sequence $`H(Z)=(h_0,h_1,h_2\mathrm{})`$ where $`h_i=codimI_n(Z)`$ is constant on the connected components of $`_{ab}`$ and that $`h_n=0`$ for $`n`$ big enough. If $`H=(h_0,\mathrm{},h_r,0,0,\mathrm{})`$ is any sequence, we denote by $`_{ab}(H)`$ the closed subscheme of $``$ parametrizing the schemes $`Z`$ verifying $`H(Z)=H`$. One can verify that $`_{ab}(H)`$ is an irreducible connected component of $`_{ab}`$ (though we won’t use it in the sequel).
Recall that a staircase is a subset of $`^2`$ whose complementary is stable by addition of $`^2`$. Staircases will be used to parametrize the stratas on $`_{ab}(H)`$. In this paper, we will identify freely the monomial $`x^py^q`$ with the couple $`(p,q)`$ and therefore the expression “staircase of monomials” will make sense. More generally, we will transpose unscrupulously the definitions between couples of integers and monomials. If $`E`$ is a staircase, then the vector space $`I^E`$ generated by the monomials which are not in $`E`$ is an ideal and reciprocally, every monomial ideal is an ideal $`I^E`$ for a unique staircase $`E`$. The subscheme $`Z(E)`$ whose ideal is $`I^E`$ is in $`_{ab}(H)`$ if and only if $`E`$ has $`h_i`$ elements in degree $`i`$.
### 2.2 The possible definitions
In this section, we explain that there are three ways to describe the stratification on $`_{ab}(H)`$. The fact that the definitions coincide is shown in \[Y\]. We call a cell of this stratification a “Schubert cell” on $`_{ab}(H)`$ as it is a Schubert cell when $`_{ab}(H)`$ is a Grassmannian (section 4.1). Moreover each cell contains a unique subscheme $`Z(E)`$ and we will denote this cell by $`C(E)`$.
The Grobner basis point of view. The theory of Gröbner basis associates with every ideal in $`k[x,y]`$ a monomial ideal (with respect to the monomial order chosen in section 2.1) called initial ideal and $`C(E)`$ is the locus in $`_{ab}(H)`$ parametrizing the ideals whose initial ideal is $`I^E`$. For the reader not familiar with Gröbner basis, we give a characterization which is sufficient for the sequel. Let $`m_1,m_2,\mathrm{}`$ be the monomials which don’t belong to $`E`$. An ideal of $`_{ab}(H)`$ is in $`C(E)`$ if, regarding it as a $`k`$-vector space, it admits a base $`f_1,f_2,\mathrm{}`$ where $`f_i=m_i+R_i`$, $`R_i`$ being a linear combination of monomials strictly smaller than $`m_i`$.
###### Remark 1.
When the product $`a.b`$ is zero, there is at most one staircase compatible with the Hilbert function $`H`$ (i.e. such that $`E`$ has $`h_i`$ elements in degree $`i`$) and $`_{ab}(H)`$ is either empty or is reduced to a unique cell (it is non empty when $`H`$ is a decreasing sequence). These cases are not relevant for the incidence problem and we suppose from now on $`ab0`$ and $`a>0`$.
The Schubert cells point of view. Let $`H=(h_0,\mathrm{},h_r,0,0,\mathrm{})`$, and $`Z_{ab}(H)`$. The ideal $`I_n(Z)`$ being a vector space of codimension $`h_n`$ in $`k[x,y]_n`$, it corresponds to a point $`p_n`$ in a Grassmannian $`G_n`$. So $`I(Z)=I_n(Z)`$ corresponds to a point $`p=(p_0,\mathrm{},p_r)`$ in the product of Grassmannianns $`G_0\times G_1\mathrm{}\times G_r`$. This set theoretical representation turns out to be a closed embedding $`g:_{ab}(H)G_0\times G_1\mathrm{}\times G_r`$. Let $`V_i`$ be the subspace generated by the $`i`$ smallest monomials of $`k[x,y]_n`$. The flag $`V_0V_1\mathrm{}V_{n+1}`$ defines by a classical construction Schubert cells on $`G_n`$ which stratify it. The stratification we consider on $`_{ab}(H)`$ is the stratification whose strata are the locally closed subschemes $`g^1(C_0\times C_1\times \mathrm{}\times C_r)`$, where $`C_i`$ is a Schubert cell on $`G_i`$.
The Bialynicki-Birula point of view. Let $`X`$ be a smooth projective variety over $`k`$ admitting an action of the torus $`k^{}`$. Suppose that the action has a finite number of fixed points $`x_1,\mathrm{},x_n`$. Let $`T_{X,x_i}^+`$ be the part of the tangent space to $`x_i`$ in $`X`$ where the weights of the $`k^{}`$-action are positive, and let $`X_i:=\{xX,lim_{t0}(t.x)=x_i\}`$. Then a theorem of Bialynicki-Birula asserts that the $`X_i`$ are a cellular decomposition of $`X`$ in affine spaces and satisfy $`T_{X_i,x_i}=T_{X,x_i}^+`$. In our case, fixing two integers $`p`$ and $`q`$ with $`ap+bq>0`$, the torus $`k^{}`$ acts on $`k[x,y]`$ by $`t.x=t^p.x`$ and $`t.y=t^qy`$. This action induces an action of $`k^{}`$ on $`_{ab}(H)`$. The fixed points of $`_{ab}(H)`$ under $`k^{}`$ are the monomial subschemes $`Z(E)`$. Applying the Bialynicki-Birula theorem to the action of $`k^{}`$ on $`_{ab}(H)`$, we get a stratification and $`C(E)`$ is the cell associated with the fixed point $`Z(E)`$.
### 2.3 Incidence relations
In this section, we recall the main theorem of \[Y\] about incidence relations.
The monomials of $`k[x,y]`$ can be ordered in an infinite sequence $`m_0<m_1<\mathrm{}`$ thanks to the monomial order. Let $`S_E`$ be the function from $``$ to $``$ defined by $`S_E(k)=`$number of monomials in $`E`$ smaller or equal to $`m_k`$.
###### Theorem 2.
If $`\overline{C(E)}C(F)\mathrm{}`$ then $`S_ES_F`$.
###### Remark 3.
In \[Y\], the result says $``$ instead of $``$ as above. The reason is that we have slightly changed the definition of $`S_E`$ to get a shorter presentation. Moreover, the paper deals with the homogeneous case $`a=1,b=1`$ but the extension is immediate.
The condition is not sufficient: if $`E`$ and $`F`$ are the staircases whose monomial ideals are $`I^E=(y^4,xy^2,x^2y,x^5)`$ and $`I^F=(y^5,xy^2,x^3)`$ then $`S_ES_F`$ but the incidence $`\overline{C(E)}C(F)\mathrm{}`$ is not fulfilled(\[Y\]).
## 3 A necessary condition for weak incidence
### 3.1 Statement of the result
In this section, we give a necessary condition on two staircases $`E`$ and $`F`$ to fulfill the relation $`\overline{C(E)}C(F)\mathrm{}`$. The condition will rely on the notions of system of arrows and of dual staircase that we introduce now.
###### Definition 4.
An arrow on $`E`$ is a couple of monomials $`(P,Q)`$, ($`P`$ being the origin of the arrow and $`Q`$ being the end of the arrow), such that $`P`$ is in $`E`$ and such that the vector $`\stackrel{}{PQ}`$ of $`^2`$ is a negative multiple of $`(a,b)`$ (recall that we have adopted the convention $`a>0`$). If $`f=(P,Q)`$, the multiplication by a monomial $`m`$ of $`f`$ is the arrow $`m.f:=(mP,mQ)`$. An arrow $`f=(P,Q)`$ is shorter that $`f^{}=(P^{},Q^{})`$ if $`\stackrel{}{PQ}=\lambda (a,b)`$, $`\stackrel{}{P^{}Q^{}}=\lambda ^{}(a,b)`$, and if the absolute values of $`\lambda `$ and $`\lambda ^{}`$ verify $`|\lambda ||\lambda ^{}|`$.
A system of arrows on $`E`$ is the data for each element $`pE`$ of an arrow $`(p,f_p)`$ on $`E`$ such that set $`S`$ of arrows parametrized by $`E`$ satisfy the following conditions:
* $`pqf_pf_q`$
* the arrows are compatible with division, which means that $`f=(P,Q)S`$, $`Q^{}`$ monomial of $`k[x,y]`$ dividing $`Q`$, $`gS`$ such that $`g`$ has end $`Q^{}`$ and such that $`g`$ is shorter than $`f`$.
With a staircase $`E`$ and a system of arrows $`S`$, one can define a new staircase $`S(E)`$ as follows.
###### Proposition 5.
Let $`E`$ be a staircase and $`S`$ a system of arrows on $`E`$. Let $`S(E)`$ be the set of monomials which are the end of an arrow. Then $`S(E)`$ is a staircase.
Proof: We have to verify that if $`m`$ and $`Q^{}`$ are two monomials with $`mQ^{}S(E)`$ then $`Q^{}S(E)`$. By definition of $`S(E)`$, $`Q:=mQ^{}`$ is the end of an arrow. The compatibility of $`S`$ with division shows that $`Q^{}S(E)`$.
###### Definition 6.
If $`E`$ and $`F`$ are two staircases, and if $`S`$ is a system of arrows on $`E`$ such that $`F=S(E)`$, we will say that $`S`$ is a system of arrows from $`E`$ to $`F`$.
###### Definition 7.
The box of size $`M\times N`$ is $`B_{M\times N}:=\{(x,y)^2,x<M,y<N\}`$. If a staircase $`E`$ is included in $`B_{M\times N}`$, the dual of $`E`$ in the box is by definition the set $`E_{MN}^\nu :=\{(x,y)^2suchthat(e_1,e_2)E,e_1+x<M1\text{ and }e_2+y<N1\}`$. In the sequel, we will often write $`E^\nu `$ instead of $`E_{MN}^\nu `$ for simplicity.
###### Theorem 8.
Let $`E`$ and $`F`$ be two staircases included in a box $`B_{M\times N}`$. Let $`E^\nu `$ and $`F^\nu `$ be their respective dual in $`B_{M\times N}`$. If the incidence $`\overline{C(E)}C(F)\mathrm{}`$ is fulfilled, then
1. there exists a system of arrows $`S`$ on $`E`$ such that $`S(E)=F`$
2. there exists a system of arrows $`S^\nu `$ on $`E^\nu `$ such that $`S^\nu (E^\nu )=F^\nu `$
Proof: see section 3.2
###### Remark 9.
One can prove that the existence of $`S^\nu `$ does not depend on the choice of the box $`B_{M\times N}`$ in which we construct the dual staircase $`E^\nu `$.
The following example shows that conditions 1) and 2) are not equivalent.
###### Example 10.
Let $`E`$ and $`F`$ be the staircases whose monomial ideals are $`I^E=(y^4,xy^2,x^2y,x^5)`$ and $`I^F=(y^5,xy^2,x^3)`$. Let $`(a,b)=(1,1)`$. There exist a system $`S`$ on $`E`$ such that $`S(E)=F`$ but there is no system of arrows $`S`$ on $`E_{\mathrm{5\; 5}}^\nu `$ such that $`S^\nu (E_{\mathrm{5\; 5}}^\nu )=F_{\mathrm{5\; 5}}^\nu `$.
The set $`S:=\{(x^3,x^2y),(x^4,y^4),(p,p),pE,px^3,px^4\}`$ is a system of arrows on $`E`$ such that $`S(E)=F`$. The dual staircases $`E^\nu `$ and $`F^\nu `$ in $`B_{5\times 5}`$ are such that $`I^{E^\nu }=(y^4,x^3y^3,x^4y,x^5)`$ and $`I^{F^\nu }=(y^5,x^2y^3,x^4)`$. Suppose that there exists a system of arrows $`S^\nu `$ on $`E^\nu `$ such that $`S^\nu (E^\nu )=F^\nu `$. Because $`x^2y^3`$ is in $`E^\nu `$ but not in $`F^\nu `$, $`S^\nu `$ must contain an arrow $`f`$ from $`x^2y^3`$ to a point $`p>x^2y^3`$, $`pF^\nu `$ ie. $`p=xy^4`$. The monomial $`x^3y^2F^\nu `$ so there is an arrow $`f^{}`$ from a monomial $`pE^\nu `$ and $`x^3y^2`$ to $`x^3y^2`$, i.e. $`p=x^3y^2`$. The compatibility with division and the arrow $`f^{}`$ show that all monomials dividing $`x^3y^2`$ admit an arrow to themselves. The compatibility with division by $`y`$ applied to the arrow $`f`$ ensures the existence of an arrow from $`xy^3`$ to itself. It follows that the compatibility with division by $`x`$ applied to $`f`$ cannot be satisfied.
### 3.2 Proof of theorem 8
#### 3.2.1 Proof of point 2
To show that there exists a system $`S^\nu `$ on $`E^\nu `$ such that $`S^\nu (E^\nu )=F^\nu `$, we will exhibit a system $`S^c`$ on the complementary $`E^c`$ of $`E`$ in $`^2`$ such that $`S^c(E^c)=F^c`$. So the first step is to explain what is a system on $`E^c`$, and to reduce the proof to a construction of such a system $`S^c`$. The reduction is given in proposition 14
###### Definition-Proposition 11.
An arrow on $`E^c`$ is a couple of monomials $`(P,Q)`$ such that $`PE`$ and such that the vector $`\stackrel{}{PQ}`$ of $`^2`$ is a positive multiple of $`(a,b)`$. A system of arrows on $`E^c`$ is the data for each element $`pE`$ of an arrow $`(p,f_p)`$ on $`E^c`$ such that the set $`S^c`$ of arrows parametrized by $`E^c`$ satisfy the following conditions:
* $`pqf_pf_q`$
* the arrows are compatible with multiplication, which means that $`f=(P,Q)S^c`$, $`m`$ monomial of $`k[x,y]`$, $`gS^c`$ such that $`g`$ has end the product $`mQ`$ and such that $`g`$ is shorter than $`f`$.
If $`S^c`$ is a system of arrows on $`E^c`$, then the set
$$S^c(E^c):=\{(x,y)\text{ which are the end of an arrow of }S^c\}$$
is the complementary of a staircase $`F`$. The system $`S^c`$ will be called a system of arrows from $`E^c`$ to $`F^c`$.
The proof is similar to the proof of proposition 5
The next lemma identifies systems on $`E^c`$ and systems on $`E^\nu `$. Let $`S^\nu `$ (resp. $`S^c`$) be the set of systems of arrows on $`E^\nu `$ (resp. on $`E^c`$). Let $`\phi :B_{M\times N}B_{M\times N}`$, $`(x,y)(M1x,N1y)`$ be the dualizing map. For an arrow $`f=(p,q)`$ on $`E^\nu `$, if $`qB_{M\times N}`$, we define $`\phi (f):=(\phi (p),\phi (q))`$. Let
$$S_{MN}^\nu :=\{S^\nu \text{ system of arrows on }E^\nu \text{ s.t. }s=(p,q)S^\nu qB_{M\times N}\}.$$
If $`S^\nu S_{MN}^\nu `$, let $`\psi (S^\nu ):=\{\phi (f),fS^\nu \}\{(p,p),pE^c\phi (E^\nu )\}`$.
###### Lemma 12.
The map $`\psi :S_{MN}^\nu S^c`$ is well defined, (ie. $`\psi (S^\nu )`$ is a system of arrows on $`E^c`$) and injective. The image $`Im\psi `$ contains the systems $`S^c`$ such that: for all monomial $`mB_{M\times N}`$, the arrow of $`S^c`$ whose origin is $`m`$ is the arrow $`(m,m)`$. Moreover if $`S^\nu `$ and $`S^c`$ are two systems such that $`\psi (S^\nu )=S^c`$, then $`S^\nu (E^\nu )=(S^c(E^c)^c)^\nu `$.
Proof: it consists in a sequence of easy combinatorial verifications which are left to the reader.
Here is a another criterion to check whether a system $`S^c`$ on $`E^c`$ is in $`Im\psi `$.
###### Lemma 13.
Let $`E`$ be a staircase included in the box $`B_{M\times N}`$, $`S^c`$ be a system of arrows on $`E^c`$, and $`F`$ be the staircase such that $`S^c(E^c)=F^c`$. If $`FB_{M\times N}`$, then $`S^cIm\psi `$.
Proof: Let $`p_0,\mathrm{},p_s`$ be the set of monomials of $`k[x,y]_s`$. Up to reordering, one can suppose that $`p_0>p_1>\mathrm{}>p_s`$. There exist integers $`l`$ and $`m`$ such that $`k[x,y]_sB_{M\times N}=\{p_0,\mathrm{},p_l\}\{p_m,\mathrm{},p_s\}`$, $`l<m`$, $`1ls`$, $`0ms+1`$. By 12, to prove the proposition, we have to show that $`i\{0,\mathrm{},l1,l,m,m+1,\mathrm{},s\}`$, the arrow $`(p_i,p_i)`$ is in $`S^c`$. If $`ms`$, $`p_sB_{M\times N}`$ so $`p_sE^c`$. It follows that $`S^c`$ contains an arrow $`(p_s,)`$ by definition of a system on $`E^c`$ and necessarily $`=p_s`$ since $`p_s`$. Now, if $`m<s`$, $`p_{s1}E^c`$, so there is in $`S^c`$ an arrow $`(p_{s1},)`$. But distinct arrows have distinct ends so $`=p_{s1}`$. By decreasing induction, one shows that $`S^c`$ contains $`(p_m,p_m),(p_{m+1},p_{m+1}),\mathrm{},(p_s,p_s)`$.
Similarly, if $`p_0B_{M\times N}`$ then $`p_0F^c`$. By definition of $`S^c(E^c)`$, this means that $`S^c`$ contains an arrow $`(,p_0)`$ and the only possibility is $`(p_0,p_0)`$. If the monomial $`p_1B_{M\times N}`$, then $`p_1`$ is the end of an arrow $`(,p_1)`$ and we must have $`(=p_1)`$ because the arrow with origin $`p_0`$ is already determined. By induction, $`S^c`$ contains the arrows $`(p_0,p_0),\mathrm{},(p_l,p_l)`$.
###### Proposition 14.
To prove the existence of a system $`S^\nu `$ such that $`S^\nu (E^\nu )=F^\nu `$, it suffices to exhibit a system $`S^c`$ on $`E^c`$ such that $`S^c(E^c)=F^c`$.
Proof: it is an immediate consequence of the last two lemmas.
The goal of the next proposition is to exhibit such a system $`S^c`$. The cell $`C(E)`$ is an affine space as a cell of a Bialynicki-Birula stratification. So, if $`\overline{C(E)}`$ meets $`C(F)`$, there exists a morphism $`f`$ from the affine line $`𝔸^1=Speck[t]`$ to $`\overline{C(E)}`$ such that the image of the generic point is in $`C(E)`$ and such that the limit of $`f(t)`$ when $`t`$ tends to infinity is $`X_{\mathrm{}}C(F)`$. Let $`p_1,\mathrm{},p_s`$ be the points in $`𝔸^1`$ whose image by $`f`$ is in $`\overline{C(E)}C(E)`$. By the universal property of $`C(E)`$, the morphism $`f`$ is defined by a closed subscheme $`U`$ of $`𝔸^1\{p_1,\mathrm{},p_s\}\times Speck[x,y]`$. Let $`I(t)k[x,y][t]`$ be the ideal defining the closure of $`U`$ in $`𝔸^1\times Speck[x,y]`$. The staircase over the generic point being $`E`$, for each monomial $`mE`$, there exists a quasi-homogeneous element $`PI(t)`$ with initial monomial $`m`$. Among all possible $`P`$, choose one as follows. Let $`d_t(P)`$ be the degree of $`P`$ in $`t`$, and $`d_{in,t}(P)`$ be the degree in $`t`$ of the initial coefficient of $`P`$. Let $`S_m`$ be the set of $`P`$ such that $`\mathrm{\Delta }(P):=d_t(P)d_{in,t}(P)`$ is minimal. For a fixed $`P`$, denote by $`val(P)`$ the greatest monomial of $`P`$ whose coefficient has degree $`d_t(P)`$. Now choose a $`P`$ in $`S_m`$ such that $`val(P)`$ is minimum. Call this element $`P(m)`$ and put $`e(m):=val(P(m))`$.
###### Proposition 15.
Let $`S^c:=\{(m,e(m)),mE^c\}`$. Then $`S^c`$ is a system of arrows on $`E^c`$ and $`S^c(E^c)=F^c`$.
To prove the proposition, we need two lemmas.
###### Lemma 16.
Let $`P`$ be a quasi-homogeneous element of $`I(t)`$ with $`in(P)=m`$ and $`val(P)=n`$. Let $`Q`$ be a quasi-homogeneous element of $`I(t)`$ such that $`in(Q)<m`$ and $`val(Q)=n`$. Let $`b(t)`$ be the coefficient of $`n`$ in $`P`$, $`d(t)`$ be the coefficients of $`n`$ in $`Q`$. Then $`R=d(t)Pb(t)Q`$ is an element of $`I(t)`$ with initial term $`m`$ satisfying $`\mathrm{\Delta }(R)<\mathrm{\Delta }(P)`$ or ($`\mathrm{\Delta }(R)=\mathrm{\Delta }(P)`$ and $`val(R)<val(P)`$).
Proof: the fact that $`R`$ is in $`I(t)`$ with initial term $`m`$ is obvious and it remains to calculate $`\mathrm{\Delta }(R)`$. By construction, $`d_t(R)d_t(P)+d_t(Q)`$ and $`d_{in,t}(R)=d_{in,t}(P)+d_t(Q)=d_t(P)\mathrm{\Delta }(P)+d_t(Q)`$ so $`\mathrm{\Delta }(R)\mathrm{\Delta }(P)`$. If $`\mathrm{\Delta }(R)=\mathrm{\Delta }(P)`$, then $`val(R)`$ is the greater monomial of $`R`$ whose coefficient in $`t`$ has degree $`d_t(P)+d_t(Q)`$. By construction, $`R`$ has no term in $`n`$ and all monomials greater than $`n`$ have degree in $`t`$ strictly less than $`d_t(P)+d_t(Q)`$. So, the greatest monomial having degree $`d_t(P)+d_t(Q)`$ is smaller than $`n`$, showing $`val(R)<val(P)`$.
###### Lemma 17.
1) The map $`e:E^c^2`$ sending $`m`$ to $`e(m)`$ is injective.
2) $`Ime=F^c`$.
3) The inverse map to $`e:E^cF^c`$ sends a monomial $`nF`$ to the smallest monomial $`mE`$ for which there exists a quasi-homogeneous polynomial $`PI(t)`$ with $`in(P)=m`$ and $`val(P)=n`$.
Proof: 1) If $`e`$ was not injective, there would exist two monomials $`m_1`$ and $`m_2`$ with $`m_1>m_2`$ and $`val(P(m_1))=val(P(m_2))`$. Then the preceding lemma would construct from $`P(m_1)`$ and $`P(m_2)`$ an element $`RI(t)`$ contradicting the minimality of $`P(m_1)`$.
2) By definition of the flat limit, for any $`mE`$, the limit $`P_{\mathrm{}}(m)`$ of $`P(m)/t^{d_t(P(m))}`$ as $`t\mathrm{}`$ is an element of $`I(X_{\mathrm{}})`$. It follows by definition of $`X_{\mathrm{}}`$ that the initial monomial $`e(m)`$ of $`P_{\mathrm{}}(m)`$ is not in $`F`$. Conversely, we show that an element $`nF`$ can be written $`e(m)`$ for some $`mE^c`$. Let $`S`$ (resp. $`T`$) be the set of elements of $`E^c`$ (resp. of $`F^c`$) with the same degree as $`n`$. The flatness shows that $`S`$ and $`T`$ have the same (finite) cardinal. By the above, $`e`$ is an injection from $`S`$ into $`T`$, so is a bijection. In particular, $`nT`$ can be written $`e(m)`$ for some $`m`$.
3) Let $`n=e(m)`$ be a monomial in $`F^c`$. The polynomial $`P(m)I(t)`$ verifies $`in(P(m))=m`$ and $`val(P(m))=n`$. Suppose that there exists $`m^{}<m`$ and a quasi-homogeneous polynomial $`P^{}I(t)`$ such that $`in(P^{})=m^{}`$ and $`val(P^{})=n`$. Then lemma 16 would construct from $`P(m)`$ and $`P^{}`$ a polynomial $`R`$ contradicting the minimality of $`P(m)`$.
Proof of proposition 15. Point 1) of the last lemma shows that each point is the end of at most one arrow. To prove that the set is a system of arrows, the compatibility with multiplication by a monomial $`m^{}`$ remains to be shown. Let $`f=(m,e(m))`$ be an arrow. Then $`P(m)I(t)`$ verifies $`in(P(m))=m`$ and $`val(P(m))=e(m)`$. The polynomial $`m^{}P(m)I(t)`$ verifies $`in(m^{}P(m))=m^{}m,val(m^{}P(m))=m^{}e(m)`$. Then point 3) shows that the arrow $`g`$ ending at $`m^{}e(m)`$ has an origin smaller or equal to $`xm`$, i.e. $`g`$ is shorter than $`f`$. So $`S^c`$ is a system of arrows on $`E`$, and $`S^c(E^c)=F^c`$ by 2).
#### 3.2.2 Proof of point 1
It will essentially rely on the proof of point 2 thanks to a notion of duality corresponding to the notion of linkage in a complete intersection, introduced in \[PS\].
###### Definition 18.
Let $`Z`$ be a zero-dimensional scheme included in the complete intersection $`Y=(x^M,y^N)`$. The scheme $`Z^\nu `$ defined by the ideal $`(I(Y):I(Z))`$ is called the scheme obtained by linkage in the complete intersection $`(x^M,y^N)`$. Let $`(Y)`$ be the reduced subscheme of $``$ parametrizing the subschemes included in $`Y`$.
###### Proposition 19.
The morphism $`\phi :(Y)(Y)`$, $`ZZ^\nu `$ is well defined and sends $`C(E)`$ to $`C(E_{MN}^\nu )`$. Moreover $`Z^{\nu \nu }=Z`$.
Proof: if $`Y`$ is a $`0`$-dimensional Gorenstein scheme, and if $`ZY`$, then the ideal $`(I(Y):I(Z))`$ defines a scheme of degree $`deg(Y)deg(Z)`$. A complete intersection being Gorenstein, the degree of $`Z^\nu `$ is constant on the connected components of $`(Y)`$ and the morphism is well defined.
Let $`ZC(E)`$, let $`in(I(Z^\nu ))`$ be the set of initial monomials of the elements of $`I(Z^\nu )`$, and let $`e=x^{e_1}y^{e_2}in(I(Z^\nu ))`$. The inclusion $`in(I(Z^\nu ))in(I(Z))in(I(Z)I(Z^\nu ))in(x^M,y^N)`$ shows that: $`(\alpha ,\beta )E`$, $`e_1+\alpha M`$ or $`e_2+\beta N`$. In other words, $`in(I(Z^\nu ))(E_{MN}^\nu )^c`$. The complementary sets $`in(I(Z^\nu ))^c`$ and $`E_{MN}^\nu `$ have the same cardinal $`MNcard(E)`$ so $`in(I(Z^\nu ))=(E_{MN}^\nu )^c`$. This shows that $`\phi `$ sends $`C(E)`$ to $`C(E_{MN}^\nu )`$.
Finally, we have by definition of the dual, $`Z^{\nu \nu }Z`$. But $`Z`$ and $`Z^{\nu \nu }`$ have the same degree so they are equal.
###### Corollary 20.
Let $`E`$ and $`F`$ be two staircases of the same cardinal $`n`$. Then $`\overline{C(E)}C(F)\mathrm{}\overline{C(E_{nn}^\nu )}C(F_{nn}^\nu )\mathrm{}`$.
Proof: apply the dualizing morphism $`\phi `$ of the last proposition with $`M=N=n`$.
Now, we come to the proof of point 1 of theorem 8. If $`\overline{C(E)}C(F)\mathrm{}`$, then $`\overline{C(E_{nn}^\nu )}C(F_{nn}^\nu )\mathrm{}`$. Then, point 2 states that there exists a system of arrows $`S`$ on $`(E_{nn}^\nu )_{nn}^\nu =E`$ such that $`S(E)=(F_{nn}^\nu )_{nn}^\nu =F`$.
## 4 Comparison to the known results
### 4.1 Grassmannians
In the case $`H=(1,2,3,\mathrm{},k1,k,k+1n,0,0,\mathrm{},0)`$, $`_{1,1}(H)`$ is isomorphic to the Grassmannian $`G(n,k+1)=G(n,k[x,y]_k)`$. The isomorphism consists in associating $`I=I_d`$ with the vector space $`I_k`$ and the inverse isomorphism associates $`I_k`$ with the ideal $`I=I_kk[x,y]_{k+1}k[x,y]_{k+2}\mathrm{}`$ . A non increasing sequence of non negative integers $`(p_1,\mathrm{},p_n)`$ with $`p_1k+1n`$ defines a staircase $`E(p_1,\mathrm{},p_n)`$ which contains the monomials of degree at most $`k1`$ and the monomials of degree $`k`$ which are not in $`\{x^{ni+p_i}y^{kn+ip_i}`$, $`1in\}`$. Let $`V_i`$ be the vector space generated by the $`i`$ smallest monomials of degree $`k`$. The Schubert cell $`C(p_1,\mathrm{},p_n)`$ on $`G(n,k[x,y]_k)`$ is by definition $`C(p_1,\mathrm{},p_n):=\{WG(n,k[x,y]_k),dim(WV_j)=i\text{ if }k+1n+ip_ij<k+1n+ip_{i+1}\}`$. By construction, the Schubert cells $`C(E(p_1,\mathrm{},p_n))`$ and $`C(p_1,\mathrm{},p_n)`$ correspond under the above isomorphism. A classical result on Grassmannians asserts that the closure of a cell $`C(p_1,\mathrm{},p_n)`$ is the union of the cells $`C(q_1,\mathrm{},q_n)`$ for which $`q_ip_i`$ for all $`i`$. The following result says that in the particular case of Grassmannians, the two points of our necessary condition coincide, and they coincide with the known incidence condition on Grassmannians.
###### Proposition 21.
The four following conditions are equivalent:
1. $`\overline{C(p_1,\mathrm{},p_n)}C(q_1,\mathrm{},q_n)\mathrm{}`$
2. $`q_ip_i`$ for all $`i`$.
3. There exists a system $`S`$ on $`E(p_1,\mathrm{},p_n)`$ such that $`S(E(p_1,\mathrm{},p_n))=E(q_1,\mathrm{},q_n)`$
4. There exists a system $`S^\nu `$ on $`E(p_1,\mathrm{},p_n)^\nu `$ such that $`S^\nu (E(p_1,\mathrm{},p_n)^\nu )=E(q_1,\mathrm{},q_n)^\nu `$
Proof: $`12`$ is a classical result on Grassmanians.
$`13`$ and $`14`$ is exactly theorem 8.
$`41`$. If 4) is true, we have a system $`S^\nu `$ on $`E^\nu `$, which can be identified with a system $`S^c`$ on $`E^c`$ (thanks to lemma 12) such that $`S^c(E(p_1,\mathrm{},p_n)^c)=E(q_1,\mathrm{},q_n)^c`$. For each arrow $`f`$ of $`S^c`$ in degree $`k`$ with origin $`o(f)`$ and end $`g(f)`$, let $`e(f):=o(f)+tg(f)k[x,y,t]`$ and denote by $`V_k`$ the $`k[t]`$-module generated by $`\{e(f),fS^c\text{ of degree }k\}`$. Let $`I(t)k[x,y][t]`$ be the ideal defined by $`I(t):=V_kk[x,y]_{k+1}[t]k[x,y]_{k+2}[t]\mathrm{}`$. The scheme whose ideal is $`lim_t\mathrm{}I(t)=I^{E(q_1,\mathrm{},q_n)}`$ is in $`\overline{C(E(p_1,\mathrm{},p_n))}C(E(q_1,\mathrm{},q_n))`$.
$`34`$. We have $`43`$. Using the duality of proposition 19 to exchange the roles of $`E(p_1,\mathrm{},p_n)`$ and $`E(p_1,\mathrm{},p_n)^\nu `$, we have $`34`$.
### 4.2 Comparison to Yameogo’s result
The next result asserts that any of the two conditions of theorem 8 implies strictly the condition of theorem 2.
###### Proposition 22.
Let $`S`$ be a system of arrows on $`E`$. Then $`S_{S(E)}S_E`$. Moreover, there exist staircases $`E`$ and $`F`$ such that:
* $`S_ES_F`$
* there is no system of arrows $`S`$ on $`E`$ verifying $`F=S(E)`$
Let $`S^\nu `$ be a system of arrows on $`E^\nu `$. Then $`S_{S^\nu (E_{MN}^\nu )_{MN}^\nu }S_E`$. Moreover, there exist staircases $`E`$ and $`F`$ such that:
* $`S_ES_F`$
* there is no system of arrows $`S^\nu `$ on $`E_{MN}^\nu `$ verifying $`F_{MN}^\nu =S^\nu (E_{MN}^\nu )`$
Proof: The passage from $`E`$ to $`S_E`$ can be described dynamically as follows. First, you suppress all the arrows which go from a point $`P`$ to itself. Then, if there is an arrow $`f_1`$ from a monomial $`P`$ to a monomial $`Q`$ and an arrow $`f_2`$ from $`Q`$ to a monomial $`R`$, you suppress the arrows $`f_1`$ and $`f_2`$ and you replace them by the arrow from $`P`$ to $`R`$. After a finite number of such operations, you come to a finite set of arrows $`g_1,\mathrm{},g_s`$ with distinct ends such that the origins of the arrows belong to $`E`$ and such that the ends of the arrows do not belong to $`E`$. Now suppress in $`E=:E_1`$ the origin of $`g_1`$ and replace it by the end of $`g_1`$. You obtain in this way a subset $`E_2`$ of $`^2`$. Replacing in $`E_2`$ the origin of $`g_2`$ by the end of $`g_2`$, a subset $`E_3`$ of $`^2`$ is obtained. Continuing in this way, you finally reach the subset $`E_{s+1}=S(E)`$. We have by construction $`S_{S(E)}=S_{E_{s+1}}S_{E_s}\mathrm{}S_{E_2}S_{E_1}=S_E`$.
The same reasoning shows that the inequality $`S_{S^\nu (E_{MN}^\nu )_{MN}^\nu }S_E`$ always holds.
For the second part of the proposition, we consider the staircases $`E`$ and $`F`$ whose monomial ideals are $`I^E=(y^4,xy^2,x^2y,x^5)`$ and $`I^F=(y^5,xy^2,x^3)`$. They verify $`S_ES_F`$ but example 10 shows that there is no system of arrows $`S^\nu `$ on $`E_{\mathrm{5\; 5}}^\nu `$ verifying $`F_{\mathrm{5\; 5}}^\nu =S^\nu (E_{\mathrm{5\; 5}}^\nu )`$. The dual example $`I^E=(y^4,x^3y^3,x^4y,x^5)`$ and $`I^F=(y^5,x^2y^3,x^4)`$ verifies $`S_ES_F`$, but there is no system $`S`$ on $`E`$ such that $`F=S(E)`$.
## 5 An example
For a small Hilbert function, it is easy to show that the necessary condition given by our theorem is in fact sufficient. It follows that we can solve the weak incidence problem on $`_{ab}(H)`$ for small $`H`$. In this section, we illustrate this fact in the case $`(a,b)=(1,1)`$ and $`H=(1,2,3,2,1,0,0,\mathrm{})`$.
To solve the problem, we will use sufficient incidence conditions given by $`[Y2]`$. We could deal without these conditions, but we use them to shorten the proofs.
There are nine possible staircases corresponding to the Hilbert function $`H`$, all of them drawn and named in the following figure.
In the figure, a square whose position is $`(i,j)`$ relative to the square at the left bottom represents the monomial $`x^iy^j`$. We know that $`\overline{C(E_1)}C(E_2)\mathrm{}S_{E_1}S_{E_2}`$. It follows that the only possible relations are
$$\{(E_{gen},);(a,c);(a,d);(a,e);(a,f);(a,g);(a,h);(b,d);(b,e);(b,f);(b,g);$$
$$(b,h);(c,e);(c,f);(c,g);(c,h);(d,f);(d,g);(d,h);(e,g);(e,h);(f,h);(g,h)\}$$
where $``$ denotes any staircase and where we have adopted the notation $`(E_1,E_2)`$ for the incidence relation $`\overline{C(E_1)}C(E_2)\mathrm{}`$.
In the figure, a thick arrow between two staircases $`E_1`$ and $`E_2`$ means that the relation $`\overline{C(E_1)}C(E_2)`$ occurs and these strong incidence relations come from \[Y2\]. It follows that if two staircases $`E_1`$ and $`E_2`$ are linked by a sequence of thick arrows, then the relation $`(E_1,E_2)`$ holds. The cases then remaining from the above list are:
$$\{(a,e);(b,d);(b,f);(c,e);(c,g);(d,f)\}$$
Example 10 shows that $`(c,g)`$ does not hold. The five remaining cases are compatible with our necessary condition. In particular, there are systems of arrows on the dual staircases $`a^\nu ,b^\nu ,b^\nu ,c^\nu ,d^\nu `$ coming from the condition 2) of theorem 8. We interpret these systems as systems on $`a^c,b^c,b^c,c^c,d^c`$ thanks to lemma 12. The following figure shows the systems obtained in this way.
To exhibit ideals “compatible” with these systems, consider the sub-$`k[t]`$-modules of $`k[x,y,t]`$ given by their generators: $`T_3=<y^3+t^2xy^2+tx^3,x^2y>`$, $`T_4=<y^4+t^2xy^3,xy^3+tx^4,x^2y^2,x^3y>`$, $`P_5=<x^\alpha y^\beta ,\alpha +\beta 5>`$. Consider the ideal $`T:=T_3T_4P_5`$ of $`k[x,y,t]`$. The quotient $`k[x,y,t]/T`$ is flat over $`k[t]`$. The associated universal morphism sends the generic point to a point in $`C(a)`$ and can be extended by sending $`\mathrm{}`$ to the point whose ideal is the monomial ideal with staircase $`e`$. This shows that $`(a,e)`$ holds.
Similarly, if $`U_3=<y^3,xy^2+tx^2y+t^2x^3>`$, $`U_4=<y^4,xy^3,x^2y^2+tx^3y,x^4>`$, $`V_3=<y^3+tx^2y,xy^2+tx^3>`$, $`V_4=<y^4,xy^3+x^3y,x^2y^2,x^4>`$, $`W_3=<xy^2,x^2y>`$, $`W_4=<y^4+tx^4,xy^3,x^2y^2,x^3y>`$, $`Z_3=<y^3+txy^2+t^2x^2y,x^3>`$, $`Z_4=<y^4,xy^3+tx^2y^2,x^3y,x^4>`$, $`U:=U_3U_4P_5`$, $`V:=V_3V_4P_5`$, $`W:=W_3W_4P_5`$, $`Z:=Z_3Z_4P_5`$, the ideals $`U,V,W,Z`$ show respectively that the relations $`(b,d),(b,f),(c,e),(d,f)`$ are true. The diagram sums up our results: the relation $`(E,F)`$ holds if and only if the staircases $`E`$ and $`F`$ are linked by a thin arrow or by a sequence of thick arrows.
Bibliography:
\[BB\]: Bialynicki-Birula, A: Some theorems on actions of algebraic groups . Ann. of Math. 98, 480-497, (1973)
\[BB1\]: Bialynicki-Birula, A: Some properties of the decompostions of algebraic varieties determined by actions of a torus. Bulletin de l’académie Polonaise des sciences, Serie des Sciences math. astr. et phys. 24, #9, 667-674, (1976)
\[Br\]: Brion, M: Equivariant Chow groups for torus actions, Transformations Groups 2, 225–267, (1997)
\[ES\]: Ellingsrud, G. and Stromme, S.: On the homology of the Hilbert scheme of points in the plane, Invent. Math. 87 (1987), no. 2, 343–352
\[GH\]: Griffiths, P. and Harris, J.: Principles of algebraic geometry, J. Wiley, New York (1977)
\[Go\]: Göttsche, L.: Betti numbers for the Hilbert function strata of the punctual Hilbert scheme in two variables. Manuscripta Math. 66, 253-259 (1990)
\[Gr\]: Granger, JM.: Géométrie des schémas de Hilbert ponctuels, Mem. Soc. Math. Fr, Nouv. ser. 7-12(1982-83)
\[I\]: Iarrobino, A: Punctual Hilbert Scheme. Memoirs of AMS, vol. 10, #188, (1977)
\[M\]: Manivel, L.: Fonctions symétriques, polynômes de Schubert et lieux de dégénérescence, Cours specialises de la SMF, 3, (1999)
\[PS\]: Peskine C., Szpiro L.: Liaison des variétés algébriques, Invent. Math. 26, (1974), 271-302
\[Y\]: Yameogo, J: Décomposition cellulaire de variétés paramétrant des idéaux homogènes de $`/C[[x,y]]`$. Incidence des cellules I. Compositio Math. 90, #1, 81–98, (1994)
\[Y2\]: Yameogo, J: Décomposition cellulaire de variétés paramétrant des idéaux homogènes de $`/C[[x,y]]`$. Incidence des cellules II. J.reine angew. Math. 450, 123-137, (1994)
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# Quasi-vacuum solar neutrino oscillations
## I Introduction
A well-known explanation of the solar $`\nu _e`$ flux deficit is provided by flavor oscillations of neutrinos along their way from the Sun $`()`$ to the Earth $`()`$. For two active neutrino states \[say, $`(\nu _e,\nu _\mu )`$ in the flavor basis and $`(\nu _1,\nu _2)`$ in the mass basis\], the physics of solar $`\nu `$ oscillations is governed, at any given energy $`E`$, by the mass-mixing parameters $`\delta m^2`$ and $`\omega `$ in vacuum,<sup>*</sup><sup>*</sup>*One can take $`\delta m^2=m_2^2m_1^2>0`$, as far as $`\omega `$ is taken in the first quadrant $`[0,\frac{\pi }{2}]`$. as well as by the electron density profile $`N_e(x)`$ in matter .
Different oscillation regimes can be identified in terms of three characteristics lengths, namely, the astronomical unit
$$L=1.496\times 10^8\mathrm{km},$$
(1)
the oscillation length in vacuum
$$L_{\mathrm{osc}}=\frac{4\pi E}{\delta m^2}=2.48\times 10^3\left(\frac{\delta m^2/E}{\mathrm{eV}^2/\mathrm{MeV}}\right)^1\mathrm{km},$$
(2)
and the refraction length in matter
$$L_{\mathrm{mat}}=\frac{2\pi }{\sqrt{2}G_FN_e}=1.62\times 10^4\left(\frac{N_e}{\mathrm{mol}/\mathrm{cm}^3}\right)^1\mathrm{km},$$
(3)
which is associated to the effective mixing angle $`\omega _m`$ ,
$$\mathrm{sin}^22\omega _m=\frac{\mathrm{sin}^22\omega }{(\mathrm{cos}2\omega L_{\mathrm{osc}}/L_{\mathrm{mat}})^2+\mathrm{sin}^22\omega }.$$
(4)
Typical solutions to the solar neutrino problem (see, e.g., ) involve values of $`L_{\mathrm{osc}}`$ either in the so-called “just-so” (JS) oscillation regime , characterized by
$$L_{\mathrm{osc}}^{\mathrm{JS}}LL_{\mathrm{mat}},$$
(5)
or in the “Mykheyev-Smirnov-Wolfenstein” (MSW) oscillation regime , characterized by
$$L_{\mathrm{osc}}^{\mathrm{MSW}}L_{\mathrm{mat}}L.$$
(6)
The two regimes correspond roughly to $`\delta m^2/EO(10^{11})`$ eV<sup>2</sup>/MeV and to $`\delta m^2/E10^7`$ eV<sup>2</sup>/MeV, respectively.
For just-so oscillations, since $`L_{\mathrm{mat}}/L_{\mathrm{osc}}0`$, the effect of matter is basically to suppress the oscillation amplitude both in the Sun and in the Earth ($`\mathrm{sin}^22\omega _m0`$), so that (coherent) flavor oscillations take place only in vacuum, starting from the Sun surface . Conversely, for MSW oscillations, $`L_{\mathrm{osc}}L_{\mathrm{mat}}`$ and flavor transitions are dominated by the detailed matter density profile, while the many oscillation cycles in vacuum $`(L_{\mathrm{osc}}L)`$ are responsible for complete $`\nu `$ decoherence at the Earth, once smearing effects are taken into account .
Therefore, it is intuitively clear that in the intermediate range between (5) and (6), corresponding approximately to $`10^{10}\delta m^2/E10^7`$ eV<sup>2</sup>/MeV, the simple vacuum oscillation picture of the JS regime becomes increasingly decoherent and affected by matter effects for increasing values of $`\delta m^2/E`$, leading to a hybrid regime that might be called of “quasi-vacuum” (QV) oscillations, characterized by
$$L_{\mathrm{mat}}L_{\mathrm{osc}}^{\mathrm{QV}}L,$$
(7)
In the past, quasi-vacuum effects on the oscillation amplitude and phase have been explicitly considered only in relatively few papers (see, e.g., ) as compared with the vast literature on solar neutrino oscillations, essentially because typical fits to solar $`\nu `$ rates allowed only marginal solutions in the range where QV effects are relevant. However, more recent analyses appear to extend the former ranges of the JS solutions upwards and of the MSW solutions downwards in $`\delta m^2/E`$, making them eventually merge in the QV range , especially under generous assumptions about the experimental or theoretical $`\nu `$ flux uncertainties.Such extended range for current solutions reflects, in part, the lack of a strong, model-independent signature of energy dependence in the solar neutrino deficit. Therefore, a fresh look at QV corrections seems warranted. Recently, a semianalytical approximation improving the familiar just-so formula in the QV regime was discussed in and, in more detail, in , where additional numerical checks were performed. In this work we revisit the whole topic, by performing accurate numerical calculations which include the exact density profile in the Sun and in the Earth, within the reference mass-mixing ranges $`\delta m^2/E[10^{10},10^7]`$ eV<sup>2</sup>/MeV and $`\mathrm{tan}^2\omega [10^3,10]`$. The variable $`\mathrm{tan}^2\omega `$ is useful to chart the first two octants of the mixing angle range . We also discuss some approximations that can simplify the computing task in present applications. We then apply such calculations to a global analysis of solar neutrino data in the range $`\delta m^210^8`$ eV<sup>2</sup>.
Our paper is structured as follows. The basic notation and the numerical techniques used in the calculations are introduced in Sec. II and III, respectively. The effects of solar matter in the quasi-vacuum oscillation regime are discussed in Sec. IV, where the results for true and exponential density profiles are compared. Earth matter effects are described in Sec. V. The decoherence of oscillations induced by energy (and time) integration is discussed in Sec. VI. The basic results are summarized and organized in Sec. VII, and then applied to a three-flavor oscillation analysis in Sec. VIII. Section IX concludes our work.
## II Notation
The $`\nu `$ propagation from the Sun core to the detector at the Earth can be interpreted as a “double slit experiment,” where the original $`\nu _e`$ can take two paths, corresponding to the intermediate transitions $`\nu _e\nu _1`$ and to $`\nu _e\nu _2`$. The global $`\nu _e`$ survival amplitude is then the sum of the amplitudes along the two paths,
$`A(\nu _e\nu _e)`$ $`=`$ $`A_{}(\nu _e\nu _1)A_{\mathrm{vac}}(\nu _1\nu _1)A_{}(\nu _1\nu _e)`$ (8)
$`+`$ $`A_{}(\nu _e\nu _2)A_{\mathrm{vac}}(\nu _2\nu _2)A_{}(\nu _2\nu _e),`$ (9)
where the transition amplitudes from the Sun production point to its surface $`(A_{})`$, from the Sun surface to the Earth surface $`(A_{\mathrm{vac}})`$ and from the Earth surface to the detector $`(A_{})`$ have been explicitly factorized. The $`\nu _e`$ survival probability $`P_{ee}`$ is then given by
$$P_{ee}=|A(\nu _e\nu _e)|^2.$$
(10)
In general, the above amplitudes can be written as
$`A_{}(\nu _e\nu _1)`$ $`=`$ $`\sqrt{P_{}}\mathrm{exp}(i\xi _{}),`$ (12)
$`A_{\mathrm{vac}}(\nu _1\nu _1)`$ $`=`$ $`\mathrm{exp}\left(im_1^2(LR_{})/2E\right),`$ (13)
$`A_{}(\nu _1\nu _e)`$ $`=`$ $`\sqrt{P_{}}\mathrm{exp}(i\xi _{}),`$ (14)
for the first path and as
$`A_{}(\nu _e\nu _2)`$ $`=`$ $`\sqrt{1P_{}},`$ (16)
$`A_{\mathrm{vac}}(\nu _2\nu _2)`$ $`=`$ $`\mathrm{exp}\left(im_2^2(LR_{})/2E\right),`$ (17)
$`A_{}(\nu _2\nu _e)`$ $`=`$ $`\sqrt{1P_{}},`$ (18)
for the second path, where $`R_{}`$ is the Sun radius.<sup>§</sup><sup>§</sup>§Although $`R_{}`$ is relatively small ($`R_{}/L=4.7\times 10^3`$), it is explicitly kept for later purposes. The Earth radius $`R_{}`$ can instead be safely neglected ($`R_{}/L=4.3\times 10^5`$). In the above equations, $`P_{}`$ and $`P_{}`$ are real numbers ($`[0,1]`$) representing the transition probability $`P(\nu _e\nu _1)`$ along the two partial paths inside the Sun (up to its surface) and inside the Earth (up to the detector). The corresponding phase differences between the two paths, $`\xi _{}`$ and $`\xi _{}`$ $`([0,2\pi ])`$, have been associated to the first path without loss of generality. The $`\nu _e`$ survival probability $`P_{ee}`$ from Eq. (10) reads then
$`P_{ee}`$ $`=`$ $`P_{}P_{}+(1P_{})(1P_{})`$ (20)
$`+2\sqrt{P_{}(1P_{})P_{}(1P_{})}\mathrm{cos}\xi ,`$
where the global oscillation phase is given by
$$\xi =\frac{\delta m^2L}{2E}(1\delta _R\delta _{}\delta _{}),$$
(21)
with the definitions
$`\delta _R`$ $`=`$ $`{\displaystyle \frac{R_{}}{L}}=4.7\times 10^3,`$ (22)
$`\delta _{}`$ $`=`$ $`{\displaystyle \frac{2E}{\delta m^2L}}\xi _{},`$ (23)
$`\delta _{}`$ $`=`$ $`{\displaystyle \frac{2E}{\delta m^2L}}\xi _{}.`$ (24)
A relevant extension of the $`2\nu `$ formula (20) is obtained for $`3\nu `$ oscillations, required to accommodate solar and atmospheric neutrino data . Assuming a third mass eigenstate $`\nu _3`$ with $`m^2=|m_3^2m_{1,2}^2|\delta m^2`$, the $`3\nu `$ survival probability can be written as
$$P_{ee}^{3\nu }=c_\varphi ^4P_{ee}^{2\nu }+s_\varphi ^4,$$
(25)
where $`\varphi `$ is the $`(\nu _e,\nu _3)`$ mixing angle, and $`P_{ee}^{2\nu }`$ is given by Eq. (20), provided that the electron density $`N_e`$ is replaced everywhere by $`c_\varphi ^2N_e`$ (see and refs. therein). Such replacement implies that the $`3\nu `$ case is not a simple mapping of the $`2\nu `$ case, and requires specific calculations for any given value of $`\varphi `$.
We conclude this section by recovering some familiar expressions for $`P_{ee}`$, as special cases of Eq. (20). The JS limit ($`L_{\mathrm{mat}}/L_{\mathrm{osc}}0`$, with complete suppression of oscillations inside matter) corresponds to $`P_{}c_\omega ^2P_{}`$ and to negligible $`\delta _{}`$, $`\delta _{}`$. Then, neglecting also $`\delta _R`$, one gets from (20) the standard “vacuum oscillation formula,”
$$P_{ee}^{\mathrm{JS}}1\mathrm{sin}^22\omega \mathrm{sin}^2(\pi L/L_{\mathrm{osc}}).$$
(26)
In the MSW limit ($`L/L_{\mathrm{osc}}\mathrm{}`$), the global oscillation phase $`\xi `$ is very large and $`\mathrm{cos}\xi 0`$ on average. Furthermore, assuming for $`P_{}`$ a well-known approximation in terms of the “crossing” probability $`P_c`$ between mass eigenstates in matter \[in our notation, $`P_{}\mathrm{sin}^2\omega _m^0P_c+\mathrm{cos}^2\omega _m^0(1P_c)`$, with $`\omega _m^0`$ calculated at the production point\], one gets from (20) and for daytime $`(P_{}=c_\omega ^2)`$ the so-called Parke’s formula ,
$$P_{ee,\mathrm{day}}^{\mathrm{MSW}}\frac{1}{2}+(\frac{1}{2}P_c)\mathrm{cos}2\omega \mathrm{cos}2\omega _m^0.$$
(27)
## III Numerical Techniques
In general, numerical calculations of the $`\nu `$ transition amplitudes must take into account the detailed $`N_e`$ profile along the neutrino trajectory, both in the Sun and in the Earth.
Concerning the Sun, we take $`N_e`$ from (“year 2000” standard solar model). Figure 1 shows such $`N_e`$ profile as a function of the normalized radius $`r/R_{}`$, together with its exponential approximation $`N_e=N_e^0\mathrm{exp}(r/r_0)`$, with $`N_e^0=245`$ mol/cm<sup>3</sup> and $`r_0=R_{}/10.54`$. For the exponential density profile, the neutrino evolution equations can be solved analytically . In order to calculate the relevant probability $`P_{}`$ and the phase $`\xi _{}`$, we have developed several computer programs which evolve numerically the familiar MSW neutrino evolution equations along the Sun radius, for generic production points, and for any given value of $`\delta m^2/E[10^{10},10^7]`$ eV<sup>2</sup>/MeV and of $`\mathrm{tan}^2\omega `$. We estimate a numerical (fractional) accuracy of our results better than $`10^4`$, as derived by several independent checks. As a first test, we integrate numerically the MSW equations both in their usual complex form (2 real + 2 imaginary components) and in their Bloch form involving three real amplitudes , obtaining the same results. We have then repeated the calculations with different integration routines taken from several computer libraries, and found no significant differences among the outputs. We have optionally considered, besides the exact $`N_e`$ profile, also the exponential profile, which allows a further comparison of the numerical integration of the MSW equations with their analytical solutions, as worked out in in terms of hypergeometric functions (that we have implemented in an independent code). Also in this case, no difference is found between the output of the different codes.
Concerning the calculation of the quantities $`P_{}`$ and $`\xi _{}`$ in the Earth, we evolve analytically the MSW equations at any given nadir angle $`\eta `$, using the technique described in , which is based on a five-step biquadratic approximation of the density profile from the Preliminary Reference Earth Model (PREM) and on a first-order perturbative expansion of the neutrino evolution operator. Such analytical technique provides results very close to a full numerical evolution of the neutrino amplitudes, the differences being smaller than those induced by uncertainties in $`N_e`$ . In particular, we have checked that, for $`\delta m^2/E10^7`$ eV<sup>2</sup>/MeV, such differences are $`10^3`$. In conclusion, we are confident in the accuracy of our results, which are discussed in the following sections.
## IV Matter effects in the Sun
Figure 2 shows, in the mass-mixing plane and for standard solar model density, isolines of the difference $`c_\omega ^2P_{}`$ (solid curves), which becomes zero in the just-so oscillation limit of very small $`\delta m^2/E`$. The isolines shape reminds the “lower corner” of the more familiar MSW triangle . Also shown are isolines of constant resonance radius $`R_{\mathrm{res}}/R_{}`$ (dotted curves), defined by the MSW resonance condition $`L_{\mathrm{osc}}/L_{\mathrm{mat}}(R_{\mathrm{res}})=\mathrm{cos}2\omega `$. The values of $`c_\omega ^2P_{}`$ are already sizable (a few percent) at $`\delta m^2/E10^9`$, and increase for increasing $`\delta m^2/E`$ and for large mixing $`[\mathrm{tan}^2\omega O(1)]`$, especially in the first octant, where the MSW resonance can occur. The difference between matter effects in the first and in the second octant can lead to observable modifications of the allowed regions in fits to the data , and to a possible discrimination between the cases $`\omega <\frac{\pi }{4}`$ and $`\omega >\frac{\pi }{4}`$ .
In the whole parameter range of Fig. 2, it turns out that, within the region of $`\nu `$ production ($`r/R_{}0.3`$), it is $`L_{\mathrm{mat}}(r)L_{\mathrm{osc}}`$ (and thus $`\mathrm{sin}^22\omega _m^00`$).This is also indicated by the fact that $`R_{\mathrm{res}}/R_{}0.55`$ in the $`\delta m^2/E`$ range of Fig. 2. As a consequence, all the curves of Fig. 2 do not depend on the specific $`\nu `$ production point (as we have also checked numerically), and no smearing over the $`\nu `$ source distribution is needed in the quasi-vacuum regime. This is a considerable simplification with respect to the MSW regime, which involves higher values of $`\delta m^2/E`$ and thus shorter (resonance) radii, which are sensitive to the detailed $`\nu `$ source distribution.
Figure 3 shows, in the same coordinates of Fig. 2, the isolines of $`c_\omega ^2P_{}`$ corresponding to the exponential density profile (dotted curves), for which we have used the fully analytical results of . (Identical results are obtained by numerical integration.) The solid lines in Fig. 3 refer to a well-known approximation (sometimes called “semianalytical”) to such results, which is obtained in the limit $`N_e0`$ at the Sun surface (it is not exactly so for the exponential profile, see Fig. 1). More precisely, the zeroth order expansion of the hypergeometric functions in terms of the small parameter $`z=i\sqrt{2}G_Fr_0N_e(R_{})`$ $`(|z|0.16)`$ gives, for $`N_e(0)\mathrm{}`$, the result $`P_{}[\mathrm{exp}(\gamma c_\omega ^2)1]/[\mathrm{exp}(\gamma )1]`$ with $`\gamma =\pi r_0\delta m^2/E`$ , which leads to the QV prescription discussed in . The differences between the solid and dotted curves in Fig. 3 are essentially due to the “solar border approximation” ($`N_e0`$) assumed in the semianalytical case; indeed, the differences would practically vanish if the exponential density profile, and thus the “effective” Sun radius, were unphysically continued for $`rR_{}`$ (not shown). Such limitations of the semianalytical approximation have been qualitatively suggested by the authors of but, contrary to their claim, our Fig. 3 shows explicitly that the semianalytical calculation of $`P_{}`$ represents a reasonable approximation to the analytical one for $`\omega `$ in both octants, as also verified in .
A comparison of the results of Fig. 2 (true density) and of Fig. 3 (exponential density) shows that, in the latter case, the correction term $`c_\omega ^2P_{}`$ tends to be somewhat overestimated, in particular when the semianalytical approximation is used. We have verified that such bias is dominantly due to the difference (up to a factor of $`2`$) between the true density profile and its exponential approximation around $`r/R_{}0.8`$ (see Fig. 1) and, subdominantly, to the details of the density profile shape at the border ($`r/R_{}1`$). As a consequence, the “exponential profile” calculation of $`P_{}`$ (either semianalytic or analytic) tends to shift systematically the onset of solar matter effects to lower values of $`\delta m^2/E`$. For instance, at $`\mathrm{tan}^2\omega 1`$ (maximal mixing), the value $`c_\omega ^2P_{}=0.05`$ is reached at $`\delta m^2/E8\times 10^{10}`$ eV<sup>2</sup>/MeV for the true density, and at $`\delta m^2/E`$ a factor of $`2`$ lower for the exponential density. In order to avoid artificially larger effects at low $`\delta m^2/E`$ in neutrino data analyses, one should numerically calculate $`P_{}`$ with the true electron density profile. The difference between the numerical calculation and the semianalytic approximation is also briefly discussed in for $`\delta m^2/E10^8`$ eV<sup>2</sup>/MeV (where $`P_cP_{}`$).
Concerning the phase factor $`\delta _{}`$, we confirm earlier indications about its smallness, in both cases of true and exponential density. In the latter case, the semianalytic approximation gives , for $`\delta m^2/E0`$, the $`\omega `$-independent result
$$\delta _{}L^1\left\{r_0\left[\mathrm{ln}(\sqrt{2}G_FN_e^0r_0)+\gamma _E\right]R_{}\right\},$$
(28)
where $`\gamma _E(0.577)`$ is the Euler constant. Numerically, $`\delta _{}^05.39\times 10^4`$, much smaller than $`\delta _R`$. As far as $`\delta m^2/E10^8`$ ($`10^7`$) eV<sup>2</sup>/MeV, we find that $`\delta _{}`$ differs from such value by less than $`20\%`$ (50%), the exact difference depending on the value of $`\mathrm{tan}^2\omega `$ (not shown). Therefore, $`\delta _{}`$ is an order of magnitude smaller than $`\delta _R`$ in the whole mass-mixing range considered. A similar behavior is found by using the true density profile $`(|\delta _{}|10^3`$ everywhere). As for $`P_{}`$, we find also for $`\delta _{}`$ no significant dependence on the neutrino production point in the Sun core.
In conclusion, we find only modest differences between the analytical calculations of $`P_{}`$, based on the exponential density profile, and its semianalytical approximation. The difference with respect to the numerical calculation of $`P_{}`$, based on the standard solar model profile, is instead more pronounced, and leads to a factor of $`2`$ difference in the value of $`\delta m^2/E`$ where QV effects start to be significant. Therefore, we recommend the use of the electron profile from the standard solar model, implying numerical calculations of $`P_{}`$.The interested reader can obtain numerical tables of $`P_{}`$, calculated for the standard solar model density, upon request from the authors. Concerning the phase $`\delta _{}`$, it simply shifts the global oscillation phase $`\xi `$ in Eq. (20) by less than one permill, and thus it can be safely neglected in all current applications.
## V Matter effects in the Earth
Strong Earth matter effects typically emerge in the range where $`L_{\mathrm{osc}}L_{\mathrm{mat}}`$ within the mantle ($`N_e2`$ mol/cm<sup>3</sup>) or the core ($`N_e5`$ mol/cm<sup>3</sup>), as well as in other ranges of mantle-core oscillation interference , globally corresponding to $`\delta m^2/E10^7`$$`10^6`$ eV<sup>2</sup>/MeV. Therefore, only marginal effects are expected in the parameter range considered in this work, as confirmed by the results reported in Fig. 4.
Figure 4 shows isolines of the quantity $`c_\omega ^2P_{}`$, which becomes zero in the just-so oscillation limit of very small $`\delta m^2/E`$. The solid curves corresponds to a nadir angle $`\eta =0^{}`$ (diametral crossing of neutrinos) and the dotted curves to $`\eta =45^{}`$ (crossing of mantle only). For other values of $`\eta `$ (not shown), the quantity $`c_\omega ^2P_{}`$ has a comparable magnitude. In the current neutrino jargon, the Earth effect shown in Fig. 4 is operative in the lowermost part of the so-called “LOW” MSW solution to the solar neutrino problem, or, from another point of view, to the uppermost part of the vacuum solutions . Concerning the phase correction $`\delta _{}`$ (not shown), it is found to be smaller than $`1.5\times 10^5`$ in the whole mass-mixing plane considered, and thus can be safely neglected.
In practical applications, the correction term $`c_\omega ^2P_{}`$ must be time-averaged. This poses, in principle, a tedious integration problem, since such correction appears, in Eq. (20), both in the amplitude of the oscillating term ($`\mathrm{cos}\xi `$) and in the remaining, non-oscillating term. While the integration over time can be transformed, for the non-oscillating term, into a more manageable integration over $`\eta `$ , this cannot be done for the oscillating term, which depends on time both through the prefactor $`\sqrt{P_{}(1P_{})}`$ and through the phase $`\xi `$ (via eccentricity effects). However, in region of the mass-mixing plane where the Earth effect is non-negligible (namely, where $`c_\omega ^2P_{}`$ few %), it turns out that $`P_{}0`$, so that the total amplitude of the oscillating term is small. Moreover, as it will be discussed in the next section, energy smearing effects strongly suppress the oscillation factor $`\mathrm{cos}\xi `$ in the same region. Therefore, for $`\delta m^2/E10^8`$ eV<sup>2</sup>/MeV, the oscillating term in Eq. (20) is doubly suppressed, and one can safely consider only the non-oscillatory terms.
## VI Damping of the interference term
In the just-so regime, the interference factor $`\mathrm{cos}\xi `$ in Eq. (20) can lead to an observable modulation both in the energy and in the time distribution of solar neutrino events. This modulation gradually disappears as $`\delta m^2/E`$ increases, as a consequence of several decoherence effects ($`\mathrm{cos}\xi 0`$) , which are typically dominated by energy smearing . The broader the neutrino spectrum, the lower the values of $`\delta m^2/E`$ where decoherence becomes important: therefore, it suffices to consider the narrowest spectra (the so-called Be neutrino “lines” ) for illustration purposes.
Let us consider a $`\nu `$ “line” energy spectrum $`s(E)`$, with $`𝑑Es(E)=1`$ and $`E=𝑑Es(E)E`$. The energy-averaged oscillating factor,
$$C=\mathrm{cos}\left(\frac{\delta m^2L}{2E}\right)_E=𝑑Es(E)\mathrm{cos}\left(\frac{\delta m^2L}{2E}\right),$$
(29)
can be written, in the narrow-width approximation $`(\mathrm{\Delta }E=EEE)`$, in terms of the Fourier transform of the spectrum,
$$\stackrel{~}{s}(\tau )=𝑑Es(E)e^{i\mathrm{\Delta }E\tau }.$$
(30)
More precisely,
$`C`$ $``$ $`{\displaystyle 𝑑Es(E)\mathrm{cos}\left(\frac{\delta m^2L}{2E}\left(1\frac{\mathrm{\Delta }E}{E}\right)\right)}`$ (31)
$`=`$ $`D\mathrm{cos}\left({\displaystyle \frac{\delta m^2L}{2E}}(1\delta )\right),`$ (32)
where $`D=|\stackrel{~}{s}(\tau )|`$, $`\delta =\tau ^1E^1\mathrm{arg}\stackrel{~}{s}(\tau )`$, and $`\tau =\delta m^2L/2E^2`$.
Using the $`s(E)`$ profile for the Be lines from , we find that the phase correction $`\delta `$ is smaller than a few permill in the region where the damping factor $`D`$ is greater than a few percent. If $`\delta `$ is neglected, the average oscillation factor can be simply written as the oscillation factor at the average energy, times a damping term $`D`$:
$$\mathrm{cos}\left(\frac{\delta m^2L}{2E}\right)_ED\mathrm{cos}\left(\frac{\delta m^2L}{2E}\right).$$
(33)
A similar approach to smearing effects for neutrino lines was developed in .
Figure 5 shows the damping factor $`D`$ for the two Be neutrino lines at $`E=0.863`$ and $`0.386`$ MeV. The factor is negligible for $`\delta m^2/E10^8`$ eV<sup>2</sup>/MeV, implying that the oscillation pattern is completely smeared out in such range.<sup>\**</sup><sup>\**</sup>\**This range corresponds approximately to $`L_{\mathrm{osc}}/Lw/E`$, where $`wO(1\mathrm{keV})`$ is the line width. The onset of smearing effects is shifted to even lower values of $`\delta m^2/E`$ for continuous energy spectra. In conclusion, the integration over neutrino energy makes the average oscillation factor $`\mathrm{cos}\xi `$ always negligible in the range where Earth matter effects are important, even for the narrow Be neutrino lines. Such accidental simplification should not make one forget that, in general, the transition from coherent to incoherent oscillations is a complex phenomenon that deserves further studies (see, e.g., for a recent ab initio approach).
Finally, we recall that integration over time acts as a further damping factor . At first order in the eccentricity $`ϵ`$, one has that $`L(t)L[1ϵ\mathrm{cos}\alpha (t)]`$ and $`\dot{\alpha }(t)1+2ϵ\mathrm{cos}\alpha `$, where $`\alpha `$ is the orbital phase, and $`t[0,2\pi ]`$ is the normalized time variable ($`t=0`$ at perihelion). Then the yearly average of $`\mathrm{cos}\xi `$ \[times the square-law factor $`L^2/L^2(t)`$\] can be performed analytically , giving
$$\frac{1}{2\pi }_0^{2\pi }𝑑t\frac{L^2}{L^2(t)}\mathrm{cos}\left(\frac{\delta m^2L(t)}{2E}\right)J_0\left(ϵ\frac{\delta m^2L}{2E}\right)\mathrm{cos}\left(\frac{\delta m^2L}{2E}\right),$$
(34)
where $`J_0`$ is the Bessel function, acting as a further damping term for large values of its argument.
Notice that the maximum fractional variation of the orbital radius, $`(L_{\mathrm{max}}L_{\mathrm{min}})/L=2ϵ=3.34\times 10^2`$, is an order of magnitude larger than $`\delta _R=R_{}/L=4.7\times 10^3`$ which, in turn, is larger than the phase corrections $`\delta _{}`$ and $`\delta _{}`$. Therefore, one can safely neglect $`\delta _R`$, $`\delta _{}`$ and $`\delta _{}`$ in practical applications involving yearly (or even seasonal) averages, as we do in this work. However, for averages over shorter time intervals, such approximation might break down. In particular, $`\delta _R`$ ($`\delta _{}`$) might be comparable to the monthly (weekly) variations of the solar neutrino signal. The observability of such short-time variations is beyond the present sensitivity of real-time solar $`\nu `$ experiments and would require, among other things, very high statistics and an extremely stable level of both the signal detection efficiency and of the background. If such difficult experimental goals will be reached in the future, some of the approximations discussed so far (and recollected in the next section) should be revisited and possibly improved.
## VII Practical Recipes
We have seen in the previous sections that, as $`\delta m^2/E`$ increases, the deviations of $`P_{}`$ (and subsequently of $`P_{}`$) from the vacuum value $`c_\omega ^2`$ become increasingly important. We have also seen that the phase corrections $`\delta _{}`$ and $`\delta _{}`$ are smaller than $`\delta _R=R_{}/L`$, which can in turn be neglected in present applications, so that one can practically take the usual vacuum value for the oscillation phase, $`\xi \delta m^2L/2E`$. We think it useful to organize known and less known results through the following approximate expressions for the calculation of $`P_{ee}`$, which are accurate to better than 3% with respect to the exact, general formula (20) valid at any $`\delta m^2/E`$.
For $`\delta m^2/E5\times 10^{10}`$ eV<sup>2</sup>/MeV, one can take $`P_{}P_{}c_\omega ^2`$, and obtain the just-so oscillation formula
$$P_{ee}^{\mathrm{JS}}c_\omega ^4+s_\omega ^4+2s_\omega ^2c_\omega ^2\mathrm{cos}\xi ,$$
(35)
with $`\xi =\delta m^2L/2E`$. For $`5\times 10^{10}\delta m^2/E10^8`$ eV<sup>2</sup>/MeV, one can still take $`P_{}c_\omega ^2`$, but since $`P_{}c_\omega ^2`$ (quasi-vacuum regime) one has that
$$P_{ee}^{\mathrm{QV}}c_\omega ^2P_{}+s_\omega ^2(1P_{})+2s_\omega c_\omega \sqrt{P_{}(1P_{})}\mathrm{cos}\xi ,$$
(36)
where $`\xi =\delta m^2L/2E`$, and $`P_{}`$ has to be calculated numerically (see also for earlier versions of the above equation).
Finally, for $`\delta m^2/E10^8`$ eV<sup>2</sup>/MeV, also Earth matter effects are important $`(P_{}c_\omega ^2)`$; however, this complication is balanced by the disappearance of the oscillating term ($`\mathrm{cos}\xi 0`$) due to unavoidable smearing effects, so that the usual MSW regime is recovered:
$$P_{ee}^{\mathrm{MSW}}P_{}P_{}+(1P_{})(1P_{}).$$
(37)
Time averages are then relatively simple to implement. In Eqs. (35) and (36), the yearly averages of $`\mathrm{cos}\xi `$ can be performed analytically \[Eq. (34)\]. In Eq. (37), the nighttime average of $`P_{}`$ can be transformed into a more manageable integration over the nadir angle, both for yearly and for seasonal averages.
To summarize, the above sequence of equations describes the passage from the regime of just-so to that of MSW oscillations, via quasi-vacuum oscillations. In the JS regime, oscillations are basically coherent and do not depend on the electron density in the Earth or in the Sun ($`N_e\mathrm{}`$). In the MSW regime, oscillations are basically incoherent ($`L\mathrm{}`$) and, in general, depend on the detailed electron density profile of both the Sun and the Earth. In particular, in the MSW regime one has to take into account the interplay between the density profile and the neutrino source distribution profile. The intermediate QV regime is instead characterized by partially coherent oscillations (with increasing decoherence as $`\delta m^2/E`$ increases), and by a sensitivity to the electron density of the Sun (but not of the Earth). Such sensitivity is not as strong as in the MSW regime and, in particular, QV effects are independent from the specific $`\nu `$ production point, which can be effectively taken at the Sun center.
For the sake of completeness, we mention that, for high values of $`\delta m^2/E`$ ($`10^4`$ eV<sup>2</sup>/MeV), corresponding to $`L_{\mathrm{osc}}L_{\mathrm{mat}}`$ in the Sun, the sensitivity to matter effects is eventually lost both in the Sun and in the Earth ($`P_{}P_{}c_\omega ^2`$), and one reaches a fourth regime sometimes called of energy-averaged (EA) oscillations, which is totally incoherent and $`N_e`$-independent:
$$P_{ee}^{\mathrm{EA}}c_\omega ^4+s_\omega ^4.$$
(38)
Such regime, which predicts an energy-independent suppression of the solar neutrino flux, seems to be disfavored (but perhaps not yet ruled out) by current experimental data on total neutrino rates. In conclusion, for $`\delta m^2/E`$ going from extremely low values to infinity, one can identify four rather different oscillation regimes,
$$\mathrm{JS}\mathrm{QV}\mathrm{MSW}\mathrm{EA},$$
(39)
each being characterized by specific properties and applicable approximations. Experiments still have to tell us unambiguously which of them truly applies to solar neutrinos.
## VIII Three-flavor oscillation analysis
As discussed in , in the QV regime the $`2\nu `$ survival probability (36) is non-symmetric with respect to the operation $`\omega \frac{\pi }{2}\omega `$, which instead holds for JS oscillations \[Eq. (35)\].<sup>††</sup><sup>††</sup>††In the MSW regime, the mirror asymmetry of the first two octants was explicitly shown in . Here we extend such observation to $`3\nu `$ oscillations, under the hypothesis $`\delta m^2m^2`$ which leads to Eq. (25). The $`2\nu `$ case is recovered for $`\varphi =0`$.
Equation (25) for $`P_{3\nu }`$ preserves the (a)symmetry properties of $`P_{2\nu }`$ under the replacement $`\omega \frac{\pi }{2}\omega `$. Therefore, while $`P_{3\nu }^{\mathrm{JS}}`$ \[obtained from Eqs. (25) and (35)\] is symmetric with respect to the $`\omega =\frac{\pi }{4}`$ value, the expression of $`P_{3\nu }^{\mathrm{QV}}`$ \[obtained from Eqs. (25) and (36)\] is not. Such properties become evident in the triangular representation of the solar $`3\nu `$ mixing parameter space discussed in , to which the reader is referred for further details.
Figure 6 shows, in the triangular plot, isolines of $`P_{3\nu }^{\mathrm{QV}}`$ (dotted lines) for $`\delta m^2/E`$ close to $`1.65\times 10^9`$ eV<sup>2</sup>/MeV, corresponding to about 100 oscillation cycles. More precisely, the six panels correspond to $`\xi =100\times 2\pi +\mathrm{\Delta }\xi `$, with $`\mathrm{\Delta }\xi `$ from 0 to $`\pi `$ in steps of $`\pi /5`$. The dotted isolines are asymmetric with respect to $`\omega =\frac{\pi }{4}`$, as is the QV correction $`P_{}c_\omega ^2`$. For increasing $`s_\varphi ^2`$ (upper part of the triangle), the symmetry tends to be restored, mainly because the asymmetric term $`c_\varphi ^4P_{2\nu }`$ in Eq. (25) is suppressed by the prefactor $`c_\varphi ^4`$ and, marginally, because the effective electron density $`c_\varphi ^2N_e`$ in $`P_{2\nu }`$ is also suppressed.<sup>‡‡</sup><sup>‡‡</sup>‡‡When the effective electron density $`c_\varphi ^2N_e`$ is small ($`s_\varphi ^21`$), the statement that the quasi-vacuum probability $`P_{2\nu }`$ does not depend on the production point in the core (see Sec. IV) is not strictly valid, and one should in principle consider also the $`\nu `$ source profile. However, since $`P_{2\nu }`$ is correspondingly suppressed by $`c_\varphi ^4`$, the source smearing effect is numerically small, and one can effectively discard it (e.g., by taking the production point at $`r=0`$ at any $`s_\varphi ^2`$). Exact mirror symmetry at any $`\varphi `$ is restored only if the QV corrections are switched off (solid lines); in such case, $`P_{3\nu }`$ is a quadratic form in the coordinates $`s_\varphi ^2`$ and $`s_\omega ^2`$, and the isolines are conical curves. In any case, $`P_{ee}`$ becomes typically too low (too high) in the central region (in the corners) of the triangle where, as a consequence, one should expect only marginal solutions to the solar neutrino deficit.
The asymmetry with respect to $`\omega =\frac{\pi }{4}`$ also shows up in solar neutrino data fits . Figure 7 reports the results of our global (rates + spectrum + day/night) three-flavor analysis in the mass-mixing range $`\delta m^2[10^{11},10^8]`$ eV<sup>2</sup> and $`\mathrm{tan}^2\omega [10^2,10^2]`$, for several representative values of $`\mathrm{tan}^2\varphi `$. We only show 99% C.L. contours<sup>\**</sup><sup>\**</sup>\**At 95% C.L. the solutions would be mostly located at $`\delta m^210^9`$ eV<sup>2</sup>. ($`N_{\mathrm{DF}}=3`$) for the sake of clarity. The theoretical and experimental inputs, as well as the $`\chi ^2`$ statistical analysis , are the same as in Ref. (where MSW solutions were studied). Here, however, the range of $`\delta m^2`$ is lower, in order to show the smooth transition from the MSW solutions to the QV and finally the JS ones as $`\delta m^2`$ is decreased. In particular, the solutions shown in Fig. 7 represent the continuation, at low $`\delta m^2`$, of the LOW MSW solutions shown in Fig. 10 of (panel by panel).<sup>\*†</sup><sup>\*†</sup>\*†A technical remark is in order. The minimum value of $`\chi ^2`$ in the plane of Fig. 7 ($`\chi _{\mathrm{min}}^2=31.8`$) is reached at $`(\delta m^2/\mathrm{eV}^2,\mathrm{tan}^2\omega ,\mathrm{tan}^2\varphi )=(4.4\times 10^{10},2.4,0.1)`$. For $`\delta m^2>10^8`$ eV<sup>2</sup> (MSW regime), the minimum value is $`\chi _{\mathrm{min}}^2=27.0`$ for the same input data. In order to match the results of Fig. 7 in this work and of Fig. 10 in , we have adopted the absolute minimum ($`\chi _{\mathrm{min}}^2=27.0`$) to draw the $`\mathrm{\Delta }\chi ^2=11.34`$ contours in Fig. 7. As anticipated in the comments to Fig. 6, the mirror asymmetry around $`\mathrm{tan}^2\omega =1`$ decreases for decreasing $`\delta m^2`$ (JS regime); a little asymmetry is still present even at $`\delta m^210^{10}`$ eV<sup>2</sup>, where the gallium rates (sensitive to $`E`$ as low as $`0.2`$ MeV) start to feel QV effects. In the region where QV effects are important, the solutions are typically shifted in the second octant $`(\omega \pi /4`$), since the gallium rate is suppressed too much in the first octant (see also ). A similar drift was found for the LOW MSW solution . At any $`\delta m^2`$, the asymmetry decreases at large values of $`\varphi `$ ($`\mathrm{tan}^2\varphi 1.5`$) which, however, are excluded by the combination of accelerator and reactor data , unless the second mass square difference $`m^2`$ turns out to be in the lower part of the sensitivity range of the CHOOZ experiment ($`m^210^3`$ eV<sup>2</sup>). For $`\varphi =0`$, the standard two-flavor case is recovered, and the results are comparable to those found in . The Super-Kamiokande spectrum plays only a marginal role in generating the mirror asymmetry of the solutions in Fig. 7, since the modulation of QV effects in the energy domain is much weaker than the one generated by the oscillation phase $`\xi `$. We find that, at any given $`\delta m^210^8`$ eV<sup>2</sup>, the $`\chi ^2`$ difference at symmetric $`\omega `$ values is less than $`1`$ for the 18-bin spectrum data fit. Therefore, QV effects are mainly probed by total neutrino rates at present.
We think it is not particularly useful to discuss more detailed features of the current QV solutions, such as combinations of spectral data or rates only, fits with variations of hep neutrino flux, etc. (which were instead given in for MSW solutions). In fact, while the shapes of current MSW solutions are rather well-defined, those of JS or QV solutions are still very sensitive to small changes in the theoretical or experimental input. Therefore, a detailed analysis of the “fine structure” of the QV solutions in Fig. 7 seems unwarranted at present.
Finally, in Fig. 8 we show sections of the allowed $`3\nu `$ solutions (at 99% C.L.) in the triangle representation, for six selected (increasing) values of $`\delta m^2`$. Solutions are absent or shrunk at $`s_\varphi ^20.5`$, where the theoretical $`\nu `$ flux underestimates the gallium and water-Cherenkov data. The lowest value of $`\delta m^2`$ ($`0.66\times 10^{10}`$ eV<sup>2</sup>) falls in the JS regime, so that the ring-like allowed region (which resembles the curves of iso-$`P_{ee}`$ in Fig. 6) is symmetric with respect to the vertical axis at $`\omega =\pi /4`$. However, as $`\delta m^2`$ increases and QV effects become operative, the solutions become more and more asymmetric, and shifted towards the second octant of $`\omega `$ .
Figures 7 and 8 in this work, as well as Fig. 10 in , show that solar neutrino data, by themselves, put only a weak upper bound on the mixing angle $`\varphi `$. Much tighter constraints are set by reactor data , unless the second mass square difference $`m^2`$ happens to be $`10^3`$ eV<sup>2</sup> (which seems an unlikely possibility). In any case, QV effects are operative also for a small (or zero) value of $`s_\varphi ^2`$.
## IX Conclusions
We have presented a thorough analysis of solar neutrino oscillations in the “quasi-vacuum” oscillation regime, intermediate between the familiar just-so and MSW regimes. The QV regime is increasingly affected by matter effects for increasing values of $`\delta m^2`$. We have calculated such effects both in the Sun and in the Earth, and discussed the accuracy of various possible approximations. We have implemented the QV oscillation probability in a full three-flavor analysis of solar neutrino data, obtaining solutions which smoothly join (at $`\delta m^210^8`$ eV<sup>2</sup>) the LOW MSW regions found in for the same input data. The asymmetry of QV effects makes such solutions different for $`\omega <\frac{\pi }{4}`$ and $`\omega >\frac{\pi }{4}`$, the two cases being symmetrized only in the just-so oscillation limit of small $`\delta m^2`$.
###### Acknowledgements.
We thank J.N. Bahcall for providing us with updated standard solar model results. We thank A. Friedland and S.T. Petcov for useful discussions.
.
FIG. 1. Radial profile of the electron density in the Sun from the standard solar model (solid line), together with its exponential approximation (dashed line).
.
FIG. 2. Isolines of the solar correction term $`c_\omega ^2P_{}`$ in the mass-mixing plane (solid curves), for standard solar model density . Isolines of MSW resonance radii (dashed curves) are also shown.
.
FIG. 3. As in Fig. 2, but for the exponential density profile. Dotted and solid lines correspond to analytical calculations and to their semianalytical approximation, respectively.
.
FIG. 4. Isolines of the Earth correction term $`c_\omega ^2P_{}`$ in the mass-mixing plane, for PREM density. The solid and dotted lines correspond to a nadir angle equal to $`0^{}`$ and $`45^{}`$, respectively.
.
FIG. 5. Damping factor for the oscillating term, due to the finite Be line width. See the text for details.
.
FIG. 6. Three-flavor oscillations: Isolines of $`P_{ee}`$ in the triangular representation of the solar $`\nu `$ mixing parameter space, for representative values of the oscillation phase $`\xi `$. The dotted and solid lines refer, respectively, to calculations with and without quasi-vacuum effects. When such effects are included, the mirror symmetry $`\omega \frac{\pi }{2}\omega `$ is broken.
.
FIG. 7. Three-flavor oscillations: Contours at 99% C.L., as derived from a global analysis of neutrino data for $`\delta m^210^8`$ eV<sup>2</sup> (quasi-vacuum regime). The asymmetry of the solutions is evident for increasing values of $`\delta m^2`$. The solutions represent the continuation (at low $`\delta m^2`$) of the MSW 99% C.L. regions reported in Fig. 10 of .
.
FIG. 8. Sections of the $`3\nu `$ allowed volume (99% C.L.) for six representative values of $`\delta m^2`$ (0.66, 2.36, 6.42, 8.76, 10.8, and $`29.5\times 10^{10}`$ eV<sup>2</sup>), shown in the solar triangle plot.
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# Untitled Document
$`p`$-ADIC AND ADELIC HARMONIC OSCILLATOR
WITH TIME-DEPENDENT FREQUENCY
Goran S. Djordjević<sup>1</sup> and Branko Dragovich<sup>2,3</sup>
<sup>1</sup>Department of Physics, University of Niš, P.O.Box 91, 18001 Niš, Yugoslavia
<sup>2</sup>Institute of Physics, P.O.Box 57, 11001 Belgrade, Yugoslavia
<sup>3</sup>Steklov Mathematical Institute, Russian Academy of Sciences, Moscow, Russia
The classical and quantum formalism for a $`p`$-adic and adelic harmonic oscillator with time-dependent frequency is developed, and general formulae for main theoretical quantities are obtained. In particular, the $`p`$-adic propagator is calculated, and the existence of a simple vacuum state as well as adelic quantum dynamics is shown. Space discreteness and $`p`$-adic quantum-mechanical phase are noted.
1. Introduction
In quantum-mechanical experiments, as well as in all measurements, numerical results belong to the field of rational numbers Q. In principle, the corresponding theoretical models could be made using only Q, but it would missed usual effectiveness and beauty of mathematical analysis. So, instead of Q one traditionally applies the field of real numbers R in classical mechanics and the field of complex numbers C in quantum mechanics. R is completion of Q with respect to the metric induced by the absolute value and C is an algebraic extension of R. In addition to R there exist the fields of $`p`$-adic numbers $`\text{Q}_{\mathrm{}}`$ as completions of Q with respect to $`p`$-adic norms ($`p`$= a prime number) . According to the Ostrowski theorem, R and $`\text{Q}_{\mathrm{}}`$ (for every $`p`$) exhaust all possible completions of Q. Thus Q is dense not only in R but also in each $`\text{Q}_{\mathrm{}}`$. Therefore, in the last decade there has been a lot and successful interest in construction of theoretical models with $`p`$-adic numbers (for a review, see, Refs. 2-5).
There is a common belief that none separated prime number $`p`$ plays a special role in physics and that $`p`$-adic models have to be taken together for all primes. It is clear that $`p`$-adic models, having some physical meaning, must be somehow connected with the ordinary (real) ones. The space of adeles A is a mathematical instrument which enables us to consider real and $`p`$-adic numbers simultaneously and as a whole. Thus it is natural to expect that adelic approach provides a more complete description of a physical system than the ordinary one.
$`p`$-Adic numbers exhibit ultrametric (non-archimedean) properties, which may be realized in quantum systems at very short distances. Possibility that space-time at the Planck scale exhibits $`p`$-adic and adelic structure is one of the main physical motivations to investigate the corresponding models.
In order to start with a systematic approach to $`p`$-adic models of quantum systems, $`p`$-adic quantum mechanics was formulated. Quantization is done along the Weyl procedure. The corresponding Hilbert space $`L_2(\text{Q}_{\mathrm{}}\mathrm{}`$ contains complex-valued square integrable functions on $`\text{Q}_{\mathrm{}}`$. Instead of the Schr$`\ddot{o}`$dinger equation, the dynamical evolution and the spectral problem of a system are related to the unitary representation of the evolution operator $`U_p(t)`$ on $`L_2(\text{Q}_{\mathrm{}}\mathrm{}`$. As a generalization and unification of $`p`$-adic and ordinary quantum mechanics, recently was formulated adelic quantum mechanics .
So far a rather small number of physical systems has been treated in $`p`$-adic and adelic quantum mechanics: a non-relativistic free particle , a harmonic oscillator , a particle in a constant field , the de Sitter minisuperspace model of the universe and a relativistic free particle . It is doubtless that evaluation of some other physical systems, which exhibit $`p`$-adic and adelic properties, will give new insights into this subject and new directions for future investigations at the Planck scale.
In this paper we show existence and some properties of $`p`$-adic and adelic harmonic oscillator with time-dependent frequency (HOTDF). Model of the HOTDF has vast applications from quantum optics to quantum cosmology . Nevertheless, many properties of classical and quantum motion can be found without specifying the time dependence of $`\omega (t)`$.
2. $`p`$-Adic numbers and adeles
To make this paper more self contained, we give here a very short review of some basic facts on $`p`$-adic numbers and adeles.
Any rational number $`x0`$ can be presented as $`x=p^\nu \frac{m}{n}`$, where $`\nu ,m,n\text{Z}`$ and $`p`$ is a given prime number which divides neither $`m`$ nor $`n`$. By definition, $`p`$-adic norm of $`x`$ is
$$|x|_p=p^\nu ,|0|_p=0,$$
$`(2.1)`$
and holds the strong triangle inequality:
$$|x+y|_p\text{max}(|x|_p,|y|_p).$$
$`(2.2)`$
A norm (valuation) with the property (2.2) is called non-archimedean or ultrametric norm. Every $`p`$-adic number $`x`$ can be uniquely presented by the canonical expansion
$$x=p^\nu \underset{i=0}{\overset{+\mathrm{}}{}}x_ip^i,x_i\{0,1,\mathrm{},p1\},x_00,\nu \text{Z}$$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
The expansion (2.3) is convergent with respect to the metric induced by $`p`$-adic norm, i.e. $`d_p(x,y)=|xy|_p`$.
There are mainly two kinds of analysis on $`\text{Q}_{\mathrm{}}`$ based on two different maps: $`\text{Q}_{\mathrm{}}\text{Q}_{\mathrm{}}`$ and $`\text{Q}_{\mathrm{}}\text{C}`$. We use both of these analyses.
Elementary $`p`$-adic functions, like $`\mathrm{exp}x`$, $`\mathrm{sin}x`$ and $`\mathrm{cos}x`$ are given by series of the same form as in the real case. However, the region of convergence is rather restricted and it is $`|x|_p<|2|_p`$ for the above functions. Derivatives of $`p`$-adic valued functions are defined as in the real case, but using $`p`$-adic norm instead of the absolute value.
For complex-valued functions of $`p`$-adic argument there is well-defined integration with the Haar measure. In particular, we use the Gauss integral
$$_{x_pp^\nu }\chi _p(\alpha x^2+\beta x)𝑑x=\{\begin{array}{cc}p^\nu \mathrm{\Omega }(p^\nu |\beta |_p),\hfill & |\alpha |_pp^{2\nu },\hfill \\ \lambda _p(\alpha )|2\alpha |_p^{1/2}\chi _p\left(\frac{\beta ^2}{4\alpha }\right)\mathrm{\Omega }\left(p^\nu |\frac{\beta }{2\alpha }|_p\right),\hfill & |4\alpha |_p>p^{2\nu }.\hfill \end{array}$$
$`(2.4)`$
$`\chi _p(u)=\mathrm{exp}(2\pi i\{u\}_p)`$ is a p-adic additive character, where $`\{u\}_p`$ denotes the fractional part of $`u\text{Q}_{\mathrm{}}`$. $`\lambda _p(\alpha )`$ is an arithmetic complex-valued function with the following basic properties :
$$\lambda _p(0)=1,\lambda _p(a^2\alpha )=\lambda _p(\alpha ),\lambda _p(\alpha )\lambda _p(\beta )=\lambda _p(\alpha +\beta )\lambda _p(\alpha ^1+\beta ^1),|\lambda _p(\alpha )|_{\mathrm{}}=1.$$
$`(2.5)`$
$`\mathrm{\Omega }(|u|_p)`$ is the characteristic function on $`\text{Z}_{\mathrm{}}`$, i.e.
$$\mathrm{\Omega }(|u|_p)=\{\begin{array}{cc}1,\hfill & |u|_p1,\hfill \\ 0,\hfill & |u|_p>1,\hfill \end{array}$$
$`(2.6)`$
where $`\text{Z}_{\mathrm{}}\mathrm{}\{\mathrm{}\text{Q}_{\mathrm{}}\mathrm{}\mathrm{𝕜}\mathrm{}\mathrm{𝕜}_{\mathrm{}}\mathrm{}\}`$ is the ring of p-adic integers.
An adele $`a\text{A}`$ is an infinite sequence
$$a=(a_{\mathrm{}},a_2,\mathrm{},a_p,\mathrm{}),$$
$`(2.7)`$
where $`a_{\mathrm{}}\text{R}`$ and $`a_p\text{Q}_{\mathrm{}}`$ with the restriction that $`a_p\text{Z}_{\mathrm{}}`$ for all but a finite set $`S`$ of primes $`p`$. The set of all adeles A can be written in the form
$$\text{A}\mathrm{}\underset{\mathrm{𝕊}}{𝕌}\mathrm{𝔸}\mathrm{}\mathrm{𝕊}\mathrm{}\mathrm{}\mathrm{𝔸}\mathrm{}\mathrm{𝕊}\mathrm{}\mathrm{}\text{R}\times \underset{\mathrm{}\mathrm{𝕊}}{}\text{Q}_{\mathrm{}}\times \underset{\mathrm{}\mathrm{𝕊}}{}\text{Z}_{\mathrm{}}\mathrm{}$$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
A is a topological space. It is a ring with respect to componentwise addition and multiplication. There is a natural generalization of analysis on R and $`\text{Q}_{\mathrm{}}`$ to analysis on A.
3. Classical oscillator: real, $`p`$-adic and adelic case
Classical HOTDF is given by the Lagrangian
$$L(x,\dot{x},t)=\frac{m}{2}\dot{x}^2\frac{m\omega ^2(t)}{2}x^2,$$
$`(3.1)`$
where $`m\text{Q}`$. Time-dependent frequency $`\omega (t)=_{n0}\omega _nt^n,`$ where $`\omega _n\text{Q}`$, is assumed to be an analytic function on $`\text{D}_{\mathrm{}}\text{R}`$ and on $`\text{D}_{\mathrm{}}\text{Z}_{\mathrm{}}`$ for all $`p`$. In other words, when $`t𝒜(S)`$ then $`\omega (t)𝒜(S^{})`$, where $`S`$ and $`S^{}`$ are some finite sets of primes $`p`$. In the real case $`m,x,\dot{x},t,\omega (t)\text{R}\text{Q}_{\mathrm{}}`$ (in the sequel index $`\mathrm{}`$ denotes quantities defined on R or C) and the analogous situation is for the $`p`$-adic counterparts. Because of formal similarity of analyses, evaluation of (3.1) is the same in real and $`p`$-adic dynamics. Thus, in real and $`p`$-adic cases, the equation of motion is
$$\ddot{x}(t)+\omega ^2(t)x(t)=0$$
$`(3.2)`$
with general solution
$$x(t)=G(t)[C_1\mathrm{cos}\gamma (t)+C_2\mathrm{sin}\gamma (t)].$$
$`(3.3)`$
The amplitude $`G(t)`$ and phase $`\gamma (t)`$ satisfy equations
$$G^3(t)\ddot{G}(t)+\omega ^2(t)G^4(t)=C^2,\dot{\gamma }(t)G^2(t)=C,$$
$`(3.4)`$
where $`C`$ is a constant $`(0<C\text{R}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\text{Z}_{\mathrm{}}\mathrm{}`$ and can be taken $`C=1`$. We are interested in analytic solution of (3.3), where $`G(t)`$ and $`\gamma (t)`$ are power series in $`t`$ with rational coefficients. Differential equation for $`G(t)`$ is non-linear. However, it does not lead to non-linear algebraic equations for unknown coefficients $`G_n`$ in expansion $`G(t)=_{n0}G_nt^n`$ and any $`G_n`$ can be presented as a rational number, which is the same in the real and all $`p`$-adic cases. Note that usual power series with rational coefficients which are convergent on $`\text{D}_{\mathrm{}}\text{R}`$ in the real case are also $`p`$-adically convergent in some region $`\text{D}_{\mathrm{}}\text{Z}_{\mathrm{}}`$.
As an illustration of analytic solutions of the equations (3.4) (with $`C=1`$) we present two simple examples.
Example 1. Let $`\omega (t)=\omega _0/(1+at)^2`$, where $`\omega _0=b^2`$ and $`a,b\text{N}`$. Then
$$G(t)=b(1+at),\gamma (t)=\frac{1}{b^2}\frac{t}{1+at}.$$
Since $`\gamma (t)`$ is argument of trigonometric functions in (3.3) one obtains that common region of convergence for all analytic expansions is $`t_p<2b^2_p`$ for each $`p`$.
Example 2. Let $`\omega (t)=\omega _0/(1+at)`$, where $`\omega _0=b^2(1+a^2b^4/4)^{\frac{1}{2}}`$ and $`a,b2\text{N}`$. Then in an analogous way to the Example 1 we get:
$$G(t)=b(1+at)^{\frac{1}{2}},\gamma (t)=\frac{1}{ab^2}\mathrm{ln}(1+at),t_p<2b^2_p.$$
Thus, there exist non-trivial adelic solutions for $`G(t)`$, and $`\gamma (t)`$, and conseqently for $`x(t)`$ in the form (3.3).
To determine constants $`C_1`$ and $`C_2`$ we use two kinds of conditions on the classical trajectory.
(3.1) Solution with the end point conditions
The classical trajectory that links two space-time points $`(x^{},t^{})`$ and $`(x^{\prime \prime },t^{\prime \prime })`$ is
$$x(t)=\frac{G(t)}{\mathrm{sin}(\gamma ^{\prime \prime }\gamma ^{})}\left[\frac{x^{}}{G^{}}\mathrm{sin}(\gamma ^{\prime \prime }\gamma (t))+\frac{x^{\prime \prime }}{G^{\prime \prime }}\mathrm{sin}(\gamma (t)\gamma ^{})\right],$$
$`(3.5)`$
where $`x^{}=x(t^{}),x^{\prime \prime }=x(t^{\prime \prime }),G^{}=G(t^{}),G^{\prime \prime }=G(t^{\prime \prime }),\gamma ^{}=\gamma (t^{})`$ and $`\gamma ^{\prime \prime }=\gamma (t^{\prime \prime })`$. Note that condition $`\gamma ^{\prime \prime }\gamma ^{}m\pi ,m\text{Z}`$, must be satisfied in the real case. Recall also that $`|\gamma ^{\prime \prime }\gamma ^{}|_p<|2|_p`$. As we shall see later it is useful to write the corresponding momentum in the form
$$k(t)=m\dot{x}(t)=m\frac{\dot{G}(t)}{G(t)}x(t)+\frac{mG(t)\dot{\gamma }(t)}{\mathrm{sin}(\gamma ^{\prime \prime }\gamma ^{})}\left[\frac{x^{\prime \prime }}{G^{\prime \prime }}\mathrm{cos}(\gamma (t)\gamma ^{})\frac{x^{}}{G^{}}\mathrm{cos}(\gamma ^{\prime \prime }\gamma (t))\right].$$
$`(3.6)`$
(3.2) Solution with the initial conditions
Imposing the initial conditions $`x^0=x(t^0),k^0=m\dot{x}(t^0)`$ we find evolution of the classical state as follows:
$$x(t)=\left[\frac{G(t)}{G^0}\mathrm{cos}(\gamma (t)\gamma ^0)\frac{G(t)\dot{G}^0}{C}\mathrm{sin}(\gamma (t)\gamma ^0)\right]x^0+\frac{G(t)G^0}{mC}\mathrm{sin}(\gamma (t)\gamma ^0)k^0,$$
$$k(t)=\left[m\left(\frac{\dot{G}(t)}{G^0}\frac{G(t)\dot{\gamma }(t)\dot{G}^0}{C}\right)\mathrm{cos}(\gamma (t)\gamma ^0)m\left(\frac{\dot{G}(t)\dot{G}^0}{C}+\frac{G(t)\dot{\gamma }(t)}{G^0}\right)\mathrm{sin}(\gamma (t)\gamma ^0)\right]x^0$$
$$+\frac{G^0}{C}\left[G(t)\dot{\gamma }(t)\mathrm{cos}(\gamma (t)\gamma ^0)+\dot{G}(t)\mathrm{sin}(\gamma (t)\gamma ^0)\right]k^0,$$
$`(3.7)`$
where $`G^0=G(t^0)`$ and $`\gamma ^0=\gamma (t^0)`$. Putting first $`t=t^{}`$ and then $`t=t^{\prime \prime }`$ in the first equation of (3.7) one can find $`x^0`$ and $`k^0`$ as functions of $`x^{}`$ and $`x^{\prime \prime }`$. Inserting these $`x^0=x^0(x^{},x^{\prime \prime })`$ and $`k^0=k^0(x^{},x^{\prime \prime })`$ into the second equation of (3.7) one gets the same formula (3.6) for $`k(t)`$.
A suitable way to calculate the corresponding classical action
$$\overline{S}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=\frac{m}{2}_t^{}^{t^{\prime \prime }}[\dot{x}^2(t)\omega ^2(t)x^2(t)]𝑑t$$
$`(3.8)`$
is integrating by parts and using the equation of motion (3.2). It leads to
$$\overline{S}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=\frac{m}{2}(x^{\prime \prime }\dot{x}^{\prime \prime }x^{}\dot{x}^{}).$$
$`(3.9)`$
In virtue of (3.6), that gives $`\dot{x}`$ as function of $`x`$, we find action in the form quadratic in $`x^{\prime \prime }`$ and $`x^{}`$, i.e.
$$\overline{S}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=\frac{m}{2}[(\frac{\dot{\gamma }^{\prime \prime }}{\mathrm{tan}(\gamma ^{\prime \prime }\gamma ^{})}+\frac{\dot{G}^{\prime \prime }}{G^{\prime \prime }})x^{\prime \prime 2}$$
$$\frac{2\sqrt{\dot{\gamma }^{\prime \prime }\dot{\gamma }^{}}}{\mathrm{sin}(\gamma ^{\prime \prime }\gamma ^{})}x^{\prime \prime }x^{}+(\frac{\dot{\gamma }^{}}{\mathrm{tan}(\gamma ^{\prime \prime }\gamma ^{})}\frac{\dot{G}^{}}{G^{}})x^2],$$
$`(3.10)`$
where we used equality
$$\frac{G^{\prime \prime }\dot{\gamma ^{\prime \prime }}}{G^{}}+\frac{G^{}\dot{\gamma ^{}}}{G^{\prime \prime }}=2\sqrt{\dot{\gamma ^{\prime \prime }}\dot{\gamma ^{}}},$$
which is derived by means of (3.4).
Note that the above $`p`$-adic formalism has the same form as its real counterpart, or in other words, the classical HOTDF is invariant under change of the number field R and $`\text{Q}_{\mathrm{}}`$, for every $`p`$. This may be regarded as a necessary condition for existence of an adelic classical HOTDF, which we construct in the following way. Let the position $`x`$, momentum $`k`$ and time $`t`$ be adelic quantities, like (2.7). The corresponding adelic Langrangian is given by
$$L(x,\dot{x},t)=(L(x_{\mathrm{}},\dot{x}_{\mathrm{}},t_{\mathrm{}}),L(x_2,\dot{x}_2,t_2),\mathrm{},L(x_p,\dot{x}_p,t_p),\mathrm{}),$$
where $`L(x_v,\dot{x}_v,t_v)=m[\dot{x}_v^2\omega ^2(t_v)x_v^2]/2`$ with $`v=\mathrm{},2,\mathrm{},p,\mathrm{}`$, and $`|L(x_p,\dot{x}_p,t_p)|_p1`$ for all but a finite number of primes $`p`$. Also, all the other above introduced quantities, regarded as real and $`p`$-adic, can be generalized to the adelic ones. For instance, adelic classical action is
$$\overline{S}(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=(\overline{S}(x_{\mathrm{}}^{\prime \prime },t_{\mathrm{}}^{\prime \prime };x_{\mathrm{}}^{},t_{\mathrm{}}^{}),\overline{S}(x_2^{\prime \prime },t_2^{\prime \prime };x_2^{},t_2^{}),\mathrm{},\overline{S}(x_p^{\prime \prime },t_p^{\prime \prime };x_p^{},t_p^{}),\mathrm{}),$$
$`(3.11)`$
where real and $`p`$-adic ingredients have the form (3.10).
4. Quantum oscillator: real, $`p`$-adic and adelic case
The unique formalism of ordinary quantum mechanics which enables $`p`$-adic and adelic generalization with complex-valued wave functions is a triple
$$(L_2(\text{R}\mathrm{}\mathrm{}\mathrm{𝕎}\mathrm{}ϝ_{\mathrm{}}\mathrm{}\mathrm{}\mathrm{𝕌}\mathrm{}\mathrm{}_{\mathrm{}}\mathrm{}\mathrm{}\mathrm{}$$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
where $`L_2(\text{R}\mathrm{}`$ is the Hilbert space, $`z_{\mathrm{}}`$ is a point of real classical phase space, $`W(z_{\mathrm{}})`$ is a unitary representation of the Heisenberg-Weyl group on $`L_2(\text{R}\mathrm{}`$, and $`U(t_{\mathrm{}})`$ is a unitary representation of the evolution operator on $`L_2(\text{R}\mathrm{}`$. Hence, under $`p`$-adic and adelic quantum mechanics we understand $`p`$-adic and adelic analogues of (4.1), i.e.
$$(L_2(\text{Q}_{\mathrm{}}\mathrm{}\mathrm{}\mathrm{𝕎}\mathrm{}ϝ_{\mathrm{}}\mathrm{}\mathrm{}\mathrm{𝕌}\mathrm{}\mathrm{}_{\mathrm{}}\mathrm{}\mathrm{}\mathrm{}$$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
$$(L_2(\text{A}\mathrm{}\mathrm{}\mathrm{𝕎}\mathrm{}ϝ\mathrm{}\mathrm{}\mathrm{𝕌}\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}$$
$`\mathrm{}\mathrm{}\mathrm{}\mathrm{}\mathrm{}`$
respectively. Thus, to find adelic eigenstate and its evolution of a HOTDF given by $`U(t^{\prime \prime },t^{})`$, one has to solve the equation
$$U(t^{\prime \prime },t^{})\mathrm{\Psi }_S^{(\alpha )}(x^{},t^{})=\chi [\alpha (\gamma ^{\prime \prime }\gamma ^{})]\mathrm{\Psi }_S^{(\alpha )}(x^{},t^{}),$$
$`(4.4)`$
where $`\alpha =(\alpha _{\mathrm{}},\alpha _2,\mathrm{},\alpha _p,\mathrm{})`$ is an adelic analogue of energy, $`\chi (u)=_v\chi _v(u_v)=\mathrm{exp}(2\pi iu_{\mathrm{}})_p\mathrm{exp}(2\pi i\{u_p\}{}_{p}{}^{})`$ and
$$\mathrm{\Psi }_S^{(\alpha )}(x,t)=\mathrm{\Psi }_{\mathrm{}}^{(\alpha _{\mathrm{}})}(x_{\mathrm{}},t_{\mathrm{}})\underset{pS}{}\mathrm{\Psi }_p^{(\alpha _p)}(x_p,t_p)\underset{pS}{}\mathrm{\Omega }(x_p_p).$$
$`(4.5)`$
The evolution operator $`U(t^{\prime \prime },t^{})=_vU_v(t_v^{\prime \prime },t_v^{})`$ acts componentwise as follows:
$$[U_v\mathrm{\Psi }_v](x_v^{\prime \prime },t_v^{\prime \prime })=_\text{Q}_{\mathrm{}}𝒦_v(x_v^{\prime \prime },t_v^{\prime \prime };x_v^{},t_v^{})\mathrm{\Psi }_v(x_v^{},t_v^{})𝑑x_v^{}.$$
$`(4.6)`$
The kernel $`𝒦_v(x_v^{\prime \prime },t_v^{\prime \prime };x_v^{},t_v^{}))`$ is defined by the Feynman path integral
$$𝒦_v(x_v^{\prime \prime },t_v^{\prime \prime };x_v^{},t_v^{})=\chi _v\left(\frac{1}{h}S[x]\right)𝒟x=\chi _v\left(\frac{1}{h}_{t_v^{}}^{t_v^{\prime \prime }}L(x_v,\dot{x}_v,t_v)𝑑t_v\right)\underset{t_v}{}dx(t_v),$$
$`(4.7)`$
where $`h`$ is the Planck constant. The kernel $`𝒦`$, also called the quantum-mechanical propagator, is of central importance not only in ordinary but also in $`p`$-adic and adelic quantum mechanics.
The $`p`$-adic Feynman path integral for classical actions quadratic in $`x^{\prime \prime }`$ and $`x^{}`$ is calculated in and has the same form as its real counterpart. Namely, if $`\overline{S}(x_v^{\prime \prime },t_v^{\prime \prime };x_v^{},t_v^{})`$ is quadratic in $`x_v^{\prime \prime }`$ and $`x_v^{}`$ then
$$𝒦_v(x_v^{\prime \prime },t_v^{\prime \prime };x_v^{},t_v^{})=\lambda _v\left(\frac{1}{2h}\frac{^2\overline{S}}{x_v^{\prime \prime }x_v^{}}\right)\left|\frac{1}{h}\frac{^2\overline{S}}{x_v^{\prime \prime }x_v^{}}\right|_v^{1/2}\chi _v\left(\frac{1}{h}\overline{S}(x_v^{\prime \prime },t_v^{\prime \prime };x_v^{},t_v^{})\right),$$
$`(4.8)`$
where $`\lambda _{\mathrm{}}(\alpha )=(1isign\alpha )/\sqrt{2}`$ , $`\lambda _{\mathrm{}}(0)=1`$ and satisfies properties (2.5).
Applying formula (4.8) to the HOTDF and using (3.10), we get
$$𝒦_v(x_v^{\prime \prime },t_v^{\prime \prime };x_v^{},t_v^{})=\lambda _v\left(\frac{m}{2h}\frac{\sqrt{\dot{\gamma }^{\prime \prime }\dot{\gamma }^{}}}{\mathrm{sin}(\gamma ^{\prime \prime }\gamma ^{})}\right)\left|\frac{m}{h}\frac{\sqrt{\dot{\gamma }^{\prime \prime }\dot{\gamma }^{}}}{\mathrm{sin}(\gamma ^{\prime \prime }\gamma ^{})}\right|_v^{1/2}$$
$$\chi _v\left\{\frac{m}{2h}\left[\left(\frac{\dot{\gamma }^{\prime \prime }}{\mathrm{tan}(\gamma ^{\prime \prime }\gamma ^{})}+\frac{\dot{G}^{\prime \prime }}{G^{\prime \prime }}\right)x^{\prime \prime 2}\frac{2\sqrt{\dot{\gamma }^{\prime \prime }\dot{\gamma }^{}}}{\mathrm{sin}(\gamma ^{\prime \prime }\gamma ^{})}x^{\prime \prime }x^{}+\left(\frac{\dot{\gamma }^{}}{\mathrm{tan}(\gamma ^{\prime \prime }\gamma ^{})}\frac{\dot{G}^{}}{G^{}}\right)x^2\right]\right\},$$
$`(4.9)`$
that contains the earlier obtained result in the real case (see, e.g. Ref. 14). One can explicitly show that the propagator (4.9) satisfies all usual properties of the probability amplitude for a quantum particle to go from a space-time point $`(x_v^{},t_v^{})`$ to a space-time point $`(x_v^{\prime \prime },t_v^{\prime \prime })`$.
The corresponding adelic propagator is
$$𝒦(x^{\prime \prime },t^{\prime \prime };x^{},t^{})=\underset{v}{}𝒦_v(x_v^{\prime \prime },t_v^{\prime \prime };x_v^{},t_v^{}),$$
$`(4.10)`$
where $`𝒦_v(x_v^{\prime \prime },t_v^{\prime \prime };x_v^{},t_v^{})`$ is given by (4.9). Product in (4.10) is divergent, but it may be regarded as an adelic functional on the space of test functions which are the adelic Schwartz-Bruhat functions (see, also ).
In $`p`$-adic quantum mechanics a significant role plays the eigenstate $`\mathrm{\Omega }(|x_p|_p)`$ (2.6), which is invariant under $`U_p(t_p^{\prime \prime },t_p^{})`$ transformation and may be regarded as $`p`$-adic vacuum state since it has $`\{\alpha _p(\gamma _p^{\prime \prime }\gamma _p^{})\}_p=0`$. Due to (4.5), the existence of $`\mathrm{\Omega }(|x_p|_p)`$ for all but a finite number of $`p`$ is a necessary condition for unification of ordinary and $`p`$-adic quantum mechanics in the form of adelic one. This $`\mathrm{\Omega }`$-state exists iff
$$_{x_p^{}_p1}𝒦_p(x_p^{\prime \prime },t_p^{\prime \prime };x_p^{},t_p^{})𝑑x_p^{}=\mathrm{\Omega }(|x_p^{\prime \prime }|_p)$$
$`(4.11)`$
is satisfied.
Inserting (4.9) with $`v=p`$ into (4.11), and using the integral (2.4) for $`\nu =0`$, we can derive some conditions on $`G(t_p),\gamma (t_p)`$ and $`m`$, which provide $`\mathrm{\Omega }`$-eigenstate. For example, if $`\gamma (t_p)=\gamma _0+\gamma _1t_p+\gamma _2t_p^2+\mathrm{}+\gamma _nt_p^n+\mathrm{}.`$ and
$$\left|\frac{\dot{G}^{}}{G^{}}\right|_p<\left|\frac{\dot{\gamma }^{}}{\mathrm{tan}(\gamma ^{\prime \prime }\gamma ^{})}\right|_p>\left|\frac{h}{2m}\right|_p$$
then $`\mathrm{\Omega }(|x_p|_p)`$ exists for all $`p2`$. It is worth noting that not every HOTDF has $`\mathrm{\Omega }`$-state and may be adelically generalized.
As the simplest illustration of the above expressions one can take frequency $`\omega (t)=\omega _0`$ and recover earlier obtained result .
5. Concluding remarks
According to (4.5) adelic wave function $`\mathrm{\Psi }(x,t)`$ offers more information on a physical system than only its standard part $`\mathrm{\Psi }_{\mathrm{}}(x_{\mathrm{}},t_{\mathrm{}})`$. Let us note here space discreteness and $`p`$-adic phase, which are generic and mainly follow from adelic quantum formalism.
For example, adelic state
$$\mathrm{\Psi }(x,t)=\mathrm{\Psi }_{\mathrm{}}(x_{\mathrm{}},t_{\mathrm{}})\underset{p}{}\mathrm{\Omega }(|x_p|_p)$$
exhibits discrete structure of the space at the length $`l_0=(hm^1\omega ^1)^{1/2}`$. Namely, according to the usual interpretation of the wave function we have to consider $`|\mathrm{\Psi }(x,t)|_{\mathrm{}}^2`$ at rational points $`x`$ and $`t`$. In the above adelic case we get
$$|\mathrm{\Psi }(x,t)|_{\mathrm{}}^2=|\mathrm{\Psi }_{\mathrm{}}(x,t)|_{\mathrm{}}^2\underset{p}{}\mathrm{\Omega }(|x|_p)=\{\begin{array}{cc}|\mathrm{\Psi }_{\mathrm{}}(x,t)|_{\mathrm{}}^2,\hfill & x\text{Z}\mathrm{}\hfill \\ 0,\hfill & x\text{Q}\text{Z}\mathrm{}\hfill \end{array}$$
Here we used the following properties of $`\mathrm{\Omega }`$-function: $`\mathrm{\Omega }^2(|x|_p)=\mathrm{\Omega }(|x|_p),_p\mathrm{\Omega }(|x|_p)=1`$ if $`x\text{Z}`$, and $`_p\mathrm{\Omega }(|x|_p)=0`$ if $`x\text{Q}\text{Z}`$. Thus, it means that position $`x`$ may have only discrete values: $`x/l_0=0,\pm 1,\pm 2,\mathrm{}`$ To verify this space discreteness experimentally one has to examine physical system in its vacuum state and at distances characterized by the length $`l_0`$. When system is in a rather mixed state
$$\mathrm{\Psi }(x,t)=\underset{S,\alpha }{}C(S,\alpha )\mathrm{\Psi }_S^{(\alpha )}(x,t)$$
the sharpness of the discrete structure disappears and space demonstrates usual continuous properties. So, this space discreteness is a quantum effect and depends on adelic quantum state.
Adelic wave function gives also a framework to investigate a new kind of phase, which may be called $`p`$-adic phase. In fact, (4.5) contains
$$\mathrm{\Psi }_p(x_p,t_p)=\chi _p[\alpha _p(\gamma (t_p)\gamma ^0)]\mathrm{\Psi }_p(x_p,0),$$
where $`\chi _p[\alpha _p(\gamma (t_p)\gamma ^0)]`$ presents a $`p`$-adic dynamical phase. This may be observed investigating the fine structure of interference phenomena.
At real distances which are very large in comparison with $`l_0`$, $`p`$-adic effects become hidden and adelic quantum mechanics reduces to the ordinary one. In such case we have to integrate $`|\mathrm{\Psi }(x,t)|^2`$ over $`p`$-adic components of adelic space. Since $`_{|x_p|_p1}𝑑x=1`$ and $`_\text{Q}_{\mathrm{}}|\mathrm{\Psi }_p(x_p,t_p)|_{\mathrm{}}^2𝑑x_p=1`$ we have
$$_{𝒜(S)\text{R}}|\mathrm{\Psi }_S(x,t)|_{\mathrm{}}^2𝑑x=|\mathrm{\Psi }_{\mathrm{}}(x_{\mathrm{}},t_{\mathrm{}})|_{\mathrm{}}^2dx_{\mathrm{}}\underset{pS}{}_\text{Q}_{\mathrm{}}|\mathrm{\Psi }_p(x_p,t_p)|_{\mathrm{}}^2𝑑x_p$$
$$\times \underset{pS}{}_\text{Z}_{\mathrm{}}\mathrm{\Omega }(|x_p|_p)dx_p=|\mathrm{\Psi }_{\mathrm{}}(x_{\mathrm{}},t_{\mathrm{}}|_{\mathrm{}}^2dx_{\mathrm{}}.$$
Hence, ordinary quantum theory may be regarded as an effective approximation of the more profound adelic one.
Acknowledgements
This work is partially supported by the Russian Foundation for Basic Research, Grant No 990100166, and leading scientific schools. We wish to thank colleagues from the Department of Mathematical Physics of the Steklov Mathematical Institute in Moscow for fruitful discussions.
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16. B. Dragovich, Integral Transforms and Special Functions 6, 197 (1998).
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# Geophysical constraints on mirror matter within the Earth
## I Introduction
Recently there has been considerable interest in the study of particle physics implications of the SuperKamiokande data on atmospheric neutrinos as well as the results of other neutrino experiments . The main purpose of such investigations has been to understand what theories could be responsible for the observed experimental features which give strong evidence for large angle neutrino oscillations. One of these theories is the Exact Parity Model (EPM) which introduces parity or “mirror” partners for all ordinary particles (except the graviton) and thus restores the parity invariance apparently broken by weak interactions . The Exact Parity Model predicts pairwise maximal mixing between ordinary and mirror neutrinos and provides a basis for interpretation of atmospheric neutrino and solar neutrino data .
An important question that arises naturally is whether or not the existence of mirror particles can lead to other observable consequences. In this work our objective is to find constraints on the possible concentration of mirror particles in the Earth. These constraints are a necessary step in the search for mirror world manifestations in terrestrial experiments. In particular, the presence of mirror matter in the Earth could lead to the regeneration of ordinary neutrinos from mirror neutrinos and consequently to the suppression of the day-night effect in the neutrino data.
Two main approaches to our problem are possible. First, one can trace the fate of the mirror particles starting from the early Universe epoch through the structure formation periods (galaxies, solar system and finally the Earth). Second, we can use geophysical data to get a more direct limit on the concentration of mirror matter in the Earth regardless of possible cosmological bounds.
It has been suggested that considerations based on the structure formation theory disfavour a significant presence of mirror matter in the Earth (see Blinnikov and Khlopov in ; Kolb, Seckel and Turner in ). However, as our knowledge of the structure formation is still incomplete, it is important to develop a geophysical approach as an independent, complementary tool of analysis exploiting the wide and rich variety of observational data accumulated in the Earth sciences.
This approach will be applicable not only to the specific EPM model, but also to any other theory predicting the existence of a new world of particles which couples to the ordinary matter only through gravitational interaction. An example is the shadow matter characteristic of superstring theories.
The plan of the paper is as follows. Section 2 summarises the main aspects of the “Standard Earth Model” called “Preliminary Reference Earth Model” (PREM). In Section 3 we analyse the possible effects of mirror matter within the context of PREM. Section 4 is devoted to the study of dynamical manifestations of mirror matter and comparison with gravimetric data. Finally, we present our conclusions in Section 5.
## II Preliminary Reference Earth Model (PREM)
The Preliminary Reference Earth Model (PREM) is a mechanical model of the average internal structure of the Earth based, mainly, on the analysis of seismological data. It gives the radial distributions of mechanical properties (such as density, elastic moduli, pressure, gravity and others) in the Earth’s interior.
The set of initial data used for constraining the model includes:
1. astronomic-geodetic data (radius, mass and moment of inertia of the Earth);
2. data on free oscillations and long-period surface waves (over 1000 eigenfrequencies are known);
3. body waves data ( $`10^6`$ arrival times have been registered).
Let us summarise the main analytical relations used in the construction of the model (for more details see e.g. ). The velocities of the elastic waves are given by
$$v_p(r)=\sqrt{\frac{K(r)+\frac{4}{3}\mu (r)}{\rho (r)}},v_s(r)=\sqrt{\frac{\mu (r)}{\rho (r)}},$$
(1)
where $`v_p`$ is the velocity of longitudinal waves, $`v_s`$ is the velocity of transverse waves, $`K`$ is the bulk modulus (or incompressibility) and $`\mu `$ is the shear modulus. These velocities can be found from seismological observations as functions of radius $`r`$. Measurements of wave velocities supply the ratios $`K/\rho `$ and $`\mu /\rho `$. To obtain $`\rho `$ independently, the Adams-Williamson equation,
$$\frac{d\rho }{dr}=\frac{G\rho (r)}{r^2(v_p^2(r)\frac{4}{3}v_s^2(r))}_0^r4\pi a^2\rho (a)𝑑a,$$
(2)
must be used. This equation expresses the condition of mechanical equilibrium between the gravitational attraction and the pressure due to elastic compression. $`G`$ is Newton’s constant, and the combination of the squared sound speeds in the denominator is called the seismic parameter $`\mathrm{\Phi }`$:
$$\mathrm{\Phi }(r)v_p^2(r)\frac{4}{3}v_s^2(r)=\frac{K(r)}{\rho (r)}.$$
(3)
Equation (2) is valid for a chemically homogeneous layer with adiabatic temperature gradient.
Further, the connections between the density profile and the profiles of pressure and gravity are given by
$$\frac{dP}{dr}=\mathrm{\Phi }(r)\frac{d\rho }{dr},$$
(4)
and
$$g(r)=\frac{G}{r^2}_0^r4\pi a^2\rho (a)𝑑a.$$
(5)
## III Static constraints
In this section, we will compute constraints on stationary mirror matter from PREM. A later section will consider dynamic manifestations of mirror matter.
### A Pedagogical warm-up exercise
In order to get a feeling for what would go wrong if a substantial amount of mirror matter were present in the Earth, we consider a simple but unrealistic scenario first. Suppose the Earth is actually two concentric spheres, one ordinary and one mirror. Suppose also, for simplicity, that the mass density ratio $`\rho _1(r)/\rho _0(r)`$ of mirror to ordinary matter is independent of radius. This is clearly unrealistic and we will relax this assumption in the next subsection.
In this case, alterations due to mirror matter of the PREM equations (1, 2, 3, 4, 5) can be accomodated either by rescaling Newton’s constant,
$$GG^{}=G(1+\frac{\rho _1}{\rho _0})$$
(6)
(throughout the paper quantities with index 0 will refer to ordinary matter while those with index 1 to mirror matter) or by keeping $`G`$ fixed and rescaling all of the other quantities used in PREM.
While the first procedure (rescaling $`G`$) seems to be the simplest, the second procedure reveals the physical modifications required in the presence of mirror matter more clearly and allows one to use PREM results directly.
From the formulas (1) and the Adams-Williamson equation (2) we deduce that the effective parameters can be defined as
$$\rho _{eff}=\rho _0+\rho _1,$$
(7)
$$P_{eff}=P_0(1+\frac{\rho _1}{\rho _0}),$$
(8)
$$\mu _{eff}=\mu _0(1+\frac{\rho _1}{\rho _0}),$$
(9)
$$K_{eff}=K_0(1+\frac{\rho _1}{\rho _0}),$$
(10)
$$g_{eff}=\frac{G}{r^2}_0^r4\pi a^2\rho _{eff}(a)𝑑a,$$
(11)
$$\mathrm{\Phi }_{eff}v_{p,eff}^2\frac{4}{3}v_{s,eff}^2=\frac{K_{eff}}{\rho _{eff}}=\frac{K_0}{\rho _0}=v_{p,0}^2\frac{4}{3}v_{s,0}^2\mathrm{\Phi }_0,$$
(12)
$$v_{p,eff}=v_{p,0}=\sqrt{\frac{K_{eff}(r)+\frac{4}{3}\mu _{eff}(r)}{\rho _{eff}(r)}},$$
(13)
$$v_{s,eff}=v_{s,0}=\sqrt{\frac{\mu _{eff}(r)}{\rho _{eff}(r)}}.$$
(14)
With these definitions the form of the Adams-Williamson equation does not change:
$$\frac{d\rho _{eff}}{dr}=\frac{\rho _{eff}g_{eff}}{\mathrm{\Phi }_{eff}}.$$
(15)
Also, the usual equation for the bulk modulus holds true in terms of the effective quantities:
$$\frac{d\rho _{eff}}{dP_{eff}}=\frac{\rho _{eff}}{K_{eff}}.$$
(16)
This demonstrates self-consistency of the rescaling procedure.
Now, for illustrative purposes, consider the case of a 50%–50% mixture of ordinary and mirror matter. The effective values for the density, incompressibility and pressure at the Earth centre can be taken directly from PREM data:
$$\rho _{eff}=13.1g/cm^3,K_{eff}=14.2Mb,P_{eff}=3.64Mb.$$
(17)
Correspondingly, the values for the ordinary matter (at the centre) are obtained by a factor of 2 rescaling:
$$\rho =6.55g/cm^3,K=7.1Mb,P=1.82Mb.$$
(18)
With these values, iron is definitely ruled out as the main component of the core and, therefore, the Earth as a whole. However, from independent evidence we know that iron is one the most abundant elements in the Earth. In addition, no other element (with significant abundance) can have the properties required by Eq. (18). For instance, silicon (at zero pressure) has a lower density than iron, but its incompessibility is greater than that of iron. Therefore the 50%–50% mixture of ordinary and mirror matter in the Earth is incompatible with the observational data encoded by the PREM model taking into account knowledge of terrestrial chemistry.
### B Modified Adams-Williamson equation
Consider now the realistic case where the mirror matter density does not follow the ordinary density. Then the density of the ordinary matter (indexed by 0) would obey the modified Adams-Williamson equation:
$$\frac{d\rho _0}{dr}=\frac{G\rho _0}{r^2(v_{p0}^2\frac{4}{3}v_{s0}^2)}_0^r4\pi a^2(\rho _0(a)+\rho _1(a))𝑑a.$$
(19)
Also, we have to require that the total mass of the Earth and the moment of inertia are equal to their observed values:
$$_0^R_{}4\pi a^2(\rho _0(a)+\rho _1(a))𝑑a=M_{},$$
(20)
$$\frac{8\pi }{3M_{}R_{}^2}_0^R_{}a^4(\rho _0(a)+\rho _1(a))𝑑a=I,$$
(21)
where
$$M_{}=5.974\times 10^{24}kg,I=0.3308.$$
(22)
Let us start with the simplest case: assume that the mirror matter forms a ball of uniform density $`\rho _1`$ with a radius equal to the radius of the inner core ($`R_{innercore}=1221kmR_{}/5`$). Within PREM the density profile in the inner core is given by the following function:
$$\rho _{PREM}(r)=\rho _{PREM}(0)q\frac{r^2}{R_{}^2},$$
(23)
where
$$\rho _{PREM}(0)=13.09g/cm^3,q=8.84g/cm^3.$$
(24)
Therefore, we look for the solution of the modified Adams-Williamson equation in the following form
$$\rho _0(r)=\rho _{PREM}(0)ϵ(q+\delta )\frac{r^2}{R_{}^2},$$
(25)
assuming that the mirror density $`\rho _1`$ is small compared to the Earth central density in PREM model $`\rho _{PREM}(0)`$. Solving the system of equations (19, 25) we find:
$$ϵ0.84\rho _1,\delta 0.9\rho _1.$$
(26)
Thus we see that if we wish to add to the Earth some amount of mirror matter then consistency with the equilibrium equation requires that the density of ordinary matter be decreased (as compared with PREM density) by approximately the same amount. Also, it can be shown that if Eq. (26) holds then the constraints due to the Earth mass and the moment of inertia are automatically satisfied as long as $`\rho _1<0.77g/cm^3`$.
Such a decrease of ordinary matter density can be achieved in one of two ways:
1) by changing the chemical composition of the core so that the new composition has lower density than the standard;
2) by decreasing the pressure so that the density of the standard core composition is lowered.
Before we consider these two possibilities let us review briefly the subject of the standard core composition (see e.g. ). Information about the chemical composition of the core is obtained by comparing the mechanical characteristics given by the PREM model with the properties of various substances under high pressure. In this way it has been established that the core characteristics are close but not equal to those of iron. There is sufficient evidence to conclude that some lighter element should be added to iron in order to satisfy the geophysical constraint on the core composition. Sulfur and oxygen appear to be the strongest candidates for that role although other elements (such as carbon, nitrogen etc) cannot be ruled out at present. As examples, cores containing 6–12% sulfur or 7–8% oxygen (by mass) have been proposed as possible compositions consistent with the PREM model.
In order to accomodate mirror matter according to the first method above, we should increase the admixture of light elements as compared with the standard levels. However, such an increase would typically raise the sound velocity in the core (for more details see ). Requiring that the sound velocity is equal within accuracy of about 0.3% to the observed value (see ) we can estimate that
$$ϵ0.18g/cm^3.$$
(27)
Next, in the second method we have to lower the pressure in order to obtain lower density of the (ordinary) matter. Lowering the pressure would lead to a decrease of the sound velocity and from the same requirement as above we can conclude that
$$ϵ0.06g/cm^3.$$
(28)
Then, using Eq. (26, 27, 28) we can obtain the upper bound on the mirror density in the inner core
$$\rho _1\frac{1}{0.84}max\{0.18,\mathrm{\hspace{0.33em}0.06}\}g/cm^3=0.21g/cm^3.$$
(29)
Translated to the upper limit on the ratio of the mirror mass to the total mass of the Earth, this becomes:
$$\frac{M_1}{M_{}}2.7\times 10^4.$$
(30)
### C Constraints from free oscillations of the Earth
Let us consider now the effect of changing the ordinary density on the Earth eigenfrequencies. Using Eq.(A2–A6) of Ref. and Eq.(41) of Ref. it can be shown that the relative change of period of the spheroidal $`{}_{0}{}^{}S_{0}^{}`$ mode as a result of changing the inner core density by $`ϵ`$ would be
$$\frac{\delta T}{T}\frac{0.026ϵ}{\rho _{PREM}(0)}.$$
(31)
Requiring that this shift of period be less than the fitting accuracy, 0.05%, we obtain a bound on the allowed change of the ordinary matter density in the inner core:
$$\frac{ϵ}{\rho _{PREM}(0)}1.9\times 10^2,ϵ0.26g/cm^3.$$
(32)
Using Eq. (26) this bound can be translated into a limit on the mirror matter density in the inner core:
$$\frac{\rho _1}{\rho _{PREM}(0)}2.3\times 10^2,\rho _10.3g/cm^3.$$
(33)
In terms of the total mass of the mirror matter $`M_1`$ we can rewrite Eq. (33) as
$$\frac{M_1}{M_{}}3.8\times 10^4.$$
(34)
Comparing Eqs. (30) and (34) we see that the difference between these upper limits is not very significant, although Eq. (30) is perhaps a less reliable estimate than Eq. (34) because of the incomplete knowledge of mechanical properties of various materials at high pressures characteristic for the Earth’s centre. For these reasons we interpret Eq. (34) as our final conservative upper bound on the mirror matter mass located inside the inner core of the Earth.
### D Arbitrary radius of the mirror matter ball
So far we have considered the case when mirror matter is contained completely inside the inner core of the Earth. To justify such an assumption one would need to know the detailed macroscopic properties of mirror matter (such as equation of state, chemical composition etc.); then using the condition of mirror matter equilibrium one could obtain the relation between the mirror matter and its radius. If we do not want to rely on such additional information than we have to regard the radius of the mirror ball $`R_1`$ as a free parameter of our model. Of course the resulting constraints will be weaker than they could be otherwise.
For simplicity we will restrict ourselves to the two characteristic values of $`R_1`$ in addition to the case $`R_1=R_{innercore}`$ considered before: first, $`R_1=R_{outercore}0.55R_{}`$, and, second, $`R_1=R_{lowermantle}0.89R_{}`$. The second choice is motivated by the fact that for radii larger than $`R_{lowermantle}`$ the Adams-Williamson equation is not valid anymore.
Also, we will assume that the mirror matter has a uniform density $`\rho _1`$ while the density of the ordinary matter differs from its PREM value by a radius-independent correction $`\delta \rho _0`$ for $`rR_1`$ and coincides with the PREM value for $`r>R_1`$. Although both these assumptions are clearly unrealistic, they nevertheless give us a self-consistent approximation scheme because both $`\rho _1`$ and $`\delta \rho _0`$ are small and therefore the effect of their radial dependence would be a second-order correction.
The precise details here depend on the chemical composition, the equation of state and other thermodynamic parameters for mirror matter. For instance, if we are given the equation of state for the mirror matter then (using some additional simplifying assumptions, such as chemical homogeneity, neglecting temperature etc.) we could find the density profile of the mirror matter (in particular, its central density) from the equation of mechanical equilibrium. However, our goal in this paper was to obtain limits on the mirror matter that would be independent of such particularities.
In any case, our method allows one to calculate the upper limit on the mirror matter mass for an arbitrary distribution of mirror matter $`\rho _1(r)`$.
As before, our constraints on the mass of the mirror matter will be based on 4 pieces of information:
1)mass of the Earth and its coefficient of inertia;
2)validity of the modified Adams-Williamson equation;
3)periods of Earth’s free oscillations;
4)velocities of elastic waves.
Generally speaking, points 1) and 2) tell us that the correction to the ordinary matter density should be approximately equal to the density of the mirror matter $`\delta \rho _0\rho _1`$. Next, using that equality we can compute the shift of the period for the $`{}_{0}{}^{}S_{0}^{}`$ mode and then obtain the upper limit on the mirror mass. From what follows it will be evident that we do not need to find an exact relation between $`\rho _1`$ and $`\delta \rho _0`$ as it would not significantly change the final constraints. Therefore, the validity of the modified Adams-Williamson equation (MAWE for short) can be analysed in a simpler manner than that of Sec. B.
As a criterion of validity of MAWE we can require that the actual mass of the Earth (that is, the ordinary mass plus the mirror mass) inside any radius $`rR_1`$ should be equal, with the accuracy of $`1\%`$ , to the PREM mass of the Earth (within the same radius):
$$w\left|\frac{M(r)}{M(r)_{PREM}}1\right|1\%,$$
(35)
where
$$M(r)=M(r)_{PREM}+_0^r4\pi a^2(\rho _1\delta \rho _0)𝑑aM(r)_{PREM}+\frac{4}{3}\pi r^3(\rho _1\delta \rho _0).$$
(36)
It is convenient to introduce the dimensionless ratio
$$f=\frac{\rho _1\delta \rho _0}{\overline{\rho }_{}},$$
(37)
where $`\overline{\rho }_{}=5.5g/cm^3`$ is the average density of the Earth. In terms of this quantity the condition (35) can be rewritten as
$$w=\left(\frac{r}{R_{}}\right)^3\left(\frac{M_{}}{M(r)_{PREM}}\right)\times f0.01,$$
(38)
for $`rR_1`$.
Next, we require that the total mass of the Earth equal the observed value with an accuracy of $`1.3\times 10^4`$ which translates into
$$\left(\frac{R_1}{R_{}}\right)^3\times f1.3\times 10^4.$$
(39)
Further, we have to demand that the coefficient of inertia of the Earth equals its PREM value $`I_{PREM}=0.3308`$ with an accuracy of $`\mathrm{\Delta }I/I3\times 10^4`$. In terms of $`f`$ this condition reads
$$\left(\frac{R_1}{R_{}}\right)^3\left(11.5\left(\frac{R_1}{R_{}}\right)^2\right)\times f3\times 10^4.$$
(40)
We observe that this condition does not give us an independent constraint because it is satisfied automatically as long as inequality (39) is fulfilled.
Finally, we compute the shift of the period for the $`{}_{0}{}^{}S_{0}^{}`$ normal mode of the Earth due to non-zero $`\delta \rho _0`$ using the same method as in Section C. In the case $`R_1=R_{outercore}`$ we find:
$$\frac{\delta T}{T}0.136\left(\frac{\rho _1}{\overline{\rho }_{}}f\right).$$
(41)
From inequality (39) we conclude that
$$f<8\times 10^4.$$
(42)
With these values of $`f`$ the criterion (38) of MAWE validity is clearly fulfilled for all radii $`rR_1`$. Next, requiring that the shift of period be less than the fitting accuracy (0.05%), we obtain the upper limit on the mirror matter density
$$\rho _10.025g/cm^3.$$
(43)
This translates into the following bound on the mirror matter mass:
$$\frac{M_1}{M_{}}7.4\times 10^4(forR_1=0.55R_{}).$$
(44)
Following the same procedure, in the case $`R_1=R_{lowermantle}`$ we obtain:
$$\frac{\delta T}{T}0.096\left(\frac{\rho _1}{\overline{\rho }_{}}f\right).$$
(45)
From (39) we find an upper bound on $`f`$:
$$f<1.8\times 10^4.$$
(46)
Again, the condition of MAWE validity, Eq. (38), is satisfied automatically with these $`f`$. Further, the upper limit on the mirror matter density $`\rho _1`$ becomes
$$\rho _10.03g/cm^3.$$
(47)
Correspondingly, the bound on the mirror matter is
$$\frac{M_1}{M_{}}3.8\times 10^3(forR_1=0.89R_{}).$$
(48)
Comparing our bounds (43) and (47) with the limit (29) we see that the latter limit is significantly weaker than the former two (although the limit (29) has been obtained for the inner core, it is also valid for the outer core because of the similarity of their chemical composition; we also would not expect substantial changes of this limit in the case of the lower mantle). Therefore we conclude that Eq. (34), Eq. (44), and Eq. (48) represent our final upper bounds on the mirror matter mass. Being the largest of the three, Eq. (48) can also be considered as the most conservative, radius-independent upper bound on the mirror matter mass in the Earth.
## IV Dynamical manifestations of mirror matter
In this section we are going to analyze constraints on mirror matter that follow from the precision measurements of the Earth’s gravitational field. Such constraints may arise if the mirror matter, for some reason, shifts away from the centre of the Earth.
The motion of the mirror matter in the Earth would be controlled mainly by the Earth’s gravity field; inside the core, in a first approximation this field grows linearly with the radius:
$$gkr,k=3.6\times 10^6s^2.$$
(49)
(The net tidal force exerted on the mirror matter by the Moon and the Sun is negligibly small as will be discussed later.) Consequently, the period of the mirror matter motion inside the Earth $`T_1`$ would be independent of radius and given by
$$T_1=\frac{2\pi }{\sqrt{k}}=3311s55min.$$
(50)
This period gives us the time scale for the variation $`\delta g`$ of gravitational acceleration at the Earth’s surface caused by the possible mirror matter motion.
To determine the amplitude of the gravity variation suppose that the amplitude of the mirror matter motion is $`h`$. Then the amplitude of the gravity variation is given by
$$\frac{\delta g}{g}\frac{M_1}{M_{}}\frac{h}{R_{}}.$$
(51)
The above equation holds exactly only for the circle on the Earth’s surface which lies in the plane of the mirror matter motion. For points outside this circle Eq.(51) can still be used for order-of-magnitude estimates. Requiring that the gravity variation should not exceed the observational limit, $`\delta g/g10^9`$ (see e.g. ) we obtain an upper bound on the amplitude of mirror matter motion:
$$h1.7\times \left(\frac{3.8\times 10^3M_{}}{M_1}\right)m.$$
(52)
What physical factors could lead to the off-centre shift of the mirror matter? Let us start by discussing the possible effect of Moon’s gravity. The two key quantities to be considered are the tidal torque and the net tidal force exerted on the mirror matter by the Moon. The effect of the tidal torque would be to slow down the spinning of the mirror matter, in analogy with the ordinary tidal torque that brakes the Earth’s rotation. The details of the effect depend on the poorly known characteristics of the mirror matter such as its angular velocity, elastic and dissipative properties etc.; we will not dwell on these.
On the other hand, the net tidal force and the corresponding off-centre shift can be estimated without the knowledge of mirror matter properties. First of all we note that if the Earth was spherically symmetric then there would be no net tidal force acting on the mirror “ball” placed in the Earth’s centre. However, if the oblateness of the Earth is taken into account then the net force at the centre is non-zero and thus the centre of the Earth is not an equilibrium position anymore. The new equilibrium position for the mirror matter can be found from the condition of balance between the net tidal force and the gravitational attraction of the mirror matter by the Earth.
To find the new equilibrium position it is convenient first to find the point in the Earth where the lunar tidal force vanishes (“the tidal centre”). Due to oblateness of the Earth the positions of the Earth’s centre of mass and the tidal centre are shifted relative to each other by a short distance $`b`$ (hereafter we ignore the $`18^{}`$ inclination of the Moon’s orbit relative to the Earth’s equatorial plane):
$$b=J_2\frac{3R_{}^2}{2R}171m,$$
(53)
where
$$J_2=\frac{CB}{M_{}R_{}^2}1.08\times 10^3$$
(54)
is the Earth’s dynamical oblateness, $`R3.84\times 10^8m`$ is the distance between the Moon and the Earth, $`C`$ and $`B`$ are the moments of inertia of the Earth with respect to the principal axes. Therefore the centre of mass of the mirror matter will move away from the Earth centre (and also away from the tidal centre) by the distance
$$h\frac{2Gmb}{kR^3},$$
(55)
where $`m`$ is the Moon’s mass. Inserting Eq.(53) into Eq.(55) we obtain
$$h\frac{3GmJ_2R_{}^2}{kR^4}8.1\times 10^6m.$$
(56)
The off-centre shift of mirror matter would create periodic variations of the gravity acceleration on the surface of Earth:
$$\frac{\delta g}{g}\frac{h}{R_{}}\frac{M_1}{M_{}}1.3\times 10^{12}\frac{M_1}{M_{}},$$
(57)
which is far beyond the observational limits. Thus we have shown that the effect of a net tidal force due to the Moon is negligible. A similar result holds for the solar effect:
$$\stackrel{~}{h}\frac{3GM_{}J_2R_{}^2}{k\stackrel{~}{R}^4}10^8m$$
(58)
where $`M_{}2\times 10^{30}kg`$ is the solar mass, $`\stackrel{~}{R}1.5\times 10^{11}m`$ is the distance between the Sun and the Earth. As expected, the effect of the net tidal force due to the Sun is even smaller than the lunar effect.
We now consider non-gravitational interactions that could possibly cause an off-centre shift of the mirror matter. Let us start by analysing the possible role of meteorites and meteor showers colliding with the Earth. Suppose that as a result of such a collision a momentum $`p`$ is transferred to the Earth. Then the mirror matter (assumed to be at rest in the centre) would receive an initial velocity $`u=p/M_{}`$ relative to the Earth. Therefore, the off- centre displacement would be equal to
$$h_{col}=\frac{u}{\sqrt{k}}=\frac{p}{M_{}\sqrt{k}}.$$
(59)
What could be the magnitude of $`p`$? The maximal velocity of a Sun-bound colliding object, relative to the Earth, is $`v_{max}73km/s`$. The heaviest meteorite found on the Earth has the mass $`m_{max}60ton`$. Inserting these values into Eq. (59) we obtain:
$$h_{col}\frac{m_{max}v_{max}}{M_{}\sqrt{k}}3.8\times 10^{13}m.$$
(60)
The variations of surface gravity acceleration caused by such displacements are many orders of magnitude beyond observational accuracy. Note that the impact of meteor showers would be much less than the estimate (60) since the total mass of even the most copious showers is significantly less than $`m_{max}`$.
In the case of still heavier meteorites which disperse after hitting the Earth the mass can be estimated only indirectly (see e.g. ). For instance, the meteorite that created the Arizona crater (with diameter of 1207 m and depth 174 m) had the estimated mass between 60 and 200 thousand tons. The mass of the Tunguska meteorite (fell in Siberia in 1908) was at least 1 million tons, its speed 30–40 km/s; if we insert these values into Eq. (59), we obtain
$$h_{col}10^8m;$$
(61)
corresponding values of $`\delta g/g`$ are still completely negligible (even in comparison with the variation of $`g`$ due to the Moon’s tidal effect, Eq. (57)).
Let us now turn to another possible mechanism of the mirror matter motion. Can earthquakes cause translational oscillations of mirror matter around the Earth centre? Large enough earthquakes are known to excite free vibrations of the Earth; the study of these vibrations has become one of the most important pieces of information about the Earth interior (for more details see e.g. ). These vibrations are classified into two categories:
—toroidal, in which only shear strain is present so that density is not perturbed; they are denoted by $`{}_{r}{}^{}T_{l}^{}`$;
—spheroidal, where both shear and volume deformations arise, denoted by $`{}_{r}{}^{}S_{l}^{}`$. The indexes $`r`$ and $`l`$ (for both $`S`$ and $`T`$ modes) have the same meaning as the radial and orbital quantum numbers of the hydrogen atom.
Toroidal modes do not lead to gravity perturbations so they cannot excite oscillations of mirror matter. On the other hand, spheroidal modes might cause the excitation of mirror matter oscillations through the gravitational coupling. Note that we should distinguish between two possible types of mirror matter oscillatons: a) bulk vibrations in which the centre of mass of mirror matter stays at rest in the Earth’s centre and b) translational oscillations where the centre of mass of mirror matter oscillates around the Earth’s centre in the Earth’s gravitational field. We cannot say much about the spectrum of bulk vibrations without knowing the detailed structure of mirror matter (i.e., its density, elastic and dissipative properties etc.); for this reason we leave them out of our consideration. On the contrary, the period of translational oscillations can be found and is given by Eq.(50).
Our next task is to find out if there are any spheroidal Earth eigenmodes that could resonate with translational oscillations of mirror matter. Note that here we deal with the case of parametric resonance so we need to look for the eigenperiod $`T_E=T_1/21655s`$ rather than $`T_E=T_1`$ (the parametric resonance in the case $`T_E=T_1`$ is weaker than for $`T_E=T_1/2`$). The closest such eigenmode is $`{}_{0}{}^{}S_{4}^{}`$ with the period of
$$T(_0S_4)=1545.6s.$$
(62)
In the time-varying gravitational field of $`{}_{0}{}^{}S_{4}^{}`$ mode the frequency of translational mirror matter oscillations also becomes time dependent according to the law
$$\omega ^2(t)=\omega _1^2(1+a\mathrm{cos}\gamma t),$$
(63)
where
$$\omega _1=\frac{2\pi }{T_1},\gamma =\frac{2\pi }{T(_0S_4)},a=\frac{\delta \rho }{\rho },$$
(64)
$`\delta \rho /\rho `$ is the amplitude of density variation in the $`{}_{0}{}^{}S_{4}^{}`$ mode.
The onset of parametric resonance is controlled by the quantity $`s`$ called amplification index (see e.g. ):
$$s=\frac{1}{2}\sqrt{\left(\frac{a\omega _1}{2}\right)^2ϵ^2},ϵ=\gamma 2\omega _1.$$
(65)
If the amplification index is real then the oscillation amplitude grows with time as $`\mathrm{exp}st`$. In the opposite case $`s^20`$ parametric resonance does not occur.
Using Eqs. (63,64,65) we can find $`s`$ to be
$$s=\frac{\omega _1}{4}\sqrt{\left(\frac{\delta \rho }{\rho }\right)^2(0.28)^2},$$
(66)
which is clearly imaginary. Thus we conclude that the condition for a parametric resonance is not satisfied and consequently there is no amplification of translational oscillations of the mirror matter.
It can be shown that the condition of parametric resonance with the $`{}_{0}{}^{}S_{2}^{}`$-mode takes the following form:
$$\frac{5}{24}a^2\omega _1<\gamma ^{}\omega _1<\frac{1}{24}a^2\omega _1,$$
(67)
where
$$\gamma ^{}=\frac{2\pi }{3233.25}s^1$$
(68)
is the frequency of the $`{}_{0}{}^{}S_{2}^{}`$-mode. One can see that the condition (67) is not satisfied and there is no resonance with the $`{}_{0}{}^{}S_{2}^{}`$-mode either.
## V Conclusion
We have investigated in detail geophysical constraints on the possible admixture of mirror matter inside the Earth. To this purpose, a method has been developed based on the Preliminary Reference Earth Model—the “Standard Model” of the Earth which describes its internal structure derived from the geophysical data in a systematic and self-consistent manner. If the density of the mirror matter is given, our method allows one to compute changes in various quantities characterising the Earth (such as its mass, moment of inertia, frequencies of its normal modes etc.). Comparing the computed and observed values of these characteristics, we can obtain for the first time the direct upper bounds on the possible concentration of the mirror matter in the Earth. In terms of the ratio of the mirror mass to the Earth mass these upper bounds range from $`3.8\times 10^4`$ to $`3.8\times 10^3`$ depending on the radius of the mirror matter ball. We then analyzed possible manifestations of mirror matter through the variations of the gravity acceleration on the Earth surface. These variations could arise as a result of an off-centre shift of the mirror matter due to several possible mechanisms such as lunar and solar tidal forces, meteorite impacts and earthquakes. Our estimates have shown that variations caused by these mechanisms are too small to be observed.
In this work we have been based on a standard premise that mirror matter interacts with ordinary matter only gravitationally <sup>*</sup><sup>*</sup>*Note that mirror matter can also couple to ordinary matter through photon—mirror-photon mixing . An analysis of this interesting possibility is beyond the scope of the present work.; we have not relied on any other specific assumptions about the mirror matter properties. Therefore our results are valid for other types of hypothetical matter coupled to ordinary matter by gravitation only; an example is shadow matter introduced in string theories. On the other hand, the use of equation of state and other macroscopic characteristic of mirror matter could lead to more severe constraints on the mirror mass inside the Earth.
###### Acknowledgements.
The authors are grateful to G.C.Joshi, R.Foot and B.H.J.McKellar for interesting discussions. This work was supported in part by the Australian Research Council.
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# FUSE Observations of O VI Absorption in the Galactic Halo
## 1 Introduction
In the fundamental paper “On a Possible Interstellar Galactic Corona,” Spitzer (1956) discussed the physical basis for the possible existence of a hot interstellar gas phase extending away from the Galactic plane into the halo. He also noted the diagnostic potential of the resonance doublet absorption lines of the Li-like ions of O VI, N V, and C IV for studying the coronal Galactic halo gas. Under conditions of equilbrium collisional ionization, these three ions peak in abundance at $`(3,2,1)\times 10^5`$ K, respectively (Sutherland & Dopita 1993).
Interstellar O VI observations were obtained by Jenkins (1978a, b) and York (1977) using the high resolution far-UV spectrometer aboard the Copernicus satellite. Those observations provided fundamental information about the hot interstellar gas in the disk of the Milky Way but were limited to stars brighter than V $``$ 7.0. The International Ultraviolet Explorer (IUE) and the Hubble Space Telescope (HST) allowed astronomers to study absorption by N V, C IV, and Si IV in the Galactic halo (Savage & de Boer 1979; Savage & Massa 1987; Sembach & Savage 1992; Savage, Sembach, & Lu 1997), but O VI was unobservable because windowless detectors and specially coated optics are required to measure the O VI $`\lambda \lambda `$1031.93, 1037.62 doublet. O VI has a special significance among the high ionization species because it is a sensitive indicator of collisionally ionized gas and is the least likely to be produced by photoionization from starlight given that 113.9 eV are required for the converison of O V to O VI.
Except for brief observing programs by the Hopkins Ultraviolet Telescope (HUT; Davidsen 1993), and the spectrographs in the Orbiting and Retrievable Far and Extreme Ultraviolet Spectrometers (ORFEUS; Hurwitz & Bowyer 1996; Hurwitz et al. 1998; Widmann et al. 1999; Sembach, Savage, & Hurwitz 1999), we have had a long hiatus in observing O VI. This has ended with the commissioning of the Far-Ultraviolet Spectroscopic Explorer (FUSE), a facility that has been specially designed to cover wavelengths from 905 to 1187 Å with an efficient 2-dimensional detector. The spectroscopic capabilities of FUSE and its in-flight performance are discussed by Moos et al. (2000) and Sahnow et al. (2000). This paper reports on initial FUSE results for O VI absorption in the Galactic halo.
## 2 FUSE Observations and Data Processing
The spectral integrations with FUSE were obtained in the time-tagged photon address mode between September and December 1999. The observations were obtained with the objects centered in the large 30″$`\times `$30″ aperture of the LiF1 channel and generally extended over spacecraft night and day. Most of the effective area in the O VI $`\lambda \lambda `$1031.93, 1037.62 wavelength region is provided by the spectra obtained with the LiF channels. We have restricted our analysis to the observations in the LiF1 channel.
To produce the final composite spectra in the region of O VI, we followed the basic data handling procedures discussed by Sembach et al. (2000). The spectra have a velocity resolution of approximately 25 km s<sup>-1</sup> (FWHM). The zero point of the wavelength scale in the vicinity of the O VI lines was determined with reference to nearby H<sub>2</sub> lines or Si II and O I lines. Sample FUSE spectra extending from 1020 to 1045 Å for Ton S210 and PG 0804+761 are shown in Figure 1.
## 3 Interstellar Analysis
The weaker O VI $`\lambda 1037.62`$ line is near strong absorption by C II $`\lambda 1037.02`$ and the H<sub>2</sub> (5–0) R(1) and P(1) lines at 1037.15 and 1038.16 Å. The stronger O VI $`\lambda 1031.93`$ line is usually relatively free of blending with other species since the nearby H<sub>2</sub> (6–0) R(3) and R(4) lines at 1031.19 and 1032.36 Å are often weak and well-separated in velocity from O VI ($`\mathrm{\Delta }v`$ = –214 and +125 km s<sup>-1</sup>, respectively). Blending from the HD (6–0) R(0) and R(1) lines at 1031.93 and 1032.36 Å is not a problem because the amount of molecular gas along these sight lines is small (Shull et al. 2000). For most paths through the Galactic halo, the contamination of the O VI $`\lambda 1037.62`$ line is severe. Therefore, we concentrate our attention on the O VI $`\lambda 1031.93`$ line in this Letter and use the $`\lambda 1037.62`$ line only to evaluate possible saturation in the $`\lambda 1031.93`$ absorption.
The O VI $`\lambda 1031.93`$ absorption is sufficiently broad that it is nearly fully resolved by FUSE. Therefore, we converted the observed absorption line profiles into measures of O VI apparent column density per unit velocity, N<sub>a</sub>($`v`$), through the relation N$`{}_{a}{}^{}(v)`$ \[ions cm<sup>-2</sup> (km s<sup>-1</sup>)<sup>-1</sup>\] = $`m_ec/\pi e^2`$ $`\tau _a(v)(f\lambda )^1`$=3.768$`\times `$10<sup>14</sup> $`\tau _a(v)(f\lambda )^1`$, where $`f`$=0.133 is the oscillator strength of the $`\lambda 1031.93`$ line (Morton 1991), $`\lambda `$ is the wavelength in Å, and $`\tau _a`$($`v`$) is the apparent absorption optical depth (see Savage & Sembach 1991). The continuum levels are well-defined by the smooth flux distributions provided by the AGNs. For cases where the O VI $`\lambda `$1037.62 line could be measured, we find the same values of $`\mathrm{log}`$N<sub>a</sub>($`v`$) near maximum absorption as those obtained for the stronger $`\lambda `$1031.93 line. Therefore, there is no evidence for unresolved saturated structure in the O VI absorption.
Values of the integrated apparent O VI column density are given in Table 1 along with errors including statistical and continuum placement uncertainties (Sembach & Savage 1992). Column densities and equivalent widths were integrated over the velocity range spanned by $`v_{}`$ to $`v_+`$ as listed in Table 1.
## 4 O VI Profiles
The O VI absorption profiles illustrated in Figure 1 and in Sembach et al. (2000) are often complex and trace a wide range of phenomena in or near the Milky Way. Those portions of the profiles between approximately –100 and +100 km s<sup>-1</sup> are likely tracing gas in the thick O VI halo of the Milky Way, which is the primary subject of this paper. However, an inspection of the O VI profiles reveals high velocity (100 $`<|v|<`$ 300 km s<sup>-1</sup> ) O VI absorption toward 7 of the 12 objects: Mrk 876, Mrk 509, PKS 2155-304, H 1821+643, NGC 7469, Ton S180, and Ton S210. Sembach et al. (2000) discuss the properties and possible origins of this high velocity O VI absorption. We comment on one case here. The O VI absorption toward H 1821+653 reveals local gas at $`v`$ = 0, gas above the Perseus spiral arm at $`v`$ = –70 km s<sup>-1</sup>, and gas in the outer warped region of the outer Milky Way at $`v`$ = –120 km s<sup>-1</sup> (Oegerle et al. 2000). In calculating the column densities listed in Table 1, the velocity limits were set to exclude high velocity O VI except for H 1821+643, where the value of N(O VI) includes the outer Galaxy gas absorption extending to $`v`$ = –160 km s<sup>-1</sup>. This high velocity O VI absorption is likely produced by differential Galactic rotation, which causes the absorption to extend to large negative velocities.
## 5 The Distribution of O VI in the Halo
Total column densities of O VI through the Galactic halo coupled with estimates of the mid-plane space density of O VI, $`n_0`$, can be used to obtain information about the stratification of O VI away from the plane of the Galaxy. We assume an exponential gas stratification with $`n(|z|)`$ = $`n_0\mathrm{exp}(|z|/h)`$, where $`h`$ is the O VI scale height. It then follows that the O VI column density perpendicular to the plane for an object with latitude $`b`$ is given by N(O VI)$`\mathrm{sin}|b|`$ = $`n_0h[1exp(|z|/h)`$\] , where N(O VI) is the line-of-sight column density to an object at a distance $`|z|`$ away from the plane. For extragalactic objects, where $`|z|h`$, N(O VI)$`\mathrm{sin}|b|`$ = $`n_0h`$.
Values of $`\mathrm{log}`$\[N(O VI)$`\mathrm{sin}|b|`$\] toward 11 extragalactic objects observed by FUSE are listed in Table 1 along with the ORFEUS value for 3C 273 from Hurwitz et al. (1998). The large value (14.80) for the 3C 273 sight line is likely a consequence passing through Radio Loops I and IV and the North Polar Spur. Such structures are local examples of the regions that likely feed hot gas into the halo. Since the 3C 273 sight line is a special case, we do not include it in our synthesis of general conclusions below. Without 3C 273, the median $`\mathrm{log}`$\[N(O VI)$`\mathrm{sin}|b|`$\] in our sample is 14.21 with a spread of 0.4 dex. The irregular nature of the distribution is highlighted by the low value of 13.80 found toward VII Zw118 which lies $`14`$ degrees from PG 0804+761 and has $`\mathrm{log}`$\[N(O VI)$`\mathrm{sin}|b|`$\] = 14.21. These irregularities must be considered when estimating the O VI scale height.
From the Copernicus observations of O VI absorption toward hot stars, Jenkins (1978b) estimated a mid-plane density $`n_0`$ = $`2.8\times 10^8`$ cm<sup>-3</sup>. However, that estimate assumed a small scale height (0.3 kpc) for the O VI absorbing layer, based on the limited data available at the time. We now find that the scale height must be about 10 times larger. Accordingly, we have obtained a new estimate of $`n_0`$ from the Copernicus O VI survey that is appropriate for a large scale height O VI absorbing layer with h $`>`$ 2 kpc. The result is $`n_0`$ = $`\mathrm{\Sigma }`$ N(O VI) / $`\mathrm{\Sigma }`$$`r_0`$ = $`2.0\times 10^8`$ cm<sup>-3</sup>, where the reduced distance $`r_0`$ = $`h`$\[$`1\mathrm{exp}(|z|/h)`$\]$`\mathrm{csc}|b|`$, compensates for the small reduction in density away from the Galactic plane. This value for $`n_0`$ was obtained when all of the upper limits for the Copernicus O VI column densities were included at their stated values. If zero is substituted for these cases, a value of $`n_0`$ that is only 5% lower is obtained. This value of $`n_0`$ has not been adjusted to allow for the fact what we live in the Local Bubble. Shelton & Cox (1994) have reanalyzed the Copernicus O VI measurements and have estimated that the mid-plane O VI density beyond the Local Bubble is $`n_0`$ = $`(1.31.5)\times 10^8`$ cm<sup>-3</sup> for an O VI absorbing layer with h $``$ 3 kpc.
For an O VI mid-plane density of $`2.0\times 10^8`$ cm<sup>-3</sup>, we obtain the scale heights listed in Table 1 from the simple relation h = N(O VI)$`\mathrm{sin}|b|/n_0`$. Ignoring 3C 273, the values range from 1.0 to 7.0 kpc. We find median and average values of 2.6 and 2.9 kpc, respectively. In deriving this average and in the subsequent calculations we treat the upper limit for PG 0052+251 as a detection. If we adopt the Shelton & Cox (1994) mid-plane density estimate of $`1.4\times 10^8`$ cm<sup>-3</sup> and reduce the extragalactic column densities by $`1.5\times 10^{13}`$ cm<sup>-2</sup> to remove the Local Bubble contribution, we obtain median and average O VI scale heights of 3.5 and 4.0 kpc. Our incomplete knowledge of the mid-plane density introduces a 35% systematic uncertainty in the derivation of the O VI scale height. Another source of uncertainty involves the irregular distribution of the gas. Edgar & Savage (1989) devised an analysis procedure for estimating scale heights that accounts for the irregular distribution. The analysis includes a logarithmic patchiness parameter, $`\sigma _p`$, that is added in quadrature to the observed logarithmic errors in the column densities. The value of $`\sigma _p`$ is varied until the minimized reduced chi square, $`\chi _\nu ^2`$(min), of the scale height fit is acceptable. Using n<sub>0</sub> = $`2.0\times 10^8`$ cm<sup>-3</sup> and the 11 FUSE values of N(O VI) sin$`|`$b$`|`$, we obtain $`\chi _\nu ^2`$(min) = 1.0 for h = 2.7 kpc and $`\sigma _p`$ = 0.21 dex. Adopting this value of $`\sigma _p`$, we can then estimate the 1$`\sigma `$ error in the scale height by determining the values of h where $`\chi ^2`$ = $`\chi ^2`$(min) + 1.0. The final result is h(O VI) = 2.7$`\pm `$0.4 kpc, where the listed 1$`\sigma `$ errors do not include the additional 35% systematic uncertainty caused by the uncertain O VI mid-plane density and Local Bubble correction.
The O VI scale height of 2.7$`\pm `$0.4 kpc can be compared with the values of h = 5.1$`\pm `$0.7, 4.4$`\pm `$0.6, and 3.9$`\pm `$1.4 kpc for Si IV, C IV, and N V, respectively, determined by Savage et al. (1997) from HST and IUE observations. The more confined distributions of O VI and N V compared to C IV are also apparent in plots of N(C IV)/N(O VI) and N(C IV)/N(N V) versus $`|z|`$ toward objects in the disk and halo of the Galaxy. In each case there is a clear increase in the ratio from the disk to the halo, suggesting that C IV is more extended than O VI and N V. For N V the ratio increases by about a factor of 2 from low to high $`|z|`$ (see Fig. 6d in Savage et al. 1997). Spitzer (1996) noted that the value of N(C IV)/ N(O VI) increases from $``$0.15 in the disk to 0.9 for objects in the low halo with $`|z|`$ 1.5 kpc. Values of N(C IV) have been measured using HST for Mrk 509, PKS 2155-304, 3C 273, H 1821+643, and ESO 141-55 by Savage et al. (1997) and by Sembach et al. (1999). Combining these with the values of N(O VI) from Table 1, we obtain N(C IV)/N(O VI) = 0.58, 0.63, 0.45, 0.63 and 1.74, respectively, for the five extragalactic objects. The increase in this ratio by about a factor of 4 from the disk to the typical extragalactic halo sight line implies a large change in the ionization state of the highly ionized gas as a function of distance from the Galactic plane.
A major goal of the FUSE O VI program is to map out the distribution of O VI in the disk and halo of the Galaxy. Once we obtain FUSE O VI observations for a substantial number of disk and halo stars and for additional extragalactic objects it will be possible to improve on our intitial estimate of the extension of O VI into the Galactic halo. These future studies will allow us to determine if a plane parallel, exponentially stratified, and patchy layer is indeed the most appropriate description of the distribution of O VI.
## 6 The Origin of Highly Ionized Atoms in the Galactic Halo
Strong O VI absorption toward 10 of 11 extragalactic objects observed by FUSE implies the widespread existence of hot gas in the halo of the Milky Way as predicted by Spitzer (1956) and also supported by the earlier observations of C IV and N V with IUE and HST. The decreasing scale heights in the sequence Si IV, C IV, N V, to O VI provides information about the changing ionization state of the gas with distance from the Galactic plane. Differences from element to element in the destructive liberation of atoms from dust grains could also influence the relative behavior of the z distributions of the highly ionized atoms.
Reviews of the many theories for the origins of the highly ionized atoms in the ISM are found in Spitzer (1990,1996) and Sembach, Savage, & Tripp (1997). The three primary types of theories involve conductive heating (CH) where cool gas evaporates into an adjacent hot medium, radiative cooling (RC) where hot gas cools as it flows into the halo or down onto the disk, and turbulent mixing layers (TML) where hot and cool gas are mixed through turbulent entrainment (Slavin, Shull, & Begelman 1993). Table 2 of Spitzer (1996) provides a summary of the many models along with the references to the theoretical literature. The different models make specific predictions for the expected values of N(C IV)/N(O VI). The CH models are compatible with the value of N(C IV)/N(O VI) $``$ 0.15 found at low $`|z|`$, while the larger values of $``$0.6 found toward the extragalactic objects observed by FUSE are better explained by a combination of the RC and TML models.
Shull & Slavin (1994) developed a hybrid model for the highly ionized gas in the Galactic halo in order to explain the smaller scale height of N V compared to Si IV and C IV suggested by the IUE and HST observations available in 1994. In their model the highly ionized ions at low $`|z|`$ are produced mainly in isolated SNRs while those at high $`|z|`$ are mainly found in radiatively cooling superbubbles that break through the disk producing Rayleigh-Taylor instabilities and turbulent mixing layers. Possible support for the origin of the high ions at low $`|z|`$ in isolated SNRs follows from the detailed SNR modeling of Shelton (1998). Confirmation that the scale height difference (smaller scale heights for ions with higher ionization potentials) first seen for N V and C IV is also clearly present in the new FUSE O VI measurements suggests that such hybrid models offer substantial promise for explaining the origin of the highly ionized species in the Galactic halo. Another example of a hybrid model is that of Ito & Ikeuchi (1988) which includes the cooling gas of a Galactic fountain flow (Shapiro & Field 1976) to provide the hot collisionally ionized gas and photoionization from the extragalactic background (Hartquist, Pettini, & Tallant 1984; Fransson & Chevalier 1985) to assist in the production of Si IV and C IV. The ionizing photons might also be provided by hot white dwarfs (Dupree & Raymond 1983). The new observations with FUSE imply several processes may be required to achieve a more complete understanding of the origins of the low density highly ionized gas extending away from the Galactic plane.
This work is based on data obtained for the Guaranteed Time Team by the NASA-CNES-CSA FUSE mission operated by the Johns Hopkins University. Financial support to U. S. participants has been provided by NASA contract NAS5-32985.
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# Doubling the number of Be/X-ray binaries in the SMC
## 1 Introduction
In high mass X-ray binaries (HMXBs) a neutron star or black hole orbits a massive early-type star and accretes matter either via Roche-lobe overflow or from the stellar wind which powers the X-ray emission (for recent reviews see Nagase 1989, White et al. 1995, Bildsten et al. 1997). One divides the class of HMXBs according to the stellar type of the mass donor star into supergiant X-ray binaries with luminosity class I-II OB star and Be/X-ray binaries with luminosity class III-IV Be star companions. Be/X-ray binaries form the larger sub-group of HMXBs. Balmer emission lines in the optical spectrum and a characteristic infrared excess are attributed to the presence of circum-stellar material, probably forming a disk in the equatorial plane of the Be star. Be/X-ray binaries often show transient behaviour with two types of outbursts. X-ray outbursts repeating with the orbital period are most likely associated with the passage of the neutron star through the circum-stellar disk in an eccentric orbit while giant outbursts, often lasting longer than a binary period, probably arise from an expansion of the disk.
Currently about 100 HMXBs and candidates are known. Nearly one third were found in the Magellanic Clouds (MCs) from which the majority is located in the Small Magellanic Cloud (SMC, Coe 1999). Most of the Be/X-ray binaries in the SMC were discovered in recent years by X-ray missions like ASCA, BeppoSAX, ROSAT and RXTE (Nagase 1999). From 20 optically identified HMXBs in the SMC only one is securely associated with a supergiant system (the X-ray pulsar SMC X–1) and from 11 of the 19 Be/X-ray binaries X-ray pulsations were detected. Five additional X-ray pulsars are yet to be identified, but are most likely also Be systems. The location of such a large number of HMXBs at a similar distance makes the SMC ideally suited for statistical and in particular spatial distribution studies of the population of HMXBs in a galaxy as a whole.
Recent surveys to look for H<sub>α</sub> emission-line objects in the SMC were performed by Meyssonier & Azzopardi (1993, hereafter MA93) and Murphy & Bessell (1999, MB99). The survey of MA93 mainly covers the main body and eastern wing of the SMC and their catalogue lists 1898 emission-line stars. The catalogue of MB99 covers nearly all the area where ROSAT PSPC observations of the SMC are available (except the southern half of the most south-east observation) but is less sensitive (372 objects, partially in common with MA93). A main goal of MA93 and MB99 was to identify planetary nebulae in the SMC, however, the catalogues also contain Be stars. MA93 state that all three at the time of publication known B\[e\] supergiants which were covered by the survey were detected. Very few Be/X-ray binaries were known in the SMC until 1993 and it was not noticed that the Be stars proposed as optical counterparts for SMC X-3, 2E 0050.1–7247 and 2E 0051.1–7304 were listed in MA93 as emission-line stars (LIN 198, AzV 111 and AzV 138, respectively). A correlation of the larger sample of Be/X-ray binaries known today in the SMC shows that most of them are found in the catalogues of MA93 and MB99. In this paper we use the identification of X-ray sources with emission-line stars to propose new very likely candidates for Be/X-ray binaries in the SMC (Sect. 2). X-ray source catalogues of the SMC which we used for our correlations with the emission-line star catalogues were published by Wang & Wu (1992) based on Einstein IPC observations (Seward & Mitchel 1981, Inoue et al. 1983, Bruhweiler et al. 1987) and by Haberl et al. (2000, HFPK00) produced from ROSAT PSPC data. We also used a preliminary version of the ROSAT HRI catalogue of Sasaki et al. (2000b).
## 2 X-ray sources and H<sub>α</sub> emission-line stars
A cross-correlation of the 517 PSPC X-ray sources in the SMC region (HFPK00) with the catalogue of H<sub>α</sub> emission-line stars published by MA93 (1898 entries) yielded 46 possible optical counterparts (distance $`<`$ $`\sqrt{r_{90}^2+10^2}`$ to account for systematic uncertainties in X-ray and optical positions, where r<sub>90</sub> denotes the 90% statistical uncertainty of the X-ray position in arc seconds). An additional object correlating with a PSPC source was found within the catalogue of candidate emission-line objects of MB99. From this sample of 47 objects one coincides with a known supernova remnant, three with supersoft sources and ten with optically identified Be stars proposed as optical counterparts for the X-ray sources. Extending the search by using additional SMC X-ray sources such as from the ROSAT HRI catalogue of Sasaki et al. (2000b), the Einstein IPC catalogue of Wang & Wu (1992) and X-ray pulsars discovered by instruments on ASCA, BeppoSAX and RXTE yields another three emission-line stars identified with known Be stars. Our correlation results are summarized in Table 1 which is sorted in right ascension.
Columns 2-4 of Table 1 give the source numbers in the X-ray catalogues for ROSAT PSPC (HFPK00), ROSAT HRI (Sasaki et al. 2000b) and Einstein IPC (Wang & Wu 1992). For sources detected by ROSAT coordinates with statistical 90% error (from HFPK00 when detected by the PSPC or from Sasaki et al. 2000b when detected by HRI only) are given in columns 5-7. For the group of IPC sources which were not detected by ROSAT, the IPC coordinates and a 40″ error as published in WW92 is given. The three digits in column 8 denote the number of PSPC detections in the energy bands 0.1 - 0.4 keV, 0.5 - 0.9 keV and 0.9 - 2.0 keV. With a few exceptions most of the sources were detected mainly in the higher energy bands which indicates a hard X-ray spectrum. HFPK00 used count ratios in the different energy bands, the hardness ratios, for a spectral classification of the PSPC sources. However, very hard sources without detection in the lower energy bands were not classified in HFPK00 because of large errors on the hardness ratios. Column 9 lists the maximum observed X-ray luminosity for Be/X-ray binaries and candidates derived in this work. The values are selected from literature or computed from ROSAT count rates using the conversion factor 1.67$`10^{37}`$ erg s<sup>-1</sup>/cts s<sup>-1</sup> (see Kahabka & Pietsch 1996, hereafter KP96), typical for X-ray binaries with hard spectrum at the distance of the SMC. HRI count rates were converted to PSPC rates using a multiplication factor of 3.0 (Sasaki et al. 2000a). Luminosities derived from count rates are indicated with colon. For all given X-ray luminosities we assume a distance of 65 kpc to the SMC.
Column 10 of Table 1 lists the entry number of the nearest object in MA93 (MB99 in one case) and column 11 the MA93 classification type (2 = SNR; 5 = planetary nebula, PN; 9 = late type star). The distance between X-ray and optical position as listed in MA93 (MB99) is listed in column 12. B, V and R magnitudes found in the literature for identified sources are given in columns 13-15. When available B and R obtained from the USNO A2.0 catalogue are listed for the remaining sources. In the last column identifications are given together with references and new proposals are marked with ’?’ behind the source type.
From the total of 18 high mass X-ray binaries in the SMC with known Be star as proposed counterpart 13 are found in the emission-line catalogues of MA93 and MB99. For completeness the remaining five Be/X-ray binaries are also listed in Table 1.
The identification of most known Be/X-ray binaries with stars in the emission-line catalogues of MA93 and MB99 suggests that un-identified X-ray sources with emission-line star counterparts are most likely also Be/X-ray binary systems. Other objects like SSSs and SNRs can be recognized by their very soft X-ray spectrum (in contrast to the Be/X-ray binaries with hard spectrum) or by their X-ray source extent, respectively. In the following section we summarize the 18 X-ray sources optically identified as Be/X-ray binary. We then propose emission-line stars as likely Be counterparts for the five un-identified pulsars and in addition for 25 hard X-ray sources.
### 2.1 Supersoft sources
Three supersoft sources detected by ROSAT were identified with emission-line objects in the catalogue of MA93. Two of them are associated with planetary nebulae (PN) while the remaining one is identified with a symbiotic star in the SMC. More detailed information on the individual sources and finding charts with X-ray error circles can be found in Sect. 6.1.
### 2.2 Optically identified Be/X-ray binaries
For 19 X-ray sources in the SMC nearby Be stars were optically identified and proposed as counterparts, suggesting a Be/X-ray binary nature. Eighteen were covered by ROSAT observations and information on the individual sources is summarized in Sect. 6 where also finding charts are found. The HEAO source 1H 0103–762 was not observed with ROSAT (see KP96 and references therein). Also we do not include the HMXB candidates RX J0106.2–7205 (Hughes & Smith 1994) and EXO 0114.6–7361 in our summary. For RX J0106.2–7205 no optical spectrum from the suggested counterpart is published yet, which would confirm its proposed Be star nature. For EXO 0114.6–7361 Wang & Wu (1992) propose the B0 Ia star AzV 488 as counterpart, however, AzV 477, also a B0 Ia star is even closer to the X-ray position. Both candidates suggest a supergiant type of HMXB. It is remarkable, that together with the only other known supergiant HMXB SMC X–1, EXO 0114.6–7361 is located in the eastern wing of the SMC.
Fourteen of the identified Be/X-ray binaries were detected by the ROSAT PSPC and their X-ray properties can be found in HPFK99. AX J0051–722, SMC X–3 and RX J0058.2–7231 were detected by the ROSAT HRI and are listed in the catalogue of HRI sources in the SMC (Sasaki et al. 2000b). Only 2E 0051.1–7304 was not detected by ROSAT. Thirteen of the proposed Be star counterparts are listed in the catalogues of MA93 and MB99 and only for five X-ray sources the Be counterparts have no entry in MA93 and MB99 (AX J0049–729, SMC X–2, RX J0032.9–7348, RX J0058.2–7231 and XTE J0111.2-7317). RX J0032.9–7348 was not covered by the MA93 survey.
### 2.3 Optically unidentified X-ray pulsars
Five X-ray pulsars in the SMC were reported for which no optical identifications are published up to day. Four of them were detected by ROSAT PSPC and/or HRI, yielding more accurate positions (HFPK00, Sasaki et al. 2000b) and for the fifth case, XTE J0054–720, several ROSAT sources are found within the large RXTE error circle. In or very close to the ROSAT error circles emission-line objects from MA93 are found and we propose these as optical counterparts. Literature, finding charts and other information on the X-ray binary pulsars is presented in Sect. 6.
Most of the Be stars proposed as optical counterparts for X-ray sources in the SMC, as summarized in the previous section, are found as emission-line objects in the catalogue of MA93. This strongly supports that the unknown emission-line objects within the error circles of the unidentified pulsars are also Be star counterparts of the X-ray pulsars forming Be/X-ray binaries.
### 2.4 New Be/X-ray binary candidates
From the correlation of X-ray source and emission-line object catalogues 34 hard X-ray sources were found with an H<sub>α</sub> emission-line object as possible optical counterpart in the X-ray error circle (see Table 1). The 34 X-ray sources were investigated in detail to obtain more information which can help to identify the nature of the object. Finding charts and notes to the individual sources are compiled in Sect. 6.
Many sources were observed more than once by ROSAT and we looked for long-term time variability. In the case of the PSPC we used the 0.9 – 2.0 keV band because of higher sensitivity for hard sources like Be/X-ray binaries. To combine detections from the different instruments we convert HRI to PSPC count rates by multiplying with 3.0 and IPC to PSPC count rates by multiplying with 1.1 (appropriate for a 5 keV Bremsstrahlung spectrum with 4.3$`10^{21}`$ cm<sup>-2</sup> absorption column density). Given the uncertainties in the count rate conversions, variability is only treated as significant above a factor of 3. None of the sources was observed with sufficient counting statistics in order to perform a detailed temporal analysis on shorter time scales (within an observation) and to detect X-ray pulsations.
We discuss all un-identified X-ray sources with emission-line object in or close to the error circle in the following and indicate very promising candidates for Be/X-ray binaries with “Be/X?” in the remark column of Table 1.
### 2.5 Chance coincidences
To estimate the number of false identifications of X-ray sources with emission-line objects in MA93 we shifted the X-ray positions of the sources in an arbitrary direction and cross-correlated again with the MA93 catalogue. To get statistically more reliable results this was repeated with different distances between 1 – 10 arc minutes. For this purpose we used the PSPC catalogue which is most complete. After application of our selection criteria for accepting Be/X-ray binary candidates (likelihood of existence for the X-ray source $`>`$ 13, no other identification, distance $`<`$ $`\sqrt{r_{90}^2+10^2}`$) we find on average about seven expected chance coincidences between PSPC sources and emission-line objects in MA93. We emphasize that these are mainly caused by the PSPC sources with the largest position uncertainties. The PSPC sources in Table 1 with the largest errors on the X-ray position and therefore most likely chance coincidences are 99, 248, 295 and 404. Indeed three of them were rejected as Be/X-ray binary candidates due to the presence of other likely counterparts. Similarly PSPC sources 77 and 253 were disregarded. Other X-ray sources with large position errors in Table 1 are the IPC sources which were not securely detected with ROSAT. Also here two were not regarded as Be/X-ray binary candidates. For the 25 new Be/X-ray binary candidates we therefore estimate that about two to three may be misidentifications, most likely among those with position error r<sub>90</sub> $`>`$ 15″.
## 3 ASCA binary pulsar candidates
The 1st ASCA Catalogue of X-ray sources in the SMC was compiled by Yokogawa (1999). The sources were classified according to their hardness ratios and Be/X-ray binaries were detected as sources with hard X-ray spectrum in the ASCA energy band. From this classification Yokogawa (1999) proposed eight new Be/X-ray binary candidates (binary pulsar candidates, BPc). We correlated our list of Be/X-ray binary candidates with the ASCA catalogue and find for five of the eight BPc a likely counterpart within 2′ (the maximum ASCA position uncertainty). In addition one out of the nine (probably because of its low flux) unclassified ASCA sources also correlates with a PSPC Be/X-ray binary candidate. In Table 2 the Be/X-ray binary candidates with likely ASCA counterpart are summarized. The first three columns show ASCA source number, classification and observed X-ray luminosity (0.7 – 10 keV, but corrected for the distance of 65 kpc used throughout this paper) from Yokogawa (1999). The next two columns contain PSPC source number and distance between ASCA and PSPC position and the last two columns give the same for HRI sources. The X-ray luminosities observed by ASCA are generally a factor of 1.2 – 2.4 higher than the ROSAT values (Table 1) which may only partly be explained by X-ray variability. Finding a systematically higher luminosity with ASCA is probably caused by the different sensitive energy bands and/or different intrinsic source spectrum. In particular the PSPC count rate to luminosity conversion is very sensitive to the assumed absorption. In the case of ASCA sources 36 (PSPC 279) and 7 (PSPC 468) the ASCA/ROSAT luminosity ratio is much higher than average (4.2 and 17, respectively), indicating strong flux variability.
## 4 Optical identifications
An optical identification campaign of a selected sample of hard ROSAT PSPC sources from the catalogue of HFPK00 was started independently to the present work (Keller et al. in preparation). From three of the X-ray sources presented here, spectra were taken which in all cases revealed a Be star nature of the proposed counterpart from the MA93 catalogue. This confirms RX J0057.8–7207 (PSPC 136, Sect. 6.4.8) as Be/X-ray binary and also the proposed counterpart for the pulsar AX J0105–722 (PSPC 163, Sect. 6.3.2, see also Filipović et al. 2000a) as Be star. RX J0051.9–7311 (PSPC 424, Sect. 6.2.8) was independently identified by Schmidtke et al. (1999). These results can be taken as further evidence that the emission-line objects we propose for counterparts of X-ray sources are indeed Be stars.
## 5 The Be/X-ray binary population of the SMC
The spatial distribution of the SMC HMXBs including the new candidates from this work is shown in Fig. 1. Nearly all new candidates are located along the main body of the SMC where most of the optically identified Be/X-ray binaries are concentrated. Only one new candidate is found in the eastern wing near the supergiant HMXB SMC X–1, where already two other Be/X-ray pulsars are known. The distribution is not biased due to incomplete coverage, neither in the optical nor in X-rays and makes the strong concentration of Be/X-ray binaries in certain areas of the SMC more pronounced.
The X-ray luminosity distribution of Be/X-ray binaries and candidates in the SMC is compared to that of systems in the Galaxy in Fig. 2. To do this we intensively searched the literature on galactic Be/X-ray binary systems. We derived 31 galactic sources with L$`{}_{}{}^{\mathrm{max}}{}_{\mathrm{x}}{}^{}`$$`>10^{33}`$ erg s<sup>-1</sup> which are summarized in Table 3 (which should be mostly complete). The luminosity estimates of galactic systems are often hampered by uncertain distances and different energy bands of the observing instrument. This may cause luminosity uncertainties by a factor of $``$10 in some cases but should not change the overall distribution drastically.
The new candidates in the SMC mainly raise the number of Be/X-ray binaries with luminosities log(L$`{}_{}{}^{\mathrm{max}}{}_{\mathrm{x}}{}^{}`$) $`<`$ 35.5 (21 out of 24 are new candidates). This can easily be explained by the high sensitivity of the ROSAT instruments which allowed to detect Be/X-ray binaries in their low-state while most of the higher luminosity Be/X-ray binaries were discovered during outburst. X-ray luminosities derived from detectors sensitive at higher energies (typically 0.5 – 10 keV) might be up to a factor $``$2 higher than those derived from ROSAT count rates (see Sect. 3) which would shift the low-luminosity end in Fig. 2 by 0.3 dex to the right. However, such a shift would not change the overall distribution significantly. Recently, in the Galaxy several likely low-luminosity Be/X-ray binaries were discovered by BeppoSAX and ASCA (1SAX J1324.4-6200, Angelini et al. 1998; AX J1820.5-1434, Kinugasa et al. 1998; AX J1749.2-2725, Torii et al. 1998; 1SAX J1452.8-5949, Oosterbroek et al. 1999; AX J1700062-4157, Torii et al. 1999). ROSAT also contributed new low-luminosity systems (RX J0440.9+4431, RX J0812.4-3114, RX J1037.5-5647, RX J0146.9+6121, Motch et al. 1997). However, the high absorption in the galactic plane makes the detection of low-luminosity X-ray sources in the ROSAT X-ray band and their optical identification difficult. This might explain the smaller number of low-luminosity Be/X-ray binaries discovered so far in the Galaxy compared to the SMC and might suggest that the luminosity distribution of Be/X-ray binaries is very similar in the SMC and our Galaxy. In this case many more such systems are expected to be found in our Galaxy which would significantly contribute to the hard X-ray galactic ridge emission (Warwick et al. 1985).
Various authors have suggested the existence of a population of low-luminosity systems which are usually persistent X-ray sources showing moderate outbursts and long pulse periods (e.g. Kinugasa et al. 1998, Mereghetti et al. 2000), somewhat different to the high-luminosity systems with strong outbursts and shorter pulse periods. The SMC results suggest that the low-luminosity sources even dominate the Be/X-ray binaries in number. From the fact that about one third of the already identified Be/X-ray binaries is not listed in current emission-line catalogues even more such systems are expected to be found in the ROSAT X-ray source catalogues of the SMC. On the other hand some of them will be observed with higher maximum luminosity in future outbursts, but if they indeed form a class of low-luminosity Be/X-ray binaries like X Per, the outbursts are expected to be small changing the luminosity distribution immaterial.
There is a large difference in the number of OB supergiant HMXBs between the Galaxy and the SMC. In the SMC at most two such systems are identified (SMC X-1 and maybe EXO 0114.6–7361) resulting in an overall ratio of Be to supergiant X-ray binaries of more than 20. In the Galaxy this proportion is more of order 2 (12 supergiant systems in the Galaxy are listed in the reviews of White et al. 1995 and Bildsten et al. 1997). It is remarkable that the SMC supergiant HMXBs are all located in the eastern wing giving rise to a local Be/supergiant HMXB ratio similar to that in the Galaxy. In contrast no supergiant HMXB is known in the SMC main body making the difference more extreme. One possible explanation is a different star formation history. Be/X-ray binaries evolve from binary star systems with typical total mass of $``$20 M within about 15 My (van den Heuvel 1983) while the more massive supergiant HMXBs are formed on shorter time scales. The latter therefore would trace more recent epochs of star formation than the Be/X-ray binaries. The comparatively large number of Be/X-ray binaries in the SMC in this view suggests a burst of star formation about 15 My ago while relatively few massive early-type stars were born during the last few million years.
It is remarkable that the Large Magellanic Cloud (LMC) may also have experienced a burst of star formation about 16 My ago as was derived from optical photometric surveys by Harris et al. (1999). LMC and SMC resemble in the relative composition of their X-ray binary populations, both rich in HMXBs but very few old low-mass X-ray binaries (Cowley et al. 1999), suggesting a common star formation history triggered by tidal interaction during close encounters of LMC, SMC and Milky Way. However, according to present day modeling the last encounter occurred $``$0.2 Gy ago (Gardiner & Noguchi 1996), too early for the formation of the Be/X-ray binaries we see today in X-rays. Therefore, the event which caused the origin of the frequent SMC (and LMC?) Be/X-ray binaries remains still unclear. Also the different numbers of HMXBs detected in LMC and SMC relative to their total mass need to be explained.
## 6 Notes to individual sources
In the following notes X-ray source numbers refer to the catalogues of Wang & Wu (1992) for Einstein IPC, of HFPK00 for ROSAT PSPC and of Sasaki et al. (2000b) for ROSAT HRI. Finding charts with X-ray error circles are shown in Figs. 3, 4, 5, and 6, for SSSs, already identified Be/X-ray binaries, unidentified X-ray pulsars and new Be/X-ray binary candidates, respectively. The order within each group of sources follows Table 1 and for faster access to the table entry the running number is given together with source name.
### 6.1 Supersoft sources
7) RX J0048.4–7332: The SSS RX J0048.4–7332 was discovered by Kahabka et al. (1994) and identified as the symbiotic M0 star SMC 3 by Morgan (1992). This star is listed in MA93 as object 218 and classified as late type star, consistent with the spectral type determined by Morgan (1992). The accurate HRI position (source 23) confirms the identification of the PSPC source (512).
39) 1E 0056.8–7154: This SSS was discovered in Einstein data (Inoue et al. 1983) and was detected with ROSAT PSPC (source 47) and HRI (79). It coincides in position with the SMC planetary nebula N67 (Aller et al. 1987) which is listed as object 1083 in MA93.
41) RX J0059.6–7138: This very soft source was discovered by HFPK00 (PSPC 51) and proposed as new SSS due to its positional coincidence with the planetary nebula LIN 357 (1159 in MA93) in the SMC.
### 6.2 Optically identified Be/X-ray binaries
1) RX J0032.9–7348: Stevens et al. (1999) identified two Be stars within the PSPC error circle of RX J0032.9–7348, discovered by KP96 as variable source with hard X-ray spectrum. The a factor of $``$5 smaller error radius obtained from a different PSPC observation by HFPK00 (source 567), however, still contains both Be stars which are very close to each other (Fig. 4). The two stars were not covered by the survey of MA93.
9) AX J0049–729: Yokogawa & Koyama (1998a) reported X-ray pulsations in ASCA data of this source. Kahabka & Pietsch (1998) suggested the highly variable source RX J0049.1–7250 (KP96) as counterpart. Stevens et al. (1999) identified two Be stars, one only 3″ from the X-ray position and one just outside the error circle given by KP96. The revised position of PSPC source 351 in HFPK00 makes the more distant Be star further unlikely as counterpart (see Fig. 4). None of the two Be stars turns up in the list of emission-line objects of MA93 with the nearest entry (279) 58″ away.
13) AX J0051–733: Yokogawa & Koyama (1998b) discovered X-ray pulsations from this source in ASCA data. The X-ray source was detected in Einstein IPC, ROSAT PSPC and HRI archival data and the 18 year history shows flux variations by at least a factor of 10 (Imanishi et al. 1999). Cowley et al. identified already 1997 a Be star as optical counterpart of the ROSAT HRI source RX J0050.8–7316 (HRI 34) which is located within the ASCA error circle. Cook (1998) reported a 0.708 d period from this star using data from the MACHO collaboration. The source was also detected by the PSPC (source 444) and coincides with object 387 in MA93.
16) AX J0051–722: Corbet et al. (1998b) reported 91 s X-ray pulsations from ASCA observations of this pulsar which was originally confused with the nearby 46 s pulsar XTE J0053-724 in XTE data. AX J0051–722 was not detected by the PSPC. An HRI detection reduced the position uncertainty and Stevens et al. (1999) identified a Be star as likely optical counterpart. The X-ray source is found as source 37 in the HRI catalogue and the star is identical to the only emission-line object from MA93 (413) in the ASCA error circle.
19) RX J0051.9–7311: This X-ray source was detected by Cowley et al. (1997) during ROSAT HRI observations of Einstein IPC source 25 and identified with a Be star by Schmidtke et al. (1999). It is identical to PSPC source 424 and HRI 41. The Be star is found as object 504 in MA93.
20) RX J0051.8–7231: This source was reported as strongly X-ray variable by KP96 and is associated with the X-ray pulsar 2E 0050.1–7247 (Israel et al. 1997). Observed X-ray luminosities range between 5$`10^{34}`$ erg s<sup>-1</sup> and 1.4$`10^{36}`$ erg s<sup>-1</sup> (Israel et al. 1997). The star AzV 111 (object 511 in MA93) was proposed as counterpart for 2E 0050.1–7247 while Israel et al. (1997) identified another H<sub>α</sub> active star within their error circle of RX J0051.8–7231 which is larger than that of KP96. Also the position error given for the corresponding PSPC source 265 by HFPK00 is large. The detection of the source in the PSPC observation 600453p (used by KP96) where the source was bright was rejected by the semi-automatic analysis of HFKP99 because the detection was close to the support structure of the detector entrance window. A careful analysis (and using the latest processed data of 600453p) of the photons of the source in the detector frame shows, however, that it moved nearly parallel to the window support structure and that it was not affected by it. In Table 1 therefore the parameters derived from this PSPC observation are given. They confirm the results of KP96 with small error circle (see Fig. 4). Both AzV 111 and star 1 of Israel et al. (1997) are outside this error circle which, however, contains a Be star identified by Stevens et al. (1999). This star is found as object 506 in MA93 and is the most probable counterpart of RX J0051.8–7231.
22) SMC X–3: This long-known X-ray source was not detected by the ROSAT PSPC but is included in the HRI catalogue as source 43. The Be star counterpart (e.g. Crampton et al. 1978) corresponds to object 531 in MA93.
23) RX J0052.1–7319: Lamb et al. (1999) reported X-ray pulsations from the variable source RX J0052.1–7319 (KP96) found in ROSAT HRI and CGRO BATSE data. Israel et al. (1999) identified a Be star as likely optical counterpart. It is found in MA93 as object 552 and was detected as X-ray source by IPC (29), PSPC (453) and HRI (44). The strong X-ray variability by a factor of $``$200 between different HRI observations (Kahabka 2000) strongly supports the identification as Be/X-ray binary.
24) 2E 0051.1–7304: For this source, listed as entry 31 in the Einstein IPC source catalogue of Wang & Wu (1992), the Be star AzV 138 (Garmany & Humphreys 1985) was proposed as optical counterpart. AzV 138 corresponds to object 618 in MA93. 2E 0051.1–7304 was not detected in ROSAT observations.
25) RX J0052.9–7158: This source was detected as X-ray transient by Cowley et al. (1997) during ROSAT HRI observations of Einstein IPC source 32 (the largest circle in the finding chart of Fig. 4. Upper limits derived from PSPC observations imply flux variations by at least a factor of $``$350 (Cowley et al. 1997). The strong variability and the hard X-ray spectrum imply a Be/X-ray transient consistent with the suggested Be star counterpart (Schmidtke et al. 1999). The Be star is identical to object 623 in MA93. The X-ray source was detected by ROSAT (PSPC 94 and HRI 46, the HRI position is most accurate as indicated by the smallest error circle in the finding chart of Fig. 4) and is located near the edge of the error circle of XTE J0054-720. Due to the large position uncertainty of XTE J0054-720 it is, however, not clear if they are identical.
28) SMC X–2: The long known Be/X-ray binary SMC X–2 was caught in outburst with 0.4 cts s<sup>-1</sup> by the ROSAT PSPC (source 547, see KP96 and references therein). Another PSPC observation yielded an upper limit indicating X-ray variability of more than a factor of 670. Optical spectra of the Be counterpart were taken by e.g. Crampton et al. (1978). The Be star is located on the rim of the PSPC error circle (Fig. 4) and is not contained in the MA93 catalogue.
31) XTE J0055–724: X-ray pulsations from this source were discovered by RXTE (Marshall & Lochner 1998) and confirmed in a SAX observation (Santangelo et al. 1998). Santangelo et al. (1998) also report pulsations from archival ROSAT data reducing the positional uncertainty. Stevens et al. (1999) identified a Be star as optical counterpart which corresponds to object 810 in MA93 and which is inside the error circle of PSPC source 241 and HRI source 58.
38) RX J0058.2–7231: RX J0058.2–7231 was detected as weak HRI source by Schmitdke et al. (1999) and identified with a Be star. It is contained in the HRI catalogue (source 76) but not found in the PSPC catalogue of HFPK00. The Be star is not detected in the emission-line star surveys of MA93 and MB99.
40) RX J0059.2–7138: This transient X-ray pulsar with peculiar soft component in the X-ray spectrum was discovered by Hughes (1994) during an outburst with a 0.2 – 2.0 keV luminosity of 3.5$`10^{37}`$ erg s<sup>-1</sup>. The X-ray source was identified with a Be star by Southwell & Charles (1996) as star 1 in their finding chart which is identical to the emission-line object 179 in MB99.
42) RX J0101.0–7206: This source was suggested as X-ray transient by KP96 with a flux variability of at least a factor of 30. Stevens et al. (1999) identified a Be star as optical counterpart. Object 1 in their Fig. 1f corresponds to entry 1240 in MA93.
49) SAX J0103.2–7209: Israel et al. (1998) reported X-ray pulsations from this source consistent in position with the Einstein source 1E 0101.5–7225. They confirm the Be star suggested as counterpart for the Einstein source by Hughes & Smith (1994) as the only object in the SAX error circle showing strong H<sub>α</sub> activity. OGLE observations presented by Coe & Orosz (2000) confirm this. The Be star corresponds to object 1367 in MA93 and was also detected by PSPC (source 143) and HRI (101).
58) XTE J0111.2–7317: Chakrabarty et al. (1998a) reported X-ray pulsations found in RXTE data from this source located about 30′ from SMC X–1. Wilson & Finger (1998) confirmed the pulsations from CGRO BATSE data and Chakrabarty et al. (1998b) derived an improved position from ASCA data. Two Be stars were identified by Israel et al. (1999) within or near the ASCA error circle of 30″. The closer of the two was concluded as most likely counterpart of XTE J0111.2–7317 by Coe et al. (1999). This Be star has no counterpart in MA93. A week source with existence likelihood of 14.5 is found in the PSPC catalogue (446). The large error circle of 61″ overlaps with the ASCA one and includes the position of the Be star. There is an additional MA93 object (1731) within the RXTE error circle and the second Be star found by Israel et al. (1999) is identical to object 1747 in MA93 but both are outside the ASCA and PSPC confidence regions (see Fig. 4).
59) RX J0117.6–7330: Similar to the previous source X-ray pulsations were discovered from the X-ray transient RX J0117.6–7330 (Clark et al. 1997) in ROSAT PSPC and CGRO BATSE data (Macomb et al. 1999). Between two PSPC observations, about 8 months apart, the count rate diminished by a factor of 270. Clark et al. (1997) identified a Be star counterpart which is identical to object 1845 in MA93 and also within the error circle of X-ray source 506 in the SMC PSPC catalogue.
### 6.3 Optically unidentified X-ray pulsars
10) AX J0049–732: AX J0049–732 was discovered as X-ray pulsar by Imanishi et al. (1998). Filipović et al. (2000b) reported two hard X-ray point sources from the catalogue of HFPK00 as possible counterparts to the ASCA pulsar. They suggest one of them (PSPC source 427) as the more likely counterpart due to its identification with an emission-line object in MA93 (number 300).
25) XTE J0054–720: The position of this X-ray pulsar could only be determined to an accuracy of 10′ radius (Lochner et al. 1998). There are at least five X-ray sources detected by the HRI within that circle (labeled 1 through 5 in Fig. 5 which correspond to the catalogue sources 55, 50, 62, 59 and 46, respectively). Object 2, 4 and 5 were also detected by the PSPC (104, 157 and 94). The southern of the three (also detected by IPC, 36) is proposed as active galactic nucleus (AGN) by HFPK00 and the northern (PSPC 94, HRI 46, IPC 32?) was identified as Be/X-ray transient RX J0052.9-7158 (see Sect. 6.2.2). The Be star counterpart of RX J0052.9-7158 coincides with object 623 in MA93. It is not clear if this Be/X-ray binary is identical to the RXTE pulsar. A final identification requires the detection of pulsations from RX J0052.9-7158.
27) XTE J0053–724: Corbet et al. (1998a) discovered this pulsar and report a ROSAT source within the error box. The pulse period, originally confused with AX J0051–722, was clarified by Corbet et al. (1998b). HFPK00 give source 242 as likely counterpart of XTE J0053–724. A single emission-line object from MA93 (717) is found inside the intersecting error circles of IPC source 34 and the PSPC source.
35) AX J0058–720: X-ray pulsations from this source were discovered by Yokogawa & Koyama (1998b) in ASCA observations. The source was detected in archival Einstein IPC, ROSAT PSPC and HRI data which span 18 years and showed flux variations by more than a factor of 100 (Tsujimoto et al. 1999). This high variability already strongly suggests a Be/X-ray binary. A single emission-line object from MA93 (1036) is found within the PSPC error circle (source 114) which is also consistent with the HRI position (73). It is not clear whether IPC source 41 originates from the same X-ray source. It may also be associated with another emission-line object (1039 of MA93) closer to the IPC position or completely unrelated.
53) AX J0105–722: Yokogawa & Koyama (1998c) reported AX J0105–722 as X-ray pulsar. From several nearby objects in MA93 number 1517 is closest to the X-ray position of PSPC source 163. This PSPC source was identified as likely counterpart of the ASCA pulsar in an area of complex X-ray emission by Filipović et al. (2000a) combining the ROSAT X-ray and radio data. The star 1517 in MA93 is the northern and bluer component of a pair of stars close to the error circles of PSPC and HRI detection (110). The nearby IPC source 53, 77″ to the north-east is most likely associated to the SNR DEM S128 (Filipović et al. 2000a).
### 6.4 New Be/X-ray binary candidates
2) RX J0041.2–7306: HFPK00 classified PSPC source 404 as foreground star based on the hardness ratios. An emission-line object in the error circle is classified as planetary nebula by MA93 indicating a chance positional coincidence. This makes the identification with the bright star just outside the error circle most likely.
3) RX J0045.6–7313: This source (PSPC 436) was detected once in the 0.9 – 2.0 keV band of the PSPC. An emission-line object in the error circle suggests an Be/X-ray binary.
5) RX J0047.3–7239: The PSPC error circle of RX J0047.3–7239 (source 295) overlaps with that of IPC source 19. An emission-line object (168 in MA93 and classified as late type star) and two radio sources from the catalogue of Filipović et al. (1998) are located in the X-ray confidence region. A point-like radio source as counterpart would favour an AGN identification leaving the nature of RX J0047.3–7239 ambiguous.
6) RX J0047.3–7312: RX J0047.3–7312 (PSPC 434) is most likely identified with the emission-line star 172 in MA93. The fluxes derived from PSPC detections show a factor of nine variations, supporting that the X-ray source is a Be/X-ray binary. RX J0047.3–7312 is probably identical to IPC source 18, which showed an intensity within the range observed by the PSPC. It is also the likely counterpart of ASCA source 2 in Yokogawa (1999; see Sect. 3), an X-ray binary candidate detected with similar intensity.
8) RX J0048.5–7302: The emission-line object 238 in MA93 is the brightest optical object in the error circle of RX J0048.5–7302 (PSPC 392). A Be/X-ray binary is suggested.
11) RX J0049.5–7331: An HRI detection (source 28) with much improved X-ray position compared to the PSPC (source 511) confirms the identification with the emission-line object 302 in MA93. RX J0049.5–7331 is the probable counterpart of ASCA source 6 in Yokogawa (1999; see Sect. 3) further supporting the likely Be/X-ray binary nature.
12) RX J0049.7–7323: This source (PSPC 468) was detected once in the 0.9 – 2.0 keV band of the PSPC. An emission-line object in the error circle suggests an Be/X-ray binary. RX J0049.7–7323 is also the likely counterpart of ASCA source 7 in Yokogawa (1999), classified as X-ray binary candidate (see Sect. 3).
14) RX J0050.7–7332: RX J0050.7–7332 was only once detected by the PSPC (514) and the emission-line object in the error circle suggests a Be/X-ray binary identification.
15) RX J0050.9–7310: HRI (source 36) and PSPC (source 421) detections are consistent with the identification of RX J0050.9–7310 with the emission-line object 414 in MA93, suggesting a Be/X-ray binary.
17) RX J0051.3–7250: Two close emission-line objects suggest RX J0051.3–7250 (PSPC 349) as Be/X-ray binary, but make the identification ambiguous.
18) RX J0051.8–7159: The emission-line object 502 (MA93) found in the error circle of RX J0051.8–7159 (PSPC 99) is classified as late type star in MA93. An active corona of this star may be producing the X-ray emission. The large error circle contains, however, another bright object which could also be responsible for the X-rays. The nature of RX J0051.8–7159 remains therefore unclear.
21) WW 26: Two emission-line objects from MA93 are found near IPC source 26 (hardness ratio 0.51, WW92). Object 521 is located inside the error circle while 487 can not be completely ruled out as counterpart. No ROSAT detection could improve on the position. A Be/X-ray binary nature is suggested.
26) RX J0053.4–7227: A precise HRI position (source 48 at the rim of the error circle of PSPC 246) with the emission-line star 667 (MA93) as brightest object in the error circle makes RX J0053.4–7227 a likely Be/X-ray binary.
29) RX J0054.5–7228: The uncertainty in the position of RX J0054.5–7228 (PSPC 248) is relatively large and six emission-line objects from MA93 are found as possible counterparts to the X-ray source. It is therefore a likely Be/X-ray binary but the optical counterpart remains ambiguous.
30) RX J0054.9–7245: Precise ROSAT X-ray positions (PSPC 324 = HRI 57) include an emission-line star (809 in MA93) with typical Be star magnitudes as brightest object in the error circles. A factor of five X-ray flux variability (the source was bright during a HRI observation) strengthens the identification as Be/X-ray binary.
32) WW 38 = 2E 0054.4–7237: An emission-line object (904 in MA93) is found inside the error circle of IPC source 38 suggesting a Be/X-ray binary. The source was not detected by ROSAT.
33) RX J0057.2–7233: This weak PSPC source (270) was marginally detected once in the hard 0.5 – 2.0 keV band with a likelihood of 10.4. Unlike all other hard sources in Table 1 it was not detected in the 0.9 – 2.0 keV band and therefore is unlikely a Be/X-ray binary.
34) WW 40 = 2E 0055.8–7229: The error circle of IPC source 40 contains two emission-line objects from MA93. Object number 1021 is identified as Be star AzV 111 while 1016, located further north, is of unknown type. ROSAT detected an X-ray source inside the IPC error circle (HRI 71 and PSPC 117 with consistent positions) which, however, is located between the two emission-line objects. The relation between the ROSAT and the Einstein source and the emission-line objects is unclear. IPC, HRI and PSPC count rates are consistent within a factor of two, which may indicate that they come from the same X-ray source. However, the accurate ROSAT positions make an association with one of the nearby objects from MA93 unlikely.
36) RX J0057.8–7207: Again small error circles from ROSAT HRI (source 74) and PSPC (source 136) observations make the identification of RX J0057.8–7207 with an emission-line star (1038 in MA93) very likely. PSPC detections with factor of eight different intensities and an HRI detection during an X-ray bright state which increases the variability to a factor of about 37, make a Be/X-ray binary nature highly probable.
37) RX J0057.9–7156: Be/X-ray binary candidate from positional coincidence of PSPC source 87 with emission-line object 1044 in MA93 which shows typical optical brightness.
43) RX J0101.3–7211: PSPC detections of this source (PSPC 159 = HRI 95) indicate flux variations by at least a factor of 15 and the source was not detected in other observations (upper limit a factor of 100 below the maximum count rate). This high variability and the presence of an emission-line star (1257 in MA93) in the small X-ray error circles likely exclude any other explanation than a Be/X-ray binary. It also is the likely counterpart of ASCA source 27 in Yokogawa (1999; see Sect. 3), an X-ray binary candidate.
44) RX J0101.6–7204: Two accurate positions from HRI (source 96) and PSPC (source 121) observations suggest the identification of RX J0101.6–7204 with object 1277 in MA93. The factor of three variability supports a Be/X-ray binary nature of RX J0101.6–7204 which is probably identical to the IPC source 46 in WW92.
45) RX J0101.8–7223: RX J0101.8–7223 (PSPC 220 = HRI 97) shows X-ray flux variations of a factor of three. The emission-line star 1288 (MA93) exhibits magnitudes typical for a Be star in the SMC and is located near the overlapping area of HRI and PSPC error circles. We suggest RX J0101.8–7223 as Be/X-ray binary as it is also the probable counterpart of ASCA source 28 in Yokogawa (1999; see Sect. 3), an X-ray binary candidate.
46) RX J0102.8–7157: This weak PSPC source (92) was only once marginally detected in the broad 0.1 – 2.4 keV band. The low detection likelihood of 10.5 and the non-detection in the hard bands indicates that it may not be real, or is at least not a hard source. A Be/X-ray binary nature is therefore unlikely.
47) WW 49: The IPC source 49 (WW92 give a hardness ratio of 0.21) contains a faint emission-line object (1357 in MA93) classified as planetary nebula. The spectral hardness of the IPC source is inconsistent with an SSS interpretation. The positional coincidence is likely by chance.
48) RX J0103.1–7151: This source was detected only once by the PSPC (source 77) and the lowest upper limit indicates variability by at least a factor of five, suggesting the detection of a single outburst. The emission-line object near the rim of the PSPC error circle is, however, the optically weakest (see Table 1), unusual in comparison with identified Be/X-ray binaries and candidates derived from this work. We therefore do not regard RX J0103.1–7151 as prime candidate for a Be/X-ray binary.
50) RX J0103.6–7201: Small error circles from HRI (source 105) and PSPC (source 106) observations make the identification with object 1393 in MA93 very likely. RX J0103.6–7201 shows variability by a factor of three between the ROSAT observations, consistent with a Be/X-ray binary.
51) RX J0104.1–7243: Two emission-line objects and a radio source from the catalogue of Filipović et al. (1998) close to RX J0104.1–7243 (PSPC 317) make the identification somewhat ambiguous. The most likely identification with emission-line star 1440 in MA93 suggests RX J0104.1–7243 as Be/X-ray binary.
52) RX J0104.5–7121: This source was not detected by the PSPC but the accurate HRI position (source 108) includes only the emission-line object 1470 from MA93 as bright object in the error circle. RX J0104.5–7121 is therefore very likely a Be/X-ray binary.
54) RX J0105.7–7226: An emission-line star (1544 in MA93) in the PSPC error circle (737) suggests RX J0105.7–7226 as Be/X-ray binary.
55) RX J0105.9–7203: A single bright object (the emission-line star 1557 in MA93) is found in the small PSPC error circle (source 120), which makes the identification of RX J0105.9–7203 as Be/X-ray very likely.
56) RX J0107.1–7235: The probable PSPC detection (279) of IPC source 56 improves the X-ray position and allows to identify it with the emission-line star 1619 in MA93. The source was a factor of 10 brighter during the Einstein observation and is also the likely counterpart of ASCA source 36 in Yokogawa (1999; see Sect. 3) detected with a factor $``$4 higher intensity. A Be/X-ray binary nature is likely.
57) RX J0109.0–7229: The emission-line object 1682 in MA93 is classified as planetary nebula. X-ray sources associated with planetary nebulae appear as SSSs which is not compatible with the hard spectrum of RX J0109.0–7229 (PSPC 253). The positional coincidence of RX J0109.0–7229 is therefore by chance and the nature of the X-ray source is unclear.
60) RX J0119.6–7330: This source (PSPC 501) was detected once in the 0.9 – 2.0 keV band of the PSPC. An emission-line object in the error circle suggests a Be/X-ray binary.
## 7 Summary
We reviewed the identification of eighteen known Be/X-ray binaries in the SMC and found that thirteen of them are listed in emission-line object catalogues of Meyssonier & Azzopardi (1993) and Murphy & Bessell (1999). From a general correlation of SMC X-ray source and H<sub>α</sub> emission-line object catalogues we propose optical counterparts for the five optically unidentified X-ray pulsars and present 25 new Be/X-ray binary candidates together with their likely optical counterparts. This more then doubles the number of know high mass X-ray binary systems in the SMC.
###### Acknowledgements.
The ROSAT project is supported by the German Bundesministerium für Bildung, Wissenschaft, Forschung und Technologie (BMBF/DLR) and the Max-Planck-Gesellschaft. The finding charts are based on photographic data obtained using The UK Schmidt Telescope. The UK Schmidt Telescope was operated by the Royal Observatory Edinburgh, with funding from the UK Science and Engineering Research Council, until 1988 June, and thereafter by the Anglo-Australian Observatory. Original plate material is copyright (c) the Royal Observatory Edinburgh and the Anglo-Australian Observatory. The plates were processed into the present compressed digital form with their permission. The Digitized Sky Survey was produced at the Space Telescope Science Institute under US Government grant NAG W-2166.
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# References
Backward emitted high-energy neutrons in hard reactions of p and $`\pi ^+`$ on carbon
A. Malki<sup>a</sup>, J. Alster<sup>a</sup>, G. Asryan<sup>c,b</sup>, D. Barton<sup>c</sup>, V. Baturin<sup>e,d</sup>, N. Buchkojarova<sup>c,d</sup>, A. Carroll<sup>c</sup>, A. Chtchetkovski<sup>e,d</sup>, S. Heppelmann<sup>e</sup>, T. Kawabata<sup>f</sup>, A. Leksanov<sup>e</sup>, Y. Makdisi<sup>c</sup>, E. Minina<sup>e</sup>, I. Navon<sup>a</sup>, H. Nicholson<sup>g</sup>, Yu. Panebratsev<sup>h</sup>, E. Piasetzky<sup>a</sup>, S. Shimanskiy<sup>h</sup>, A. Tang<sup>i</sup>, J.W. Watson<sup>i</sup>, H. Yoshida<sup>f</sup>, D. Zhalov<sup>e</sup>
<sup>a</sup>School of Physics and Astronomy, Sackler Faculty of Exact Sciences, Tel Aviv University. <sup>b</sup>Yerevan Physics Institute, Yerevan, Armenia. <sup>c</sup>Brookhaven National Laboratory. <sup>d</sup>Dept. of Physics, St. Petersburg Univ., St. Petersburg, Russia. <sup>e</sup>Physics Department, Pennsylvania State University. <sup>f</sup>Dept. of Physics, Kyoto Univ., Kyoto, Japan. <sup>g</sup>Mount Holyoke College. <sup>h</sup>J.I.N.R., Dubna, Russia. <sup>i</sup>Dept. of Physics, Kent State University.
Abstract
Beams of protons and pions of 5.9 GeV/c were incident on a C target. Neutrons emitted into the back hemisphere, in the laboratory system, were detected in (triple) coincidence with two emerging $`p_t>`$0.6 GeV/c particles. We present the momentum spectra of the backward going neutrons, which have the same universal shape observed in earlier (inclusive) reactions induced by hadrons, $`\gamma `$, $`\nu `$, and $`\overline{\nu }`$ beams. We also integrated the spectra and determined the fraction of the hard scattering events which are in coincidence with at least one neutron emitted into the back hemisphere, with momenta above 0.32 GeV/c. Contrary to the earlier measurements which found that only a small fraction (of the order of 10$`\%`$) of the total inelastic cross section for light nuclei was associated with backward going nucleons, we find that about half of the events are of this nature. We speculate that the reason for the large difference is the strong total center of mass (s) dependence of the hard reaction and short range nucleon correlations in nuclei.
PACS: 21.30.-x, 25.40.-h, 24.50.+g
There has been an intensive experimental program directed toward the systematic characterization of the emission of backward going nucleons from nuclei in collisions with high energy ($`>`$ 1 GeV) hadrons , real photons , virtual photons , neutrinos and antineutrinos . In these experiments the nucleons were emitted into angles larger than 90<sup>o</sup> in the laboratory system. Thus, the kinematical conditions were such that the observed backward nucleon could not result from a single two-body scattering of the incident particle with a nucleon at rest in the nucleus. Even though these experiments involved a large variety of interactions, energies and nuclei, a common universal spectrum of the backward nucleons was observed. This spectrum can be parameterized by the expression:
$$(E/p)\frac{d\sigma }{d(p^2)}=Ce^{Bp^2}$$
(1)
where $`E`$ and $`p`$ are the energy and momentum of the backward going nucleon.
For incident beams above about 1 GeV/c and for backward going nucleons above about 0.3 GeV/c the slope parameter $`B`$ was found to be independent of incident energy and beam type and target nucleus and only weakly dependent on the angle of the backward going particle ( see below for details). The absolute scale parameter $`C`$ depends on the nucleus and only weakly on the incident energy and projectile. The fraction of high energy (above the Fermi sea level) backward nucleons with respect to the total inelastic events for light nuclei (C, Ne) is of the order of 10$`\%`$ (see Table I).
In this paper we present the first measurement of high energy backward neutrons emitted from a nucleus in coincidence with two high-$`p_t`$ particles, ($`p_t>`$ 0.6 GeV/c). We find that, while the universal shape of the momentum spectra is maintained, the measured fraction of events with two high $`p_t`$ particles in coincidence with a backward going neutron is substantially higher than the ratios measured for the more inclusive reactions. A detailed description of our experiment follows.
We present the first results from a measurement which was performed during 1998 with the rebuilt EVA spectrometer at the AGS accelerator of Brookhaven National Laboratory. The spectrometer consisted of a super-conducting solenoidal magnet operated at 0.8 Tesla (see Fig. 1). The scattered particles were tracked by 4 sets of 4-layer straw tube cylindrical drift chambers (not shown) which surrounded the beam axis cylindrically. These straw tubes measured the transverse momenta of the particles by drift time and the polar scattering angle by charge division. Details on EVA spectrometer straw tube system are given in refs. . The major change to the spectrometer in addition to the improved performance of the solenoid and straw tubes was the installation of two new neutron counter arrays which increased the acceptance by a factor of 2.5 over the 1994 configuration .
At a momentum of 5.9 GeV/c the beam consisted of about 40$`\%`$ protons and 60$`\%`$ pions, identified by a sequence of two differential Cerenkov counters. The beam entered along the symmetry axis ($`z`$) of the magnet with an intensity of $``$1$`\times 10^7`$ particles over a one second spill, every 3 seconds. A scintillator hodoscope in the beam served as timing reference. Three solid targets consisting of various combinations of CH<sub>2</sub> and C were placed on the $`z`$ axis, separated by about 20 cm. They were 5.1$`\times `$5.1 cm<sup>2</sup> wide and 6.6 cm long in the $`z`$ direction. Their positions were interchanged several times at regular intervals. Some of the runs were with three C targets and some with two C targets and one CH<sub>2</sub> target.
As indicated in Fig 1, we triggered the spectrometer on two positively charged particles which emerged from the downstream end of the solenoid at polar angles of (27.5$`\pm `$3)<sup>o</sup> which corresponds to about 90<sup>o</sup> in the pp center of mass. For this analysis, events with two particles with a $`p_t>`$ 0.6 GeV/c which originated from one of the C targets were selected. The trigger required that one particle go to the left of the beam, and the other to the right. In addition we required that there were no additional charged tracks in the straw tubes. The polar angle coverage of the inner straw tube cyclinder extended from about 10<sup>o</sup> to 150<sup>o</sup>. Three scintillator arrays measured the direction and energy of neutrons, in coincidence with these two particles.
In Fig. 1 we present a schematic picture of the setup. We show the magnet of the EVA spectrometer and the positions of the targets. Below the targets we placed a series of 12 scintillator bars (ARRAY 1 in Fig. 1) covering an area of 0.6$`\times `$1.0 m<sup>2</sup> and 0.25 m (2 layers $`0.125`$ m each) deep. They spanned a polar angular range of 84<sup>o</sup> to 110<sup>o</sup> and an azimuthal range from 165<sup>o</sup> to 195<sup>o</sup>. A similar array of 16 scintillator bars (ARRAY 2), covering an area of 0.8$`\times `$1.0 m<sup>2</sup> and 0.25 m deep, spanned a polar angular range of 110<sup>o</sup> to 132<sup>o</sup> and the same azimuthal range as the first one. Each individual counter in these two arrays was 10$`\times `$12.5$`\times `$100 cm<sup>3</sup> in size and had a 5.1 cm photomultiplier at each end. The third array (ARRAY 3) was constructed from one layer of eight 10$`\times `$25$`\times `$100 cm<sup>3</sup> counters with a 12.7 cm photomultiplier at each end. This third array covered an area of 2.0$`\times `$1.0 m<sup>2</sup> and spanned a polar angular range of 72<sup>o</sup> to 120<sup>o</sup> and an azimuthal range from 120<sup>o</sup> to 150<sup>o</sup>. A set of veto counters (not shown in Fig 1.) served to discriminate against charged particles. Lead sheets (not shown in Fig 1.) with a total thickness of 1.7 radiation lengths were placed in front of the veto counters in order to reduce the number of photons entering the time of flight (TOF) spectrum. All counters were set to an electron equivalent detection threshold of 1 MeV by fixing a discriminator at the Compton edge of a <sup>60</sup>Co gamma source. This procedure is similar to the description in ref. . The detection efficiency was determined by the Monte Carlo method described in ref . The efficiencies depend on the neutron momentum and they ranged from about 30$`\%`$ at 0.15 GeV/c to about 15$`\%`$ at 0.5 GeV/c, for a typical single counter. We considered only neutrons above 0.15 GeV/c to avoid large uncertainties in the efficiency calculations of the neutrons at low momenta. A fraction of the neutrons gets absorbed on the trajectory from the target to the counters. The probability that a neutron was removed was calculated by assuming that the removal cross section was equal to the non-elastic cross sections in the materials . We assigned an uncertainty of 25$`\%`$ to these removal cross sections. The attenuation depends on the neutron momentum and the values ranged from about 35$`\%`$ at the low momenta to about 20$`\%`$ at the high momenta. The neutron momenta were determined from their TOF. A clearly identified peak due to remaining photons from the targets was used for calibration and to measure the timing resolution. That resolution was $`\sigma `$ 1 ns which corresponds to a momentum resolution of $`\sigma `$= 30 MeV/c at the highest momentum (0.5 GeV/c). We applied a cut in the TOF spectrum at 7 ns/m, keeping neutrons below 0.5 GeV/c, in order to eliminate the remaining photons.
In Fig 2. we present the measured invariant momentum spectra $`(E_n/p_n)\times \frac{N_3}{d(p_n^2)}`$ in arbitrary units for pion and proton incident beams, where $`E_n`$ and $`p_n`$ are the energy and momentum of the neutron detected in the backward hemisphere ($`90^o<\theta _n<130^o`$) and $`N_3`$ is defined below. We call $`N_2`$ the number of events with two high transverse momentum charged particles with $`p_t>0.6`$ GeV/c, each and no other charged particles seen in the detector. We applied software cuts to allow better determination of the target position from the track reconstruction and to obtain a better separation between incident protons and pions. The number of triple coincidence events which fulfill all the conditions of the $`N_2`$ events and, in addition, have a single neutron in the scintillator bars is indicated by $`N_3`$. Since the efficiency and attenuation corrections depend on the neutron momentum they were done event by event. The resulting spectra are plotted on a semi-logarithmic scale as a function of $`p_{n}^{}{}_{}{}^{2}`$. The error bars represent the statistical errors only. The curves above $`p_{n}^{}{}_{}{}^{2}>0.1`$ (GeV/c)<sup>2</sup> were fitted to a first order polynomial. The slope parameters ($`B`$ in equation 1) for the proton and pion incident beams are shown in the figure with their fitting error. For comparison we quote the slope parameter obtained from the p+C$``$ n+X data at 7.5 GeV/c. Their value at 119<sup>o</sup>, which lies within the range of our measurement, is B=13.0$`\pm `$0.8 (GeV/c)<sup>-2</sup>. The neutrino measurements report a value of $`B`$=9.5-10.7 (GeV/c)<sup>-2</sup> for protons emitted into the whole backward hemisphere with uncertainties ranging from 0.3 to 2 (GeV/c)<sup>-2</sup>. Ref. presents a compilation of other experiments at several different energies of hadron and photon beams incident on a variety of targets, including C and other light nuclei. For particles in the common angular range of 120<sup>o</sup>-150<sup>o</sup> of the different experiments, the slope parameters all lie within the range of 10-12.5 (GeV/c)<sup>-2</sup> with a typical error of 2 (GeV/c)<sup>-2</sup>. Our first conclusion is that the slopes we measured in this experiment agree, within the measured uncertainties, with the slopes obtained using hadrons, photons and neutrino beams of different energies incident on various targets. The angular ranges are not the same for all experiments but this does not modify the values of $`B`$ sufficiently to affect the conclusion (see ).
We also wish to compare our yield of backward scattered neutrons above 0.32 GeV/c to those of the other experiments. In Fig. 3 we show the (triple coincidence) backward yield per unit solid angle divided by the (double coincidence) two high-$`p_t`$ particle yield vs. the neutron angle. Each point in this figure contains 10-20 combinations of target positions and a neutron-counter (with the exception of the most forward angle point which includes 3 such combinations only). The errors in the figure are the statistical errors combined with up to 10$`\%`$ systematic errors due to software cuts, uncertainties in the neutron detection thresholds and uncertainties related to determining the exact geometrical positions of the counters. To obtain the yield into the backward hemisphere we fitted the ratios to a constant (see fig 3). For the proton induced reaction this gives a value of (7.4$`\pm `$0.4)$`\%`$ /sr and for the pion induced reaction the value of (6.5$`\pm `$0.6)$`\%`$/sr above $`90^o`$ and $`100^o`$, respectively where the ratios become constant. We used the parameters of the fit to extrapolate up to $`\theta _n=180^o`$, assuming that the ratio remains constant. The results for the pion and proton induced reactions are shown in Table 1 compared with the earlier mentioned experiments. Our results are given for the integration up to $`\theta _n=130^o`$ with the proper error as well as for the integration to $`\theta _n=180^o`$ which assumes the constant ratio out to that angle. Our second conclusion from this experiment is, that in this measurement the ratio is substantially larger (3-4 times) than for the more inclusive measurements. This is true even if we include only our yield out to $`\theta _n=130^o`$. Integrating out to 180<sup>o</sup> can only increase that ratio.
The earlier experiments have led to several theoretical interpretations. The models that have been discussed can be divided into two main classes. The first class deduces from the universality that the measurement must provide direct information on the nuclear ground state, especially on the high-momentum part of the wave function. These models assume that the projectile interacts with a single nucleon and that the backward yield of nucleons is due to the break-up of pre-existing two, or more, correlated nucleon clusters . One of the more extreme models in this class assumes that the high momentum of the struck nucleon is balanced coherently by the residual nucleus . A second class of models assumes that the backward yield is due to rescattering processes in the nucleus of the incident and outgoing particles. These initial and final interactions include true $`\pi `$-absorption and $`\mathrm{\Delta }`$-rescattering in the intermediate states.
The results of our experiment that the ratio of backward going neutrons above a momentum of 0.32 GeV/c in coincidence with two high p<sub>t</sub> protons, is substantially larger adds new information that has to be accomodated in the theoretical models. We speculate that the reason for the difference is the strong total center of mass (s) dependence of the hard reaction cross section and it’s sensitivity to the short range nucleon correlations in nuclei. The strong s dependence of the hard reaction (for example 1/s<sup>10</sup> for pp elastic scattering) selectively chooses the high momentum protons in the nuclei. Those protons most likely have a correlated partner at short range which are the neutrons that dominate the backward going yield . This speculation needs to be checked with detailed calculations.
We wish to thank Drs. L. Frankfurt M. Strikman and M. Sargsyan for their theoretical input. We are pleased to acknowledge the assistance of the AGS staff in building and rebuilding the detector and supporting the experiment, particularly our liaison engineers, J. Mills, D. Dayton, C. Pearson. We acknowledge the continuing support of D. Lowenstein and P. Pile. This research was supported by the U.S. - Israel Binational Science foundation, the Israel Science Foundation founded by the Israel Academy of Sciences and Humanities, the NSF grants PHY-9501114, PHY-9722519 and the U.S. Department of Energy grant DEFG0290ER40553.
Tables
Table 1
| incident/ | energy | target | integration | integration | Ref. | RATIO |
| --- | --- | --- | --- | --- | --- | --- |
| backward emitted | | | range | range | | |
| particle | (GeV) | | (GeV/c) | (Deg.) | | ($`\%`$) |
| p/n | 5.9 | C | 0.32-0.5 | 90-130 | a | 29.9$`\pm `$1.6 |
| $`\pi `$/n | 5.9 | C | 0.32-0.5 | 90-130 | b | 26.2$`\pm `$2.4 |
| $`\overline{\nu }`$/p | $`WB^c`$ | Ne | 0.2-0.7 | 90-180 | d1 | 7$`\pm `$1 |
| p/p | 3.66 | Ne | 0.2-0.7 | 90-180 | d2 | 8$`\pm `$3 |
| $`\overline{\nu }`$/p | $`WB^e`$ | Ne | 0.2-0.8 | 90-180 | f | 10$`\pm `$1 |
| $`\nu `$/p | | | | | | 12$`\pm `$4 |
| NC/p | | | | | | 10$`\pm `$1 |
| $`\overline{\nu }`$/p | $`WB^e`$ | Ne | 0.35-0.8 | 90-180 | g | 6.0$`\pm `$0.2 |
| $`\nu `$/p | | | | | | 8.5$`\pm `$0.3 |
| p/p | 0.64 | C | 0.311-0.54 | 90-180 | h | 6.0$`\pm `$1.5 |
a- This work. If we extrapolate to 180<sup>o</sup> under the assumptions in the text, we get a value of (46.5$`\pm `$2.5)$`\%`$.
b- This work. If we extrapolate to 180<sup>o</sup> under the assumptions in the text, we get a value of (40.8$`\pm `$3.7)$`\%`$.
c- A wide band $`\overline{\nu }`$ beam from 300 GeV/c incident protons.
d1- Results of ref. .
d2- Deduced in ref. from measurements of ref. to satisfy the same selection criteria as used in ref. .
e- A wide band $`\overline{\nu }`$ beam from 400 GeV/c incident protons.
f- ”(BP+2BP)/TOTAL” from Table I of ref.
g- ”1 Backward proton rate” from Table 1 of ref. .
h- Integral of I($`\theta _3`$) from Table 2 divided by $`\sigma _t`$ from ref.
Figure captions.
Fig. 1 The set up of the neutron counter arrays. Only the magnet of the spectrometer is shown with the position of the targets. The lead shield and veto counters above the neutron counter arrays are not shown.
Fig. 2 Proton and pion induced neutron invariant momentum spectra. The vertical axis is $`ln[(E_n/p_n)\times \frac{N_3}{d(p_n^2)}]`$, The horizontal axis is $`p_{n}^{}{}_{}{}^{2}`$. $`E_n`$ and $`p_n`$ are the energy and momentum of the neutron. $`N_2`$ is the number of events with exactly two charged particles, each with $`p_t>0.6`$ GeV/c, detected in the spectrometer. $`N_3`$ is the number of $`N_2`$ events that also have a single neutron entering the neutron counters. The neutron yield is corrected for the detection efficiency and attenuation. Above $`p_{n}^{}{}_{}{}^{2}>0.1`$ (GeV/c)<sup>2</sup> the points are fitted to a straight line to obtain the slope parameter defined in Eq. 1. The resulting slopes with the fitting errors are shown.
Fig. 3 The relative yield per solid angle $`\frac{dN_3}{N_2d\mathrm{\Omega }_n}`$ of backward going neutrons above 0.32 GeV/c as a function of the neutron angle. $`N_3`$ and $`N_2`$ are defined in Fig. 2 and the text. The data are for the proton and pion induced reactions. The lines represent fits to a constant which is used to estimate the total backward emission yield, see text.
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# A note on topological brane theories
## 1 Introduction
In the present note we consider a $`p`$-brane propagating on a $`(p+1)`$-dimensional background manifold. In this model all degrees of freedom can be gauged away locally. However there may still be non trivial non-local (topological) degrees of freedom. The motivation for the study of this model is twofold. First this theory is interesting by itself as a topological field theory , and we shall see that there is something one can learn about its canonical and BRST structures. Second the present theory may serve as a toy model for the study of general extended objects which play an important role in modern string theory. Thus the present model can give some insight into the more general case of $`p`$-branes propagating on background manifolds of dimension higher than $`(p+1)`$.
The $`p`$-brane theory describes the embedding of a $`(p+1)`$-dimensional world-volume into a $`d`$-dimensional space-time manifold $``$. The action is given by the volume of the embedded $`(p+1)`$-dimensional manifold
$$S[X]=Td^{p+1}\xi \sqrt{det\left(G_{\mu \nu }(X)_\alpha X^\mu _\beta X^\nu \right)},$$
(1)
where $`G_{\mu \nu }`$ is the metric in the $`d`$-dimensional space-time manifold. Throughout the paper we shall look at the case of Minkowskian signature, and we shall see that the results can be generalized to the Euclidean case in a straightforward way. It is assumed that metric $`G_{\mu \nu }`$ is not degenerate at any point of manifold $``$, $`det(G_{\mu \nu })0`$. The action (1) is called the Nambu-Goto action. The parameter $`T`$ is called tension and it is a direct generalization of the concept of mass to $`p`$-branes. The $`p`$-brane action can be equivalently written in the Polyakov form
$$S[X,h]=\frac{T}{2}d^{p+1}\xi \sqrt{h}\left[h^{\alpha \beta }_\alpha X^\mu _\beta X^\nu G_{\mu \nu }(X)(p1)\right],$$
(2)
where the $`h_{\alpha \beta }`$ transform as a world-volume metric ($`hdet(h_{\alpha \beta })`$) and play the role of Lagrange multipliers. The equivalence can be checked by varying this action with respect to $`h_{\alpha \beta }`$, solving the resulting equations of motions and putting the resulting solution for $`h_{\alpha \beta }`$ into equation (2). The result is the action (1). One should notice that within the equivalence between actions (1) and (2) the induced world-volume metric must be nondegenerate.
Now let us discuss the symmetries of the theory. There is a local diffeomorphism invariance for the action (1)
$$\delta X^\mu =\mathrm{\pounds }_\zeta X^\mu ,$$
(3)
where $`\mathrm{\pounds }_\zeta `$ stands for the Lie derivative along $`\zeta ^\alpha `$. For the Polyakov action (2) the transformation (3) should be supplemented by the appropriate transformation for the auxiliary world-volume metric $`\delta h^{\alpha \beta }=\mathrm{\pounds }_\zeta h^{\alpha \beta }`$. For the case $`p=1`$ there is an extra local symmetry $`\delta h^{\alpha \beta }=\mathrm{\Lambda }h^{\alpha \beta }`$ (Weyl rescaling). In addition both actions (1) and (2) are invariant under arbitrary diffeomorphisms on $``$, if $`G_{\mu \nu }`$ is transformed properly.
The local symmetry (3) allows one to choose locally the following gauge
$$X^\mu =\xi ^\mu ,\mu =0,1,\mathrm{},p,$$
(4)
which is usually called the static gauge. The existence of static gauge can be argued from the picture of embedding of a manifold into another one. Thus in the case of interest $`d=p+1`$ one can locally gauge away all degrees of freedom. However on a nontrivial background manifold $``$ one cannot do this globally and therefore there are nontrivial global (topological) degrees of freedom. One can also see that degenerate situations (when $`det\left(G_{\mu \nu }(X)_\alpha X^\mu _\beta X^\nu \right)=0`$) do not appear because of (4). In the static gauge the determinant of the induced metric is equal the to the determinant of the background metric which is assumed to be nondegenerate. Therefore as soon as we want to keep the local diffeomorphism symmetry (3) (i.e. the picture of embedding of one manifold into another) it is assumed everywhere that
$$det\left(G_{\mu \nu }(X)_\alpha X^\mu _\beta X^\nu \right)0.$$
(5)
Throughout the paper we use the following notation: $`\mu `$, $`\nu `$ denote space-time indices, $`\alpha `$, $`\beta `$ world-volume indices and $`a`$, $`b`$, $`c`$ spatial world-volume indices.
Let us assume that the space-time manifold $``$ is compact and oriented. When we come to the Hamiltonian treatment we also assume that $`=R\times \mathrm{\Sigma }`$ where $`\mathrm{\Sigma }`$ is a compact and oriented spatial manifold. $`p`$-branes can be closed or open. For closed $`p`$-branes periodicity conditions must be imposed along all spatial directions. The analysis of open $`p`$-branes is more involved since the theory should be supplemented by appropriate boundary conditions. In this note we look only at closed $`p`$-branes. In the case of closed branes the Nambu-Goto action is a constant for all field configurations and it is equal to the volume of the background manifold $``$. We hope to come back to the case of open $`p`$-branes elsewhere.
In this paper we study mainly the classical aspects of the theory. The paper is organized as follows: In section 2 we go through the Hamiltonian treatment of the closed brane theory. Three equivalent sets of constraints are presented. In section 3 we take a look at the construction of BRST generators for these three sets of constraints. Three different BRST generators are related to each other through canonical transformations in the extended phase space. In section 4 we look at equivalent forms of the action and specifically discuss the case of locally flat backgrounds. The degrees of freedom are briefly considered and the subtleties related to the degenerate solutions are pointed out. In the last section we discuss the results and outline possible generalizations of the model.
## 2 Hamiltonian treatment
In this section we take a look at the Hamiltonian treatment of the system. The $`p`$-brane theory is a generally covariant system and therefore the naive Hamiltonian vanishes identically. In this theory the full Hamiltonian is given by a linear combination of the corresponding constraints which are first class. Our goal is to write down three different sets of constraints for the model .
In order to carry out the Hamiltonian formulation of the theory we choose one of the integration variables $`\xi ^\alpha `$ as the evolution parameter (in the case of a relativistic metric with signature $`(1,1,\mathrm{},1)`$, the one-parameter group of diffeomorphisms defined by the translations in that variable should be generated by a timelike vector field), which we take to be $`\xi ^0`$; the remaining integration variables, which parametrize the brane itself, are represented by $`\xi ^a`$, with the small Latin letters taking values from $`1`$ to $`p`$. The system is totally constrained since the theory is invariant under redefinitions of the evolution parameter.
Denoting by $`P_\mu `$ the momenta conjugate to the $`X^\mu `$ and starting from either the Polyakov action (2) or the Nambu-Goto action (1) the constraints can be worked out as
$$_I=\left(\begin{array}{c}\\ _a\end{array}\right)=\left(\begin{array}{c}G^{\mu \nu }(X)P_\mu P_\nu +T^2det[q_{ab}]\\ P_\mu _aX^\mu \end{array}\right),$$
(6)
where
$$q_{ab}=G_{\mu \nu }(X)_aX^\mu _bX^\nu $$
(7)
is the induced spatial metric on the brane. The constraints (6) are first class and obey the algebra
$`\{_a[M^a],_b[N^b]\}`$ $`=`$ $`_a[\mathrm{\pounds }_MN^a],`$ (8)
$`\{_a[M^a],[N]\}`$ $`=`$ $`[\mathrm{\pounds }_MN],`$ (9)
$`\{[M],[N]\}`$ $`=`$ $`_a[qq^{ab}(M_bNN_bM)],`$ (10)
where $`\mathrm{\pounds }_N`$ stands for the Lie derivative along the vector field $`N^a`$, and $`q=det\left(q_{ab}\right)`$. Since there are $`d`$ pairs of canonical conjugate variables and $`p+1`$ constraints, the theory possesses $`(dp1)`$ degrees of freedom per brane point. Therefore in the case of a $`(d1)`$-brane one has got no dynamical degrees of freedom and the theory is purely topological. The algebra (8)-(10) is called the algebra of many-fingered time (the name is due to Wheeler). The constraints (6) and their algebra (8)-(10) are true for a $`p`$-brane in any space-time dimension $`d`$. The algebra (8)-(10) is closed only for the case $`p<2`$. Now let us analyze the specific properties for a $`p`$-brane propagating on a $`(p+1)`$-dimensional space-time.
Starting from the Nambu-Goto action (1) one can see that the constraints can be written in a form in which all of them are linear in the momenta. In order to do so, we observe that for $`p`$-branes the dimension of the world-volume in equation (1) is the same as the dimension of the embedding space-time, namely $`p+1`$. Consequently $`_\alpha X^\mu `$ is a square matrix, and one can write
$$det\left(G_{\mu \nu }(X)_\alpha X^\mu _\beta X^\nu \right)=G\left(X\right)det{}_{}{}^{2}\left(_\alpha X^\mu \right),$$
(11)
with
$$G\left(X\right)=det\left(G_{\mu \nu }(X)\right).$$
(12)
The action (1) becomes then
$$S=\pm Td^{p+1}\xi \sqrt{𝒢}\frac{1}{(p+1)!}ϵ^{\alpha _0\mathrm{}\alpha _p}ϵ_{\mu _0\mathrm{}\mu _p}_{\alpha _0}X^{\mu _0}\mathrm{}_{\alpha _p}X^{\mu _p},$$
(13)
where $`\pm `$ corresponds to the two possible solutions of the square root. Let us keep both signs in all calculations and eventually one can see that the sign ambiguity corresponds to the two possible orientations on the manifold. The equations of motion for (13) are somewhat trivial. They tell us that the exterior derivative of the volume form is zero. The $`(p+1)`$ decomposition of (1) is straightforward,
$$S=\pm T𝑑\xi ^0d^p\xi n_\mu \dot{X}^\mu ,$$
(14)
where a dotted quantity represents its derivative with respect to the evolution parameter $`\xi ^0`$, and the vector $`n_\mu `$ is a function of the configuration variables $`X^\mu `$ given by
$$n_\mu =\frac{1}{p!}\sqrt{G}ϵ_{\mu \nu _1\mathrm{}\nu _p}ϵ^{a_1\mathrm{}a_p}_{a_1}X^{\nu _1}\mathrm{}_{a_p}X^{\nu _p}.$$
(15)
The vector $`n_\mu `$ satisfies the following properties
$`G^{\mu \nu }n_\mu n_\nu `$ $`=`$ $`det[q_{ab}],`$ (16)
$`n_\mu _aX^\mu `$ $`=`$ $`0.`$ (17)
In this form it is clear that the momenta conjugate to the $`X^\mu `$ are
$$P_\mu =\pm Tn_\mu .$$
(18)
The Hamiltonian vanishes and one must include the primary constraints given by equation (18) with the aid of some Lagrange multiplier functions $`\lambda ^\mu `$,
$`L^\pm `$ $`=`$ $`{\displaystyle d^p\xi \left[P_\mu \dot{X}^\mu \lambda ^\mu 𝒞_\mu ^\pm \right]},`$ (19)
$`𝒞_\mu ^\pm `$ $`=`$ $`P_\mu Tn_\mu (X).`$ (20)
The obtained constraints are linear in the momenta and its Poisson bracket algebra vanishes strongly. The two sets of constraints $`𝒞_\mu ^+`$ and $`𝒞_\mu ^{}`$ correspond to two different branches of the constraint surface. These two sets intersect only on the degenerate solutions
$$𝒞_\mu ^+=𝒞_\mu ^{}=0n_\mu =0det\left(q_{ab}\right)=0.$$
(21)
Thus exluding degenerate solutions one has two independent branches of the theory with different Lagragians $`L^{}`$. The constraints $`𝒞_\mu ^\pm `$ generate the following transformation
$$\delta X^\mu =\{X^\mu ,𝒞_\mu ^\pm [N^I]\}=N^\mu .$$
(22)
In fact one can see that this is the real symmetry of the action (13). The natural question might arise about the relation between (22) and the local diffeomorphism invariance (3). There is a one to one map between the two transformations
$$\delta X^\mu =\mathrm{\pounds }_\zeta X^\mu =(_\alpha X^\mu )\zeta ^\alpha =N^\mu ,$$
(23)
when the quadratic matrix $`(_\alpha X^\mu )`$ is assumed to be nondegenerate. Thus for every vector $`\zeta `$ there is unique vector $`N`$ and vice versa. However it should be stressed that in general the gauge symmetries (3) and (22) have different properties. By using transformation (22) one can bring a nondegenerate solution to a degenerate one. One cannot do this by using the transformation (3).
To get another set of constraints one can contract $`𝒞_\mu ^\pm `$ with the set of independent vectors $`G^{\mu \nu }n_\nu `$ and $`_aX^\mu `$, resulting respectively in
$$_I^\pm =\left(\begin{array}{c}^\pm \\ _a\end{array}\right)=\left(\begin{array}{c}G^{\mu \nu }n_\mu P_\nu \pm Tdet[q_{ab}]\\ P_\mu _aX^\mu \end{array}\right).$$
(24)
Comparing with equations (6) we see that only the scalar constraint is modified, now being linear in the momenta, not quadratic. The constraints (24) obey the same Poisson bracket algebra as (8)-(10). Again the two sets of constraints $`_I^+`$ and $`_I^{}`$ describe the two separate branches if the degenerate solutions are excluded, and they obey two many-fingered algebras on the corresponding independent branches of the theory. The constraints (24) basically tell us that the system can be thought as a parametrized field theory (parametrized cosmological constant term) . The equation of motion for $`X^\mu `$ is given by
$$\dot{X}^\mu =\{X^\mu ,_a[N^a]+^\pm [M]\}=N^a_aX^\mu +MG^{\mu \nu }n_\nu ,$$
(25)
which is nothing else but the geometrodynamical canonical decomposition with respect to basic vectors ($`G^{\mu \nu }n_\nu `$, $`_aX^\mu `$).
Now one can check explicitly the equivalence of these three sets of constraints, $`𝒞_\mu ^\pm `$, $`_I^\pm `$ and $`_I`$. It is clear that $`𝒞_\mu ^\pm `$ implies both $`_I^\pm `$ and $`_I`$. To check the converse we note that the second equations $`_a=0`$ in the sets of constraints, (6) and (24), have the general solution
$$P_\mu =\alpha n_\mu ,$$
(26)
where $`\alpha `$ is any function which does not carry indices. Plugging this result into the last equation of $``$ one gets
$`(T^2\alpha ^2)det\left(q_{ab}\right)=0`$ (27)
$``$ $`T=\alpha \mathrm{or}\mathrm{T}=\alpha \mathrm{or}det[\mathrm{q}_{\mathrm{ab}}]=0.`$
The first solution is just $`𝒞_\mu ^+`$, the second is $`𝒞_\mu ^{}`$ and the last one corresponds to a degenerate solution. Applying the same procedure to $`^\pm `$ one finds that
$$𝒞_\mu ^+=0_I^+=0,𝒞_\mu ^{}=0_I^{}=0,$$
(28)
if $`det\left(q_{ab}\right)0`$ is assumed.
Thus we have shown that all three sets of constraints are equivalent
$$𝒞_\mu ^+=0\mathrm{and}𝒞_\mu ^{}=0_\mathrm{I}=0_\mathrm{I}^+=0\mathrm{and}_\mathrm{I}^{}=0,$$
(29)
if the case of degenerate metric is excluded and therefore these three sets of constraints describe the same constrained surface. Since the manifold $`\mathrm{\Sigma }`$ is assumed oriented the two branches of the theory are dynamically independent.
## 3 BRST generators
In this section we construct the BRST generators which correspond to the different sets of constraints discussed in the previous section.
We have a first class constrained system and its Hamiltonian is a linear combination of the constraints $`\mathrm{\Psi }_I`$. Introducing the ghost variables $`\eta ^I`$ and the ghost momenta $`𝒫_I`$ one can define the classical Grassmann-odd BRST generator (charge) $`Q`$ in the extended phase space
$`Q`$ $`=`$ $`{\displaystyle d^p\xi \eta ^I(\xi )\mathrm{\Psi }_I(\xi )}+{\displaystyle \underset{n=0}{\overset{r}{}}}{\displaystyle d^p\xi _1\mathrm{}d^p\xi _n}`$ (30)
$`Q^{I_1\mathrm{}I_n}(\xi _1,\mathrm{},\xi _n)𝒫_{I_1}(\xi _1)\mathrm{}𝒫_{I_n}(\xi _n),`$
such that $`Q`$ is nilpotent and real. The ghost number of $`\mathrm{\Omega }`$ should be equal to $`1`$. The BRST construction is important because it reveals that the different representations of the constraint surface can be thought of as being obtained from each other by a canonical transformation in the extended phase space. In expression (30) the number $`r`$ is called the rank of the BRST generator. The concept of rank is not intrinsic and can be made equal to zero by appropriate redefinitions of constraints .
For the general case of $`p`$-branes with constraints $`_I`$ given by (6) the classical BRST generator $`Q`$ has been constructed by Henneaux . The rank of $`Q`$ is equal $`p`$. Now we can construct the BRST generator for the other two sets of constraints $`𝒞_\mu ^\pm `$ and $`_I^\pm `$. Since the Poisson bracket algebra for $`𝒞_\mu ^\pm `$ vanishes strongly (thus the algebra is commutative) the BRST operator $`Q^\pm `$ has rank 0
$$Q^\pm =d^p\xi \eta ^\mu (\xi )𝒞_\mu ^\pm (\xi ),$$
(31)
where $`Q^+`$ and $`Q^{}`$ are defined for the two different sectors. In this case the BRST transformation ($`\delta _\pm A=\{A,Q_\pm \}`$) in the extended phase space has a simple form
$`\delta _\pm X^\mu `$ $`=`$ $`\eta ^\mu ,\delta _\pm \eta ^\mu =0,\delta _\pm 𝒫_\mu =\mathrm{\Phi }_\mu ,`$
$`\delta _\pm P_\mu `$ $`=`$ $`T{\displaystyle \frac{1}{(p1)!}}\sqrt{G}ϵ_{\mu \nu _1\nu _2\mathrm{}\nu _p}ϵ^{a_1\mathrm{}a_p}_{a_1}\eta ^{\nu _1}_{a_2}X^{\nu _2}\mathrm{}_{a_p}X^{\nu _p}.`$ (32)
The BRST generator $`𝒬^\pm `$ for the constraints $`_I^\pm `$ can also be worked out, being given by
$`𝒬^\pm `$ $`=`$ $`{\displaystyle }d^p\xi [\eta H^\pm +\eta ^aH_a+(\eta ^a_a\eta +\eta _a\eta ^a)𝒫+`$ (33)
$`+(q_{}q^{ab}\eta _a\eta +\eta ^a_a\eta ^b)𝒫_b],`$
where $`(\eta ,𝒫)`$ and $`(\eta ^a,𝒫_b)`$ are the ghost pairs associated with $`^\pm `$ and $`_a`$ respectivelly. Its rank is 1. Thus we have constructed three different BRST generators for the same theory. They should relate to each other by canonical transformations in the extended phase space. We were unable to construct these canonical transformations in any simple closed form. However one can certainly construct them perturbatively in same fashion as in and .
As we saw the three sets of constraints $`_I`$, $`_I^\pm `$ and $`𝒞_\mu ^\pm `$ describe the same constraint surface if we exclude degenerate solutions. There should be the following relation among the sets of constraints which describe the same constrain surface
$$_I=(S^\pm )_I^\mu 𝒞_\mu ^\pm ,_I^\pm =(S)_I^\mu 𝒞_\mu ^\pm ,$$
(34)
where $`(S^\pm )_I^\mu `$ and $`(S)_I^\mu `$ must be non degenerate. It is not difficult to construct these matrices explicitly. Thus for $`S^\pm `$ we have the following expression
$$S^\pm =\left(\begin{array}{ccc}_1X^0\hfill & \mathrm{}\hfill & _1X^p\hfill \\ _2X^0\hfill & \mathrm{}\hfill & _2X^p\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ G^{0\nu }(P_\nu \pm Tn_\nu )\hfill & \mathrm{}\hfill & G^{p\nu }(P_\nu \pm Tn_\nu )\hfill \end{array}\right),$$
(35)
and for $`S`$ the following
$$S=\left(\begin{array}{cccc}_1X^0\hfill & _1X^1\hfill & \mathrm{}\hfill & _1X^p\hfill \\ _2X^0\hfill & _2X^1\hfill & \mathrm{}\hfill & _2X^p\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ G^{0\nu }n_\nu \hfill & G^{1\nu }n_\nu \hfill & \mathrm{}\hfill & G^{p\nu }n_\nu \hfill \end{array}\right).$$
(36)
Let us calculate the determinants of these matrixes
$`det(S^\pm )`$ $`=`$ $`(1)^pn_\mu G^{\mu \nu }(P_\nu \pm Tn_\nu )=(1)^p\stackrel{~}{}^\pm ,`$ (37)
$`det(S)`$ $`=`$ $`(1)^{p+1}det[q_{ab}].`$ (38)
We see that $`S`$ is non degenerate if degenerate solutions ($`det[q_{ab}]=0`$) are excluded. The matrix $`S^+`$ is not degenerate either as long as we stay at the branch defined by $`𝒞_\mu ^+=0`$ (or equivalently by $`_I^+`$) and $`S^{}`$ is not degenerate at the branch defined by $`𝒞_\mu ^{}`$ (or equivalently by $`_I^{}`$). Therefore using these matrices one can construct perturbatively the relevant canonical transformations in the extended phase space.
## 4 $`p`$-brane theory in locally flat background
To understand the theory better we would like to study alternative representations of this model. In many cases alternative representations of a theory may help to analyze their degrees of freedom. In this section we study some classically equivalent actions and analyze the degrees of freedom corresponding to the topological $`p`$-brane theory. It is hard to say anything explicit about the degrees of freedom when the theory is formulated in the form (1) or (2). Intuitively we understand that the number of degrees of freedom is related to the number patches needed to cover the manifold $``$. However it is hard to count them explicitly. Therefore we can try to reformulate the theory in a more transparent way. One can reach this goal by using new variables. Since the task is difficult for generic curved background manifolds, we look first at the case of locally flat manifolds $``$
$$G_{\mu \nu }=\eta _{\mu \nu }.$$
(39)
At the end of this section we will take a brief look on equivalent actions for the generic case. However it is still problematic to analyze the degrees of freedom in all generality.
Now we are assuming that (39) holds. Let us enlarge the gauge symmetry of the system defining the tetrad fields
$$e_a{}_{}{}^{\mu }=_aX^\mu $$
(40)
as the new configuration variables. They are subject to the constraints
$$_{[a}e_{b]}{}_{}{}^{\mu }=0.$$
(41)
One can easily see that there is a one to one correspondence between new and old variables in the locally flat space-time. The static gauge (4) in new variables corresponds to $`e_a{}_{}{}^{\mu }=\delta _a^\mu `$.
The action in these new variables and their canonical conjugate momenta $`\pi ^a_\mu `$ can be obtained from the generating functional depending on the old coordinates and the new momenta
$$S_{X\pi }=d^p\xi _aX^\mu \pi ^a{}_{\mu }{}^{}.$$
(42)
One has
$`e_a^\mu `$ $`=`$ $`{\displaystyle \frac{\delta S_{X\pi }}{\delta \pi ^a_\mu }}=_aX^\mu ,`$ (43)
$`P_\mu `$ $`=`$ $`{\displaystyle \frac{\delta S_{X\pi }}{\delta X^\mu }}=_a\pi ^a{}_{\mu }{}^{}.`$ (44)
Plugging this result into equation (20) one gets
$`S`$ $`=`$ $`{\displaystyle }d^{p+1}\xi \pi ^a{}_{\mu }{}^{}\dot{e}_{a}^{}{}_{}{}^{\mu }+\varphi ^{ab}{}_{\mu }{}^{}_{[a}^{}e_{b]}{}_{}{}^{\mu }+`$
$`+\lambda ^\mu \{_a\pi ^a{}_{\mu }{}^{}\pm T{\displaystyle \frac{1}{p!}}ϵ^{a_1\mathrm{}a_d}ϵ_{\mu \nu _1\mathrm{}\nu _p}e_{a_1}{}_{}{}^{\nu _1}\mathrm{}e_{a_p}{}_{}{}^{\nu _p}\},`$
where $`\varphi ^{ab}_\mu `$ are the Lagrange multiplier functions for the constraints (41). In the case of a nonflat metric the Lagrangian (LABEL:Lagr) would be nonlocal in the new variables since it involves the original coordinates $`X^\mu `$ present in the determinant of the metric $`G_{\mu \nu }`$. But this problem does not arise in the case of a flat metric. We have then the following action
$`S`$ $`=`$ $`{\displaystyle }d^{p+1}\xi [\pi ^a{}_{\mu }{}^{}_{0}^{}e_a{}_{}{}^{\mu }+\varphi ^{ab}{}_{\mu }{}^{}_{[a}^{}e_{b]}{}_{}{}^{\mu }+\lambda ^\mu _a\pi ^a{}_{\mu }{}^{}\pm `$ (46)
$`\pm T\lambda ^\mu {\displaystyle \frac{1}{p!}}ϵ^{a_1\mathrm{}a_p}ϵ_{\mu \nu _1\mathrm{}\nu _p}e_{a_1}{}_{}{}^{\nu _1}\mathrm{}e_{a_p}{}_{}{}^{\nu _p}],`$
which can be given in a covariant form if one identifies the Lagrange multipliers $`\lambda ^\mu `$ with the time components of the tetrad fields,
$$\lambda ^\mu =e_0{}_{}{}^{\mu },$$
(47)
and writes the momenta $`\pi ^a_\mu `$ and the Lagrange multipliers $`\varphi ^{ab}_\mu `$ as the components of a $`(p1)`$-form $`F_\mu `$,
$`\pi ^a_\mu `$ $`=`$ $`ϵ^{ab_1\mathrm{}b_{p1}}F_{b_1\mathrm{}b_{p1}\mu },`$ (48)
$`\varphi ^{ab}_\mu `$ $`=`$ $`(p1)ϵ^{abc_1\mathrm{}c_{p2}}F_{0c_1\mathrm{}c_{p2}\mu }.`$ (49)
Equation (46) then becomes
$`S`$ $`=`$ $`{\displaystyle }d^{p+1}\xi [ϵ^{\alpha _0\mathrm{}\alpha _p}_{\alpha _0}e_{\alpha _1}{}_{}{}^{\mu }F_{\alpha _2\mathrm{}\alpha _p\mu }^{}\pm `$ (50)
$`\pm T{\displaystyle \frac{1}{(p+1)!}}ϵ^{\alpha _0\mathrm{}\alpha _p}ϵ_{\nu _0\mathrm{}\nu _p}e_{\alpha _0}{}_{}{}^{\mu _0}\mathrm{}e_{\alpha _p}{}_{}{}^{\mu _p}],`$
which can be compactly written in the differential form language as
$$S=F_\mu de^\mu \pm T\frac{1}{(p+1)!}ϵ_{\nu _0\mathrm{}\nu _p}e^{\nu _0}\mathrm{}e^{\nu _p},$$
(51)
where $`e^\nu `$ is a one-form and $`F_\mu `$ is a $`(p1)`$-form. The action (51) is explicitly topological since it does not involve the metric. After all one can see just at level of actions that the actions (13), (51) are equivalent to each other. This equivalence can be established by integrating out the field $`F_\mu `$.
Now let us take a look at the symmetries and equiations of motions of the action (51). The action has the following obvious symmetry
$$\delta F_\mu =dw_\mu ,$$
(52)
which is the shift $`F_\mu `$ by any exact $`(p1)`$-form. There is one extra symmetry which is less obvious
$`\delta e^\mu `$ $`=`$ $`df^\mu ,`$ (53)
$`\delta F_\mu `$ $`=`$ $`\pm {\displaystyle \frac{T}{(p1)!}}ϵ_{\mu \nu _1\nu _2\mathrm{}\nu _p}f^{\nu _1}e^{\nu _2}\mathrm{}e^{\nu _p},`$ (54)
where $`f^\mu `$ is an arbitrary zero-form (function). The equations of motion are the following
$`de^\mu `$ $`=`$ $`0,`$ (55)
$`dF_\mu `$ $`=`$ $`{\displaystyle \frac{T}{p!}}ϵ_{\mu \nu _1\mathrm{}\nu _p}e^{\nu _1}\mathrm{}e^{\nu _p}.`$ (56)
The classical moduli space is given by gauge non-equivalent solutions of equations (55). Thus we were able to reformulate the topological $`p`$-brane theory in a locally flat background as an abelian BF-like model with the action given by (51). The model has a bunch of U(1) fields $`e^\mu `$ and the nontriviality comes from the last ”mass” term which mixes different gauge fields. For the case $`p2`$ the action (51) can be thought as the zero gravitational constant limit for the general relativity with cosmological constant in $`(p+1)`$ dimensional space-time. This limit should be taken in the first order formalism .
The degrees of freedom (the classical moduli space) for the action (51) can be analyzed in a straightforward fashion through the cohomology groups. In general the situation depends on the details of the topology of the background manifold or more precisely, on the structure of the first cohomology group $`H^1(,R)`$. Since the one-forms $`e^\mu `$ are closed and any two solutions that differ by an exact one-form are gauge equivalent, $`e^\mu `$ is a element of $`H^1(,R)`$. The equations for $`F_\mu `$ are more difficult to analyze since the right hand side involves $`e^\mu `$. If $`dim(H^1(,R))<p`$ then the last equation of motion in (55) reduces to $`dF_\mu =0`$. We have not enough elements of the first cohomology group to construct a non-zero right-hand side. Thus in this case the model coincides with $`(p+1)`$ copies of an abelian BF system . Therefore the space of solutions for $`e^\mu `$ and $`F_\mu `$ is given by $`p+1`$ copies of $`H^1(,R)H^{p1}(,R)`$. In the case $`=R\times \mathrm{\Sigma }`$ we have
$$H^1(,R)H^1(\mathrm{\Sigma },R)H^{p1}(\mathrm{\Sigma },R)H^{p1}(,R)$$
(57)
where we used Poincaré duality on the $`p`$-dimensional manifold $`\mathrm{\Sigma }`$. Thus the space of gauge inequivalent solutions is even dimensional and it is given by the product of $`2(p+1)`$ copies of the first cohomology group: $`H^1(\mathrm{\Sigma },R)`$. The situation with $`dim(H^1(,R))p`$ is more involved. One should analyze what kind of right hand side in the last equation (55) can be constructed from $`e^\mu `$. For instance in the case $`=R\times \mathrm{\Sigma }`$ it might be possible to construct out of $`e^\mu `$ the volume form for $`\mathrm{\Sigma }`$: $`e^1\mathrm{}e^p`$. Since the volume form cannot be exact the corresponding equation has no solution. We will not analyze this situation in all generality. However the task might be solved straightforwardly as soon as we know explicitly the content of $`H^1(,R)`$. Above analysis of degrees of freedom is appropriate for the actions (13) and (51) where the degenerate solutions are included. However to incorporate into the analysis the restriction of nondegeneracy can be hard since the removal of degennerate solutions from the phase space might destroy the gauge orbits. The similar problem appears in the relation between $`2+1`$ gravity and Chern-Simons theory .
On a curved space-time manifold there is no such simple BF-like action as in locally flat case. However one can write the following action
$$S=(dX^\mu \eta ^\mu )B_\mu \pm T\frac{1}{(p+1)!}\sqrt{G}ϵ_{\mu _0\mathrm{}\mu _p}\eta ^{\mu _0}\mathrm{}\eta ^{\mu _p},$$
(58)
which is classically equivalent to the Nambu-Goto action (13). In the action (58) $`\eta ^\mu `$ and $`B_\mu `$ are $`1`$-forms and $`p`$ -forms respectively. The action is nonlinear in $`X^\mu `$ and therefore it is difficult to analyze it in the same fashion as before. The case $`p=1`$ is definitly special. By itself the Nambu-Goto action (13) can be interpreted as topological sigma model in two dimensions since $`\sqrt{G}ϵ_{\mu \nu }`$ might serve as closed symplectic form on $``$. Also the Nambu-Goto action is equivalent to the following action
$$S=𝑑X^\mu \eta _\mu \pm \frac{T}{2}(𝒢)^{1/2}ϵ^{\mu \nu }\eta _\mu \eta _\nu $$
(59)
which is the Poisson sigma model on two dimensional $``$ . Therefore we see that two dimensional topological string theory is classically equivalent to other known theories up to some subtleties related to degenerate configurations.
## 5 Discussion and outline
In the present work we considered the classical aspects of closed $`p`$-brane theory defined on $`(p+1)`$ dimensional background manifolds $``$. We analyzed the hamiltonian and BRST structure of the theory. We saw that model has different equivalent realizations. However the classical equivalence between the constraints and the actions might fail at the quantum level due to normal ordering problem (different regularizations). One can look at the most familiar example $`p=1`$. For the case of quadratic constraints there is an anomaly in the Virasoro algebra and therefore the system is not first class anymore. In the case of linear constraints (20) there is no anomaly possible since the constraints are completely linear. This discussion gives us an example that at the quantum level the Nambu-Goto action (natural source for linear constraints) and the Polyakov action (the natural source for quadratic constraints) are not equivalent to each other. As well at the classical level different status of degenerate solutions can bring extra problems into identification of two theories.
The actions (13) and (58) have a straightforward generalization to the following topological models
$$S=T\frac{1}{(p+1)!}d^{p+1}\xi C_{\mu _0\mathrm{}\mu _p}(X)ϵ^{\alpha _0\mathrm{}\alpha _p}_{\alpha _0}X^{\mu _0}\mathrm{}_{\alpha _p}X^{\mu _p}$$
(60)
and
$$S=(dX^\mu \eta ^\mu )B_\mu +T\frac{1}{(p+1)!}C_{\mu _0\mathrm{}\mu _p}(X)\eta ^{\mu _0}\mathrm{}\eta ^{\mu _p},$$
(61)
where $`C`$ is a $`(p+1)`$-form defined on the $`d`$-dimensional background manifold $``$ ($`d`$ might be any value equal or greater than $`(p+1)`$). If the form $`C`$ is closed the model has many similarities with the topological $`p`$-brane studied in the present work. We shall consider the classical and quantum aspects of these theories in coming work.
Acknowledgments
We are grateful to Ansar Fayyazuddin for discussions.
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# Magnetization of ferrofluids with dipolar interactions: A Born–Mayer expansion
## I Introduction
Ferrofluids are suspensions of ferromagnetic particles of about 10 nm diameter in a carrier fluid. The particles are stabilized against aggregation by coating with polymers or by electrostatic repulsion of charges brought on their surface. On macroscopic scales, ferrofluids can be described as liquids with intrinsic superparamagnetic properties.
In this paper, we are concerned with the equilibrium magnetization $`M`$ as a function of the internal magnetic field $`H`$ for given temperature and particle concentration. Sufficiently low concentrated ferrofluids behave like a paramagnetic gas. Therein the interaction between the particles can be neglected and the equilibrium magnetization can be described properly by the Langevin function. The magnetic properties are then necessarily weak. To produce ferrofluids with strong magnetic properties one has to have either a higher particle concentration or one has to use ferromagnetic material with a large bulk magnetization, e. g., cobalt instead of magnetite. In both cases the magnetization is strongly influenced by dipole–dipole and other interactions between the particles.
Several models of dipolar interacting systems have been studied in the literature. Numerical investigations were based on density functional approaches and Monte Carlo simulations . The models differ in the treatment of the short range interactions, which were described by hard sphere , other hard core potentials , soft sphere , or Lenard–Jones potentials . These investigations were mainly undertaken to reveal the phase transition properties. These properties are substantially different for different short range interactions. Thus, for example the question whether a system of particles interacting via long range dipolar forces shows without any dispersive energy, e. g., from attractive van der Waals energy a ”liquid–vapor” phase coexistence of a dense and a less dense phase is currently being discussed .
Analytical models focus mainly on the equilibrium magnetization in the gas phase (were the term ”gas” refers, as far as ferrofluids are addressed, to the magnetic particle subsystem within the liquid carrier). Such models are the Onsager model , the Weiss model , the mean spherical approximation , and an approach by Buyevich and Ivanov (called high temperature approximation in ). These models were tested experimentally for ferrofluids . Especially the mean spherical model and the high temperature approximation showed good results .
Our approach assumes the magnetic particles in the ferrofluid to be hard spheres with a common diameter $`D`$ and dipolar moment $`m`$. We use the technique of the Born–Mayer expansion together with an expansion in the strength of the dipolar coupling to get analytical approximations. They are obtained via series expansions of the free energy in terms of two parameters: (i) the volume fraction of the hard core particles $`\varphi `$ and (ii) a dimensionless dipolar coupling constant $`ϵ`$, given by the ratio between a typical dipolar energy for particles in hard core contact and the thermal energy $`kT`$. Our result for the magnetization goes beyond the high temperature approximation and reduces to it in linear order in $`\varphi `$ and $`ϵ`$.
Dipolar forces fall off as $`r^3`$ and are thus of long range nature. This long range character requires great care when invoking the thermodynamic limit . To circumvent the problem we model the dipolar fields that are generated by distant particles by a magnetic continuum field (similar to the treatment in the Weiss model) while incorporating the near field contributions explicitly in a statistical mechanical description. The magnetization $`M`$ is then derived as a function of the internal magnetic field $`H`$. The so obtained relation $`M(H)`$ is independent of the probe geometry. Once $`M(H)`$ is known, the magnetization for a given geometry can in principle be derived by solving the macroscopic Maxwell equations. This may practically still be a difficult task, at least as long as the external field is small or absent. In this case it is known that the magnetization will show for general shaped probes a nontrivial spatial variation at high enough densities .
Since our method yields an expression for the free energy of the model system we can in principle calculate also other thermodynamic quantities and in particular address the questions of phase transition, e. g., between gas and liquid or between ferromagnetic and non–ferromagnetic phases. We have not addressed the question of a gas–liquid transition of the magnetic particles suspended in the ferrofluid since it is believed that short range van der Waals–like attractions would have to be incorporated to model real ferrofluids appropriately in this regard . However, the question whether a strong dipolar coupling induces in zero external field a spontaneous magnetization that is currently debated in the literature is briefly touched upon in Sec. VI A of this paper.
The paper is organized as follows: In Sec. II we discuss the connections between the various fields that are of relevance in a ferrofluid. We present the model to treat the long range dipolar forces. In Sec. III we present the expansion method to get analytical solutions in terms of the two small parameters $`ϵ`$ and $`\varphi `$. In Sec. IV we calculate an expression for the magnetization that contains only linear terms in $`\varphi `$ but, at least in principle, arbitrary high orders in $`ϵ`$. In Sec. V a different expression is derived containing also quadratic terms in $`\varphi `$ but also only up to second order terms in $`ϵ`$. In Sec. VI we discuss our findings and investigate the applicability of the results in the $`\varphi `$$`ϵ`$ plane. Sec. VII contains a short conclusion.
## II Magnetic fields and magnetization
We are interested in the effect of dipolar interactions of the magnetic particles in a ferrofluid on the equilibrium magnetization of the ferrofluid. To that end we consider the ferrofluid as an ensemble of identical spherical particles of diameter $`D`$, each carrying a magnetic moment of magnitude $`m`$. These particles interact with each other via magnetic dipole–dipole interactions and a hard core repulsion with hard core diameter $`D`$. We assume $`DD_{mag}`$ where $`D_{mag}`$ is the diameter of the magnetic core of the particles thus allowing for a surfactant surface layer that provides a steric repulsion.
The particles can lower their potential energy by orienting their magnetic moments parallel to a local magnetic field. However any interaction of the particles with the fluid medium they are suspended in is ignored. The latter is taken to be magnetically inert.
### A Different magnetic fields
Before we outline in Sec. II B how we determine in principle the magnetization of the ferrofluid we should like to review briefly the different magnetic fields that one has to distinguish and that are of importance in a system with dipolar interactions. The first field is the external magnetic field $`𝐇_e`$ that is applied outside of the probe. If dipolar interactions can be neglected, $`𝐇_e`$ is also the field at the position of the particles – at least as long as the carrier fluid can be treated a magnetic vacuum which we will assume throughout the paper. In the presence of dipolar interactions additional fields have to be considered. One of them is the internal field $`𝐇`$, that is the macroscopic field inside the probe. By assuming that the equation of state $`𝐌(𝐇)`$ is known, $`𝐇`$ can be calculated by the common methods of continuum magnetostatics. But the macroscopic field $`𝐇`$ differs in general from the field $`𝐇_{local}`$ that the magnetic particles feel.
So far one has employed in the ferrofluid literature two models to calculate $`𝐇_{local}`$ from $`𝐇`$ that are similar in spirit, namely the Weiss model and the Onsager model . Both introduce a virtual hollow sphere inside a magnetic continuum such that the sphere contains a single magnetic particle in its center. The Weiss model assumes the magnetization $`𝐌`$ and the internal field $`𝐇`$ to be constant everywhere in the magnetic continuum surrounding the sphere. Then the field inside the sphere is given by
$$𝐇_s=𝐇+\frac{𝐌}{3}.$$
(1)
This is the field that the single magnetic particle feels within the Weiss model, i. e. $`𝐇_{local}=𝐇_s`$.
The Onsager model, on the other hand, is restricted to linearly responding fluids and calculates the field inside the sphere on the assumption that it is really hollow and that therefore $`𝐇`$ and $`𝐌`$ differ near the sphere from its bulk values. In that case the field within the sphere is
$$𝐇_{local}=\frac{3\chi +3}{2\chi +3}𝐇.$$
(2)
$`\chi `$ is the susceptibility. In both models the magnetization is calculated as the magnetization of a system of noninteracting dipoles in the magnetic field $`𝐇_{local}`$, i. e.,
$$M=M_{sat}\left(\frac{m}{kT}H_{local}\right).$$
(3)
Here
$$M_{sat}=\frac{N}{V}\frac{m}{\mu _0}$$
(4)
is the saturation magnetization of the fluid, $``$ the Langevin function, $`m`$ the magnetic moment of the particles, and $`N/V`$ their number density. In the Onsager model the Langevin function is consistently used only in linear order. Letting $`M=\chi H`$ on the left hand side of (3) and using (2) allows to calculate $`\chi `$.
In the Weiss model the selfconsistent solution $`M(H)`$ is determined using (3) and (1). The Onsager and Weiss models differ in the treatment of the back reaction of the particle inside the sphere on the magnetic continuum near the sphere’s boundary.
### B Decomposition of fields
Since the magnetic continuum is a macroscopic concept one should be careful when using it on the mesoscopic length scales of the interparticle distances and particle diameters. A first–principle statistical mechanical calculation of the magnetization would start with expressing the energy of the system in terms of the statistical variables of the constituents. In this context the local magnetic field $`𝐇_{local}`$ that a magnetic moment, say, at position $`𝐱_i`$ feels is of importance. It is composed of two different magnetic fields, the external field $`𝐇_e`$ and the dipolar contribution $`𝐇_{dipole}`$ from the other $`N1`$ particles at positions $`𝐱_j`$, possessing a magnetic moment $`𝐦_j`$. Thus within the first–principles approach one has $`𝐇_{local}=𝐇_e+𝐇_{dipole}`$, where
$$𝐇_{dipole}(𝐱_𝐢)=\underset{j}{}\frac{3\widehat{𝐫}_{ij}(𝐦_j\widehat{𝐫}_{ij})𝐦_j}{4\pi \mu _0r_{ij}^3}.$$
(5)
Here $`𝐫_{ij}=𝐱_𝐢𝐱_𝐣`$, $`r_{ij}=|𝐫_{ij}|`$, and $`\widehat{𝐫}_{ij}=𝐫_{ij}/r_{ij}`$.
The long range character of the dipolar forces requires special care when invoking the thermodynamic limit $`V\mathrm{}`$ and $`N\mathrm{}`$ ($`N/V=`$ const.) — for a critical discussion see, e.g., Ref. . The reason is that the dipolar contribution (5) will in general depend on the geometry of the ferrofluid probe and the location of $`𝐱_𝐢`$ within it. Thus the equilibrium magnetization of a probe in an external field will in general depend on the geometry of the latter and furthermore it will be spatially varying. We therefore use here an approach similar to the one that has been used successfully in solid state theory to determine, e. g., the crystal field splitting caused by local fields. It properly accounts for the contributions from microscopic and macroscopic scales.
Consider some magnetic particle $`i`$ in a ferrofluid probe in thermodynamic equilibrium. The particles beyond some distance $`R_s`$ from $`𝐱_i`$ can be considered as independent from particle $`i`$ if $`R_s`$ is larger than the correlation length induced by the dipolar interactions. Furthermore, if $`R_s`$ is large enough, their contribution to the local field at $`𝐱_i`$ can be approximated by a contribution from a magnetic continuum with equilibrium magnetization $`𝐌`$ and macroscopic field $`𝐇`$. We assume that the distance $`R_s`$ is still small compared to the length scale on which the macroscopic fields $`𝐌`$ and $`𝐇`$ vary. Thus we introduce a virtual sphere of radius $`R_s`$ around particle $`i`$ (dark particle in the center of Fig. 1) to separate the dipolar field into a ”far” and a ”near” contribution
$$𝐇_{dipole,far}(𝐱_i)=\underset{r_{ij}>R_s}{}\frac{3\widehat{𝐫}_{ij}(𝐦_j\widehat{𝐫}_{ij})𝐦_j}{4\pi \mu _0r_{ij}^3},$$
(6)
$$𝐇_{dipole,near}(𝐱_i)=\underset{r_{ij}<R_s}{}\frac{3\widehat{𝐫}_{ij}(𝐦_j\widehat{𝐫}_{ij})𝐦_j}{4\pi \mu _0r_{ij}^3}.$$
(7)
Then
$$𝐇_{local}=𝐇_e+𝐇_{dipole,far}+𝐇_{dipole,near}.$$
(8)
If the sphere would be empty, $`𝐇_e+𝐇_{dipole,far}=𝐇_s`$ would be the local field inside the sphere.
A key point of our treatment is to express the far field $`𝐇_{dipole,far}`$ within the continuum approximation. Using this approach the field in the empty sphere is given by $`𝐇_s=𝐇+𝐌/3`$ (1). Note that this approximation is not valid near the sphere’s boundary. But at the center of the sphere $`𝐇_{local}`$ consists of the continuum contribution $`𝐇_s`$ from the far region and the contribution for the dipoles within the sphere:
$$𝐇_{local}=𝐇_s+𝐇_{dipole,near}=𝐇+\frac{𝐌}{3}+𝐇_{dipole,near}.$$
(9)
This result agrees with Eq. (27.26) of Ref. where it has been derived with slightly different arguments for electric dipoles. Due to the long range character of the dipolar forces $`𝐇_{dipole}`$ will be in general geometry dependent and spatially varying. $`𝐇_s`$, respectively $`𝐇`$ and $`𝐌`$ will then also show these features as mentioned above.
If the dipolar coupling between the particles is so weak that even the dipolar fields of the nearest neighbours of particle $`i`$ can be described by a continuum field, i. e., if $`R_s`$ can be chosen as being smaller than the mean distance between the particles, we can drop the contribution $`𝐇_{dipole,near}`$ altogether and arrive at the Weiss model, where a single particle is located inside the hollow sphere in the continuum.
### C Equilibrium magnetization
We want to determine the thermodynamic equilibrium relation $`M(H)`$ between the magnetization
$$𝐌=\frac{1}{\mu _0V}\underset{i}{}𝐦_i=\frac{N}{\mu _0V}𝐦$$
(10)
and the macroscopic magnetic field $`𝐇`$ in the thermodynamic limit. Instead of considering the dipolar interaction of all particles in an external field in the statistical mechanical problem (10) we take explicitly only interactions between those particles into account whose distance is smaller than the sphere radius $`R_s`$. The other interactions are represented by the far–field continuum approximation $`𝐇_s=𝐇+𝐌/3`$. The magnetization $`M_{sphere}(H_s)`$ resulting from this decomposition of fields is then identified with the equilibrium magnetization $`M(H)`$ of the ferrofluid
$$M_{sphere}(H+M/3)=M(H).$$
(11)
Thus after having obtained the approximate expressing for $`M_{sphere}`$ as a function of $`H+M/3`$ we then obtain from solving (11) for $`M`$ an approximation for the sought after equilibrium relation $`M(H)`$. The functional dependence of $`M`$ on $`H`$ is independent of the probe geometry.
In the limit of weak dipolar coupling or when $`R_s`$ becomes smaller than the mean distance between the particles we find $`M_{sphere}=M_{sat}\left[\frac{m}{kT}\left(H+M/3\right)\right]`$ so that the Weiss model is recovered as discussed above.
The magnetization $`M_{sphere}`$ in (11) depends on two dimensionless parameters that characterize the thermodynamic state of the ferrofluid. One of these parameters is the volume concentration of the particles
$$\varphi =\frac{N}{V}\frac{\pi D^3}{6}.$$
(12)
The ratio $`\varphi _{mag}`$ of the volume of the magnetic material to the total volume is $`\varphi _{mag}=(D_{mag}/D)^3\varphi `$. The other parameter is
$$ϵ=\frac{m^2}{4\pi \mu _0kTD^3},$$
(13)
the ratio between a typical energy of dipole–dipole interaction of particles in contact (i. e., at the distance of the hard core diameter $`D`$) and the thermal energy $`kT`$.
## III Canonical partition function
We use the canonical ensemble average to evaluate $`M_{sphere}`$. Given a system of $`N`$ interacting particles with an interaction potential $`V_{ij}`$ ($`1<i,j<N`$) and external potential per particle $`V_i`$, the canonical partition function is given by
$$Z=e^{_kv_k_{i<j}v_{ij}}𝑑\mathrm{\Gamma }.$$
(14)
Here $`v_i=V_i/kT`$, $`v_{ij}=V_{ij}/kT`$, and $`d\mathrm{\Gamma }`$ means integration over the configuration space. In our case, a configuration is characterized by specifying the position vector $`𝐱_i`$ and two angles for each magnetic particle. The two angles define the direction of the magnetic moment $`𝐦_i`$. The modulus $`m`$ is assumed to be constant and the same for all particles. Note that we ignore any translational and rotational degrees of freedom of the particles that carry the magnetic moments, since they have no effect on the magnetization. Only the locations of the moments, i. e., of the particles and the orientations of the moments are considered as statistical variables.
In the first–principles statistical mechanical problem identified by a superscript $`0`$, the external potential would be the energy of a dipole in the external magnetic field
$$V_i^0=𝐦_i𝐇_e.$$
(15)
The interparticle potential is modelled by a dipole–dipole (DD) interaction plus hard core (HC) repulsion. Thus
$$V_{ij}^0=V_{ij}^{0,DD}+V_{ij}^{HC},$$
(16)
$$V_{ij}^{0,DD}=\frac{3(𝐦_i\widehat{𝐫}_{ij})(𝐦_j\widehat{𝐫}_{ij})𝐦_i𝐦_j}{4\pi \mu _0r_{ij}^3},$$
(17)
$$V_{ij}^{HC}=\{\begin{array}{cc}0& \text{for }r_{ij}>D\hfill \\ \mathrm{}& \text{for }r_{ij}<D\hfill \end{array}.$$
(18)
The replacement of the dipolar magnetic fields from far–away particles by a field that has its origin in a magnetic continuum results in a new canonical partition function with the ”external” potential of a dipole
$$V_i=𝐦_i𝐇_s,$$
(19)
in the field $`𝐇_s=𝐇+𝐌/3`$ and a dipolar interaction term with a cutoff, i. e.,
$$V_{ij}^{DD}=\{\begin{array}{cc}V_{ij}^{0,DD}& \text{for }r_{ij}<R_s\hfill \\ 0& \text{for }r_{ij}>R_s\hfill \end{array}.$$
(20)
Note that when using (19) and (20) in the expression (14) for the partition function we describe every particle $`i`$ as being at the center of a sphere of radius $`R_s`$ inside a magnetic continuum such that each particle feels the ”external” field $`𝐇_s`$ and explicit dipolar fields $`𝐇_{dipole,near}(𝐱_i)`$ (7) from the particles whose distance is smaller than $`R_s`$.
The aforementioned statistical mechanical problems with the long range nature of the bare dipolar interactions is thus circumvented by the cutoff at $`R_s`$ in (20) that results from decomposing dipolar fields into a near and a far contribution. Dipolar forces appear explicitly only as forces with a finite range. Their influence on the magnetization in our approach is therefore independent of the geometry of the probe. The geometry dependence enters only via the effective ”external” field $`H_s`$ from the far field contribution.
### A Born–Mayer expansion method
Since an integral such as (14) is hard to solve even numerically we use the Born–Mayer expansion method to get analytical results. The key point of this method is to write
$$Z=\underset{k}{}e^{v_k}\underset{i<j}{}(1+f_{ij})d\mathrm{\Gamma },$$
(21)
where
$$f_{ij}=e^{v_{ij}}1.$$
(22)
If the typical interaction energy is small compared to $`kT`$, the $`f_{ij}`$ can be considered as small parameters, and $`Z`$ can be expanded into a series:
$`Z`$ $`=`$ $`{\displaystyle \underset{m}{}e^{v_m}d\mathrm{\Gamma }}+`$ (26)
$`{\displaystyle \underset{m}{}e^{v_m}\underset{i<j}{}f_{ij}d\mathrm{\Gamma }}+`$
$`{\displaystyle \underset{m}{}e^{v_m}\underset{i<j}{}f_{ij}\underset{k<l}{}f_{kl}d\mathrm{\Gamma }}+`$
$`\mathrm{}.`$
These integrals can be factorized and are therefore easier to handle.
The first order of (26) contains terms like
$$\underset{m}{}e^{v_m}e^{v_{12}^{DD}v_{12}^{HC}}d\stackrel{}{𝐱}d\stackrel{}{\mathrm{\Omega }}.$$
(27)
Here $`d\stackrel{}{𝐱}d\stackrel{}{\mathrm{\Omega }}`$ is an abbreviation for $`d𝐱_1\mathrm{}d𝐱_Nd\mathrm{\Omega }_1..d\mathrm{\Omega }_N`$, and $`d\mathrm{\Omega }_i`$ means the integration over the possible orientations of $`𝐦_i`$.
### B Expansion in powers of $`v^{DD}`$
Obviously even the first order still cannot be evaluated analytically. Therefore a second series expansion is made
$$f_{ij}=e^{v_{ij}^{HC}}e^{v_{ij}^{DD}}1=f_{ij}^{(0)}+f_{ij}^{(1)}+f_{ij}^{(2)}+\mathrm{},$$
(29)
$`f_{ij}^{(0)}`$ $`=`$ $`\left(e^{v_{ij}^{HC}}1\right)`$ (30)
$`f_{ij}^{(\alpha )}`$ $`=`$ $`{\displaystyle \frac{\left(v_{ij}^{DD}\right)^\alpha }{\alpha !}}e^{v_{ij}^{HC}}\mathrm{},\alpha 1.`$ (31)
So we expand $`f_{ij}`$ in powers of the reduced dipolar interaction $`v_{ij}^{DD}`$. The integrals in (26) that remain to be solved are of the form
$$\underset{m}{}e^{v_m}f_{ij}^{(\alpha )}f_{kl}^{(\beta )}\mathrm{}d\stackrel{}{𝐱}d\stackrel{}{\mathrm{\Omega }}.$$
(32)
We now introduce a modification of the common graphical representation of Born–Mayer integrals as follows
1. Every distinct particle that appears via interaction terms of the form $`f_{ij}^{(\alpha )}`$ is represented by a circle.
2. A zeroth order interaction term $`f_{ij}^{(0)}`$ is represented by an overlap of the circles $`i`$ and $`j`$.
3. First, second, … order interaction is represented by one, two, … lines connecting the circles.
Note that the representation of the zeroth order dipolar interaction by two overlapping circles is a reminder that in this case the integrand is nonzero only if the particles are assumed to be in a configuration in which they would indeed overlap.
It turns out, that the expansion in terms of the $`f_{ij}^{(\alpha )}`$ means an expansion in powers of the two parameters $`ϵ`$ and $`\varphi `$, that define the thermodynamical system. Every line in a representing graph, i. e., every power of $`v_{ij}^{DD}`$ results in a factor $`ϵ`$. Every $`n`$–particle subgraph in which all circles are connected to each other directly or indirectly gives a factor of $`\varphi ^{n1}`$. In the next two sections we will present two expansions considering different terms.
## IV Expansion up to first order in $`\varphi `$
In this section only terms up to $`O(\varphi )`$ will be taken into account.
### A Partition function
In $`O(\varphi )`$ the canonical partition function reads
$`Z`$ $`=`$ $`{\displaystyle \underset{k}{}e^{v_k}d}\stackrel{}{𝐱}d\stackrel{}{\mathrm{\Omega }}+`$ (34)
$`{\displaystyle \underset{k}{}e^{v_k}\underset{i<j}{}f_{ij}d}\stackrel{}{𝐱}d\stackrel{}{\mathrm{\Omega }}+O(\varphi ^2).`$
The $`f_{ij}`$ have yet to be expanded in powers of $`v_{ij}^{DD}`$. Figure 2 shows the corresponding graphs. There are $`N(N1)/2N^2/2`$ ways to choose $`i`$ and $`j`$. Because all particles are identical one can write
$`Z`$ $`=`$ $`{\displaystyle \underset{k}{}e^{v_k}d}\stackrel{}{𝐱}d\stackrel{}{\mathrm{\Omega }}+`$ (36)
$`{\displaystyle \frac{N^2}{2}}{\displaystyle \underset{k}{}e^{v_k}f_{12}d}\stackrel{}{𝐱}d\stackrel{}{\mathrm{\Omega }}+O(\varphi ^2).`$
Integrating over most degrees of freedom results in
$`Z`$ $`=`$ $`Z_0+{\displaystyle \frac{N^2}{2}}z_0^{N2}{\displaystyle e^{v_1v_2}f_{12}𝑑𝐱_1𝑑𝐱_2𝑑\mathrm{\Omega }_1𝑑\mathrm{\Omega }_2}`$ (38)
$`+O(\varphi ^2).`$
Here
$$Z_0=z_0^N;z_0=4\pi V\frac{\mathrm{sinh}\alpha _s}{\alpha _s}$$
(39)
is the partition function of a paramagnetic gas of noninteracting particles in the field $`H_s`$ defining the Langevin parameter
$$\alpha _s=\frac{mH_s}{kT}.$$
(40)
### B Expansion in the dipolar interaction
Now we expand $`f_{12}`$ appearing in (38) in a power series in $`ϵ`$. The $`n`$–th summand of this series contains integrals of the form
$$A_n=e^{v_1v_2}f_{12}^{(n)}𝑑𝐱_1𝑑𝐱_2𝑑\mathrm{\Omega }_1𝑑\mathrm{\Omega }_2.$$
(41)
$`A_0`$ is special. Here one gets
$`A_0`$ $`=`$ $`\left(4\pi {\displaystyle \frac{\mathrm{sinh}\alpha _s}{\alpha _s}}\right)^2{\displaystyle \left(e^{v_{12}^{HC}}1\right)𝑑𝐱_1𝑑𝐱_2}`$ (42)
$`=`$ $`{\displaystyle \frac{1}{V}}z_0^2{\displaystyle \left(e^{v_{12}^{HC}}1\right)𝑑𝐫_{12}}.`$ (43)
The integrand vanishes if $`r_{12}>D`$. Otherwise its value is $`1`$. Thus
$$A_0=\frac{4}{3}\pi \frac{D^3}{V}z_0^2,$$
(44)
or by expressing the result in terms of $`\varphi `$
$$A_0=\frac{8}{N}\varphi z_0^2.$$
(45)
For $`n1`$ we have
$`A_n`$ $`=`$ $`{\displaystyle \frac{1}{n!}}{\displaystyle e^{v_1v_2}}`$ (47)
$`\times \left(v_{12}^{DD}\right)^ne^{v_{12}^{HC}}d𝐱_1d𝐱_2d\mathrm{\Omega }_1d\mathrm{\Omega }_2.`$
Switching from $`𝐫_2`$ to the relative coordinate $`𝐫_{12}`$ and integrating over $`𝐫_1`$ gives a factor of $`V`$. Then $`𝐫_{12}`$ runs over the sphere volume. We introduce spherical coordinates, i. e.,
$`𝐦_1`$ $`=`$ $`m(\mathrm{cos}\phi _1\mathrm{sin}\vartheta _1,\mathrm{sin}\phi _1\mathrm{sin}\vartheta _1,\mathrm{cos}\vartheta _1)`$ (48)
$`𝐦_2`$ $`=`$ $`m(\mathrm{cos}\phi _2\mathrm{sin}\vartheta _2,\mathrm{sin}\phi _2\mathrm{sin}\vartheta _2,\mathrm{cos}\vartheta _2)`$ (49)
$`𝐫_{12}`$ $`=`$ $`r_{12}(\mathrm{cos}\phi \mathrm{sin}\vartheta ,\mathrm{sin}\phi \mathrm{sin}\vartheta ,\mathrm{cos}\vartheta ).`$ (50)
The direction of the magnetic field defines the $`z`$–axis. Then, the integral assumes the form
$`A_n`$ $`=`$ $`{\displaystyle \frac{V}{n!}}{\displaystyle e^{\alpha _s\mathrm{cos}\vartheta _1+\alpha _s\mathrm{cos}\vartheta _2}}`$ (53)
$`\times \left({\displaystyle \frac{m^2}{4\pi \mu _0kTr_{12}^3}}\right)^nP^n(\phi _1,\vartheta _1,\phi _2,\vartheta _2,\phi ,\vartheta )`$
$`\times e^{v_{12}^{HC}}r_{12}^2dr_{12}d\omega _{12}d\mathrm{\Omega }_1d\mathrm{\Omega }_2.`$
The new spherical angle $`\omega _{12}`$ represents $`\phi `$ and $`\vartheta `$. The exact form of the function $`P`$ is not important, but $`P`$ and therefore $`P^n`$ is a polynomial in the $`\mathrm{cos}`$ and $`\mathrm{sin}`$ of the six angles. Integration over four of them can be done analytically. Finally this can also be done for $`\vartheta _1`$ and $`\vartheta _2`$ by substituting $`u_{1,2}=\mathrm{cos}\vartheta _{1,2}`$. One gets an expression of the form
$$A_n=\frac{V}{n!}G_n^{}(\alpha _s)_D^{R_s}\left(\frac{m^2}{4\pi \mu _0kTr_{12}^3}\right)^nr_{12}^2𝑑r_{12}.$$
(54)
Here we have introduced the correct bounds of the last remaining integral explicitly. By setting the lower bound to $`D`$ we have incorporated the hard core factor. The upper bound is given by the cutoff radius $`R_s`$ for the near-field dipolar contribution. While the evaluation of $`G_n^{}`$ can be done analytically, it is quite difficult to do this by hand even for $`n=2`$. We therefore used the computer algebra system mathematica to perform the integrations. See Appendix A for the form of the $`G_n^{}`$.
For $`n2`$ one can safely set $`R_s=\mathrm{}`$ (see below). For $`n=1`$ this would result in a logarithmic divergence of the integral. But $`G_1^{}0`$ anyway, because the calculation of $`G_1^{}`$ involves an averaging over a dipolar field. So by using $`ϵ`$ and $`\varphi `$ one finally has
$$A_n=\frac{2V^2}{N\pi (n1)n!}G_n^{}(\alpha _s)ϵ^n\varphi n2,$$
(55)
$$A_1=0.$$
(56)
Note that $`A_1`$ vanishes only in our spherical configuration with finite $`R_s`$. The divergence of $`A_1`$ in the general, spatially unrestricted case is just an expression of fact that the dipolar forces are long range. By treating the distant parts of the ferrofluid as a continuum we incorporate any long-range effects and the resulting geometry dependence via the field $`H_s=H+M/3`$. Into this field enters the relation between the external and the macroscopic internal field.
Two further comments should be made here. A generalization of our calculation for central symmetric interactions other than a hard sphere potential is possible. It requires an analytical or numerical evaluation of integrals of the form $`r^{23n}e^{v^{SR}}𝑑r`$ in (54), with $`v^{SR}=V^{SR}/kT`$ and $`V^{SR}`$ denoting the $`r`$–dependent short range potential in question.
The second thing is that we can now make quantitative statements about how large the virtual sphere has to be chosen. To ensure $`A_1`$ to vanish unambigously in (54) and to introduce $`H_s`$ instead of $`H_e`$ as ”external” field $`R_s`$ has only to be finite. The larger $`R_s`$ the better is the modeling of large–distance particle correlations entering into $`A_n`$ for $`n>1`$. Taking the limit $`R_s\mathrm{}`$ as the final step in the calculation of the $`A_n`$ is therefore appropriate from this point of view. On the other hand, the requirement of uniformity of the fields $`𝐇`$ and $`𝐌`$, that allows us to write $`H_s=H+M/3`$ restricts $`R_s`$ to values below the scale on which $`H`$ and $`M`$ vary. If one would use a finite radius $`R_s`$ one would get instead of (55) for $`n2`$
$$A_n=\frac{2V^2}{N\pi (n1)n!}G_n^{}(\alpha _s)ϵ^n\varphi \left[1(D/R_s)^{3n3}\right]$$
(57)
which allows an error estimate: Consider a system where $`𝐌`$ and $`𝐇`$ do not vary on the scale of, say, $`\mu `$m. For ferrofluids, $`D10`$ nm. Choosing $`R_s=10D`$ or $`100D`$ is then both allowed and implies a difference in $`A_2`$ of about 0.1 percent. The result for $`R_s=100D`$ is better than for $`R_s=10D`$, because in the latter case particles in a distance range between 100 nm and 1 $`\mu `$m are treated in the continuum approximation and not correctly. But the error that is made by treating the ferrofluid as a continuum already beyond $`R_s=10D`$ is only about 0.1 percent. We can safely assume that the macroscopic, magnetic properties do not vary on this scale. Thus 100 nm is an appropriate medium scale on which both requirements hold: The continuum approximation works well beyond this cutoff radius and the macroscopic fields $`𝐇`$ and $`𝐌`$ should be constant on this scale. Except for the calculation of $`A_1`$, it is then possible to set $`R_s=\mathrm{}`$ in the calculations of the integrals.
Using the results (55), (56), (45), and (41) in (38), one gets the following expression for $`Z`$:
$$Z=Z_0\left[14N\varphi +N\varphi \underset{n=2}{\overset{\mathrm{}}{}}G_n(\alpha _s)ϵ^n\right]+O(\varphi ^2).$$
(58)
Here we introduced the functions
$$G_n(\alpha _s)=\frac{1}{16\pi ^3(n1)n!}\left(\frac{\alpha _s}{\mathrm{sinh}\alpha _s}\right)^2G_n^{}(\alpha _s),$$
(59)
some of which are given in Appendix A.
### C Free energy and magnetization
The next step is to compute the free energy
$`{\displaystyle \frac{F}{kT}}`$ $`=`$ $`\mathrm{ln}Z=N\mathrm{ln}z_0`$ (62)
$`\mathrm{ln}\left[14N\varphi +N\varphi {\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}G_n(\alpha _s)ϵ^n\right]`$
$`+O(\varphi ^2).`$
In $`O(\varphi )`$, we can use $`\mathrm{ln}(1+x)=1+x`$ here:
$`{\displaystyle \frac{F}{kT}}`$ $`=`$ $`N\mathrm{ln}z_0+4N\varphi N\varphi {\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}G_n(\alpha _s)ϵ^n`$ (64)
$`+O(\varphi ^2).`$
The magnetization turns out to be
$`M_{sphere}(\alpha _s)`$ $`=`$ $`{\displaystyle \frac{1}{\mu _0V}}{\displaystyle \frac{F}{H_s}}={\displaystyle \frac{m}{\mu _0VkT}}{\displaystyle \frac{F}{\alpha _s}}`$ (65)
$`=`$ $`{\displaystyle \frac{Nm}{\mu _0V}}\left[(\alpha _s)+\varphi {\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}G_n^{}(\alpha _s)ϵ^n\right]`$ (67)
$`+O(\varphi ^2).`$
The leading term is the Langevin function $``$ times the saturation magnetization $`M_{sat}=\frac{Nm}{\mu _0V}`$ of the fluid.
In order to determine $`M(H)`$ we identify, according to (11), $`M_{sphere}(\alpha _s)`$ with $`M(\alpha )`$, i. e.,
$`{\displaystyle \frac{M}{M_{sat}}}`$ $`=`$ $`\left(\alpha +{\displaystyle \frac{mM}{3kT}}\right)`$ (69)
$`+\varphi {\displaystyle \underset{n=2}{\overset{\mathrm{}}{}}}G_n^{}\left(\alpha +{\displaystyle \frac{mM}{3kT}}\right)ϵ^n+O(\varphi ^2),`$
where $`\alpha `$ is the usual Langevin parameter,
$$\alpha =\frac{mH}{kT}.$$
(70)
Instead of trying to find the function $`M(\alpha )`$ that solves this equation exactly we expand the functions $``$ and $`G_n^{}`$ for small $`\varphi `$ into a series around $`M=0`$ and reinsert it on the right hand side. Using the fact that $`\frac{mM_{sat}}{3kT}=8\varphi ϵ`$ grows linearly in $`\varphi `$ we arrive at
$$\frac{M(\alpha )}{M_{sat}}=L_{0,0}+\varphi \underset{n=1}{\overset{\mathrm{}}{}}L_{1,n}ϵ^n+O(\varphi ^2)$$
(72)
with
$`L_{0,0}`$ $`=`$ $`(\alpha )`$ (73)
$`L_{1,1}`$ $`=`$ $`8(\alpha )^{}(\alpha )`$ (74)
$`L_{1,n}`$ $`=`$ $`G_n^{}(\alpha )\text{for}n2.`$ (75)
This is a consistent approximation in terms of $`\varphi `$. On the other hand, solving Eq. (69) in a formally exact manner for $`M`$ would introduce higher orders in $`\varphi `$ which we already neglected to arrive at Eq. (69).
Note also that $`M_{sphere}(\alpha _s)/M_{sat}`$ (65) contains explicitly a term $`\varphi ϵ^2`$ as lowest nontrivial power coming from the expansion in the near-field dipolar coupling strength. On the other hand the self-consistent solution (70) that solves Eq. (69) starts out with a contribution $`\varphi ϵ`$. The latter arises from the far-field dipolar continuum via the magnetization $`M`$ in the dipole-induced shift of the argument, $`\alpha +\frac{mM}{3kT}`$, of the Langevin function in Eq. (69) — in the absence of any dipolar interactions in the system one would have $`H_s=H_e=H`$ leading to ideal paramagnetism.
### D Comparison with previous results
The Onsager model, the Weiss model, and our calculation agree that up to order $`\varphi ϵ`$
$$\frac{M}{M_{sat}}=(\alpha )+8\varphi ϵ(\alpha )^{}(\alpha )+O(\varphi ^2)+O(ϵ^2).$$
(76)
This expression was also derived by Buyevich and Ivanov with a calculation similar to ours. However, they did not introduce a magnetic continuum approximation. Instead, they assumed a special probe geometry of a long cylinder parallel to the external magnetic field and performed an integration over all the particle’s dipolar fields in the cylinder explicitly. The magnetization was therefore given in terms of the external field. Their result agrees with ours because for the cylindrical geometry chosen in $`𝐇_e`$ equals $`𝐇`$.
A second paper that deals with our problem in a similar way was published by Kalikmanov . In section 4, the author arrives at an equation for the magnetization that reads in our notation
$$\frac{M}{M_{sat}}=(\alpha )+3\varphi ϵ^2G_2^{}(\alpha )_1^{\mathrm{}}\frac{g_0(x)}{x^4}𝑑x.$$
(77)
Here $`g_0(x)`$ is the hard sphere correlation function. In our $`O(\varphi )`$–approximation this function has to be set to one. Then the $`\varphi ϵ^2`$–term agrees with ours. Note, however, that the above result (77) of Kalikmanov does not contain the $`\varphi ϵ`$–term resulting from the magnetic field from the continuum.
## V Expansion up to second order in $`\varphi `$ and $`ϵ`$
It is possible to calculate $`O(\varphi ^2)`$–terms of the Born–Mayer expansion when $`ϵ`$ is taken into account up to second order only. A more elegant way to calculate the magnetization in this approximation makes use of the grand canonical rather than the canonical ensemble. This approach allows to avoid the determination of some terms that can be factorized into already known integrals and cancel out in the calculation of the free energy. However, the grand canonical approach has the disadvantage that it yields the magnetization as function of the chemical potential $`\mu `$ rather than the particle number $`N`$. Some more algebra is then required to find out the function $`\mu (N)`$. Here we continue to work with the canonical ensemble.
### A The graphs
Figure 3 shows the 12 additional graphs that are of second order in $`\varphi `$ and of less than third order in $`ϵ`$. Four of them vanish because they contain at least one first–order dipolar interaction term between otherwise unrelated particles. Integration over the relative position of these particles while leaving the relative positions between all other particles and the direction of the magnetic moments fixed yields zero since it involves a spatial averaging over a dipolar field. The graph labelled with the letter F vanishes for similar reasons that are explained in Appendix B where we calculate the integrals one by one. Their respective contribution to the partition function is
$`Z_A/Z_0`$ $`=`$ $`32N\varphi ^2`$ (79)
$`Z_B/Z_0`$ $`=`$ $`16N\varphi ^2ϵ^2G_2(\alpha _s)`$ (80)
$`Z_C/Z_0`$ $`=`$ $`8(N^26N)\varphi ^2`$ (81)
$`Z_D/Z_0`$ $`=`$ $`4(N^26N)\varphi ^2ϵ^2G_2(\alpha _s)`$ (82)
$`Z_E/Z_0`$ $`=`$ $`5N\varphi ^2`$ (83)
$`Z_F/Z_0`$ $`=`$ $`0`$ (84)
$`Z_G/Z_0`$ $`=`$ $`{\displaystyle \frac{1+6\mathrm{ln}2}{4}}N\varphi ^2ϵ^2G_2(\alpha _s)`$ (85)
$`Z_H/Z_0`$ $`=`$ $`N\varphi ^2ϵ^2K(\alpha _s).`$ (86)
The functions $`G_2`$ and $`K`$ are given in Appendix A.
### B Free energy and magnetization
Now we have all necessary terms at hand to calculate the canonical partition function up to the desired order:
$`{\displaystyle \frac{Z}{Z_0}}=14(N1)\varphi +(N1)\varphi ϵ^2G_2(\alpha _s)`$ (87)
$`+32N\varphi ^216N\varphi ^2ϵ^2G_2(\alpha _s)+8(N^26N)\varphi ^2`$ (88)
$`4(N^26N)\varphi ^2ϵ^2G_2(\alpha _s)5N\varphi ^2`$ (89)
$`+{\displaystyle \frac{1+6\mathrm{ln}2}{4}}N\varphi ^2ϵ^2G_2(\alpha _s)N\varphi ^2ϵ^2K(\alpha _s)+H.O.T..`$ (90)
The terms in $`O(\varphi )`$ appear already in (58). They are presented here including the next higher order in $`N`$. The other terms come from $`Z_A`$$`Z_H`$. To include all terms of $`O(\varphi ^2,ϵ^2)`$ in the free energy one has to approximate the logarithm $`\mathrm{ln}(1+x)`$ by $`xx^2/2`$. The quadratic order is necessary only for the $`O(\varphi )`$–terms. New terms of $`O(N^2)`$ appear and cancel against those of the terms from $`Z_C`$ and $`Z_D`$. One gets
$`{\displaystyle \frac{F}{kT}}=N\mathrm{ln}z_0+4N\varphi +5N\varphi ^2`$ (91)
$`N\varphi ϵ^2G_2(\alpha _s){\displaystyle \frac{1+6\mathrm{ln}2}{4}}N\varphi ^2ϵ^2G_2(\alpha _s)`$ (92)
$`+N\varphi ^2ϵ^2K(\alpha _s)+H.O.T..`$ (93)
The result is proportional to $`N`$ as it has to be.
The magnetization of the sphere is
$`{\displaystyle \frac{M_{sphere}(\alpha _s)}{M_{sat}}}=(\alpha _s)+\varphi ϵ^2G_2^{}(\alpha _s)`$ (94)
$`+{\displaystyle \frac{1+6\mathrm{ln}2}{4}}\varphi ^2ϵ^2G_2^{}(\alpha _s)\varphi ^2ϵ^2K^{}(\alpha _s)+H.O.T..`$ (95)
To calculate the magnetization as a function of $`\alpha `$ we identify (94) with $`M`$ and use again $`\alpha _s=\alpha +\frac{mM}{3kT}`$. The right hand side of (94) has now to be expanded around $`\alpha `$ up to second order and the resulting equation has to be iterated twice to take into account all important terms up to $`ϵ^2\varphi ^2`$. The result is
$$\frac{M(\alpha )}{M_{sat}}=L_{0,0}+\varphi ϵL_{1,1}+\varphi ϵ^2L_{1,2}+\varphi ^2ϵ^2L_{2,2}+\mathrm{}$$
(97)
with $`L_{0,0},L_{1,1}`$, and $`L_{1,2}`$ defined in eqs. (73) - (75) and
$`L_{2,2}`$ $`=`$ $`64(\alpha )^{}(\alpha )^2+32(\alpha )^{\prime \prime }(\alpha )`$ (99)
$`+{\displaystyle \frac{1+6\mathrm{ln}2}{4}}G_2^{}(\alpha )K^{}(\alpha ).`$
For the discussion in the next Sec. we decompose
$$L_{2,2}(\alpha )=L_{2,2}^{sphere}(\alpha )+L_{2,2}^{iterative}(\alpha ).$$
(101)
The function
$$L_{2,2}^{sphere}=\frac{1+6\mathrm{ln}2}{4}G_2^{}K^{}$$
(102)
occurs already in the expression (94) for the magnetization $`M_{sphere}(\alpha _s)`$ of the sphere. The contribution
$$L_{2,2}^{iterative}=64(^{})^2+32^{\prime \prime }$$
(103)
arises in obtaining the selfconsistent solution of the equation $`M=M_{sphere}`$ with an expansion and iteration.
## VI Discussion of the results
We will first show that our result (94) for $`M_{sphere}(H_s)`$ does not lead to a ferromagnetic solution in contradistinction to the Weiss model. Then we discuss the behavior of the different terms contributing to (70) and to (V B) and we delineate the range of reliability of the simplest approximation. Finally, we address problems arising when comparing with experiments.
### A Spontaneous magnetization?
Investigations based on density functional methods by Groh and Dietrich and on Monte Carlo methods by Weis and Levesque provided support for the existence of magnetized phases for absent external field $`H_e`$, i. e., ferromagnetism, in the system of dipolar hard spheres we consider in this work. Groh and Dietrich consider a ferrofluid probe of needle–like shape where $`H=H_e`$ and find a transition to a magnetized phase at $`\varphi ϵ0.35`$. But they consider this value as being overestimated and refer to . Weis and Levesque study a case without demagnetizing fields, i. e., again $`H=H_e`$. They find a transition to a magnetized phase at $`ϵ=6.25`$ for $`\varphi 0.35`$. As discussed in detail below, these values are outside the range of reliability of our results.
The Weiss model does also show ferromagnetic behavior. It is recovered from (94) by keeping only the leading–order term $`(\alpha _s)`$ describing a single moment in the field $`H_s=H+M/3`$. The resulting self–consistency equation
$$M=M_{sphere}^{Weiss}\left(H+\frac{M}{3}\right)=M_{sat}\left[\frac{m}{kT}\left(H+\frac{M}{3}\right)\right]$$
(104)
allows for zero field a solution with finite magnetization when $`kT<mM_{sat}/9`$. Using (4) combined with (12) and (13) this condition is equivalent to $`\varphi ϵ>3/8`$, about the same value as in . So according to the Weiss model the ferrofluid will show ferromagnetic behavior below a critical temperature that grows linearly with the saturation magnetization $`M_{sat}`$ of the ferrofluid. But even for a ferrofluid consisting of cobalt particles with a magnetic core diameter of 10 nm and a magnetic volume fraction of $`\varphi _{mag}=0.1`$ the critical temperature would be as low as 90 K.
While the transition combination $`ϵ=6.25`$, $`\varphi 0.35`$ of is outside the range of reliability of our results, the threshold location $`\varphi ϵ=8/3`$ of the Weiss model may be not. However, in agreement with we do not find selfconsistent ferromagnetic solutions of the equation (94) $`M=M_{sphere}(H+M/3)`$ within this range. We have numerically confirmed that for $`H=0`$ the equation $`M=M_{sphere}(M/3)`$ allows always only the trivial solution $`M=0`$.
### B Contribution from different orders
Now we will take a closer look on the functions of $`\alpha `$ involved in (70) and (V B). All these functions are odd as it has to be for reasons of symmetry. For $`\alpha \mathrm{}`$ they vanish as $`1/\alpha ^2`$ or faster. Because $`1(\alpha )1/\alpha `$ that means that the predicted magnetization is always smaller than the saturation magnetization for $`\alpha \mathrm{}`$. Nevertheless the magnetization can assume unphysical values $`>M_{sat}`$ for intermediate $`\alpha `$ if $`ϵ`$ or $`\varphi `$ is big enough for the approximations to become invalid.
#### 1 Behavior in linear order of $`\varphi `$
We will first discuss the result (70) for the magnetization that was obtained up to linear order in the volume fraction $`\varphi `$. In figure 4 the functions $`L_{1,1}`$ and $`L_{1,2}`$ are plotted. The values of the higher–order functions are smaller, but their shape remains more or less the same as the logarithmic plot in figure 5 shows. Because $`L_{1,n}`$ and $`L_{1,n+2}`$ differ by about one order of magnitude one can conclude that by including higher and higher orders of $`ϵ`$ the series (V B) for the magnetization converges, as long as $`ϵ`$ is smaller than $`3`$. For this large value of $`ϵ`$ strong agglomeration can already be expected.
For small $`\alpha `$, $`L_{1,n}`$ is proportional to $`\alpha `$ ($`\alpha ^3`$) for odd (even) $`n`$. The initial susceptibility can therefore be written as
$`\chi (H=0)=`$ (105)
$`\chi _0(H=0)\left[1+\varphi {\displaystyle \underset{n=0}{}}s_{1,2n+1}ϵ^{2n+1}+O(\varphi ^2)\right].`$ (106)
Here
$$\chi _0(H=0)=\frac{mM_{sat}}{3kT},$$
(107)
is the initial susceptibility of the ideal paramagnetic gas, and the nonvanishing $`s_{1,n}`$ we calculated are
$`s_{1,1}={\displaystyle \frac{8}{3}};s_{1,3}={\displaystyle \frac{8}{75}};s_{1,5}={\displaystyle \frac{32}{3675}}`$ (108)
$`s_{1,7}={\displaystyle \frac{8}{19845}};s_{1,9}={\displaystyle \frac{148}{12006225}}.`$ (109)
Figure 6 shows $`\chi _0(H=0)`$ (thick dashed line), and the susceptibility $`\chi (H=0)`$ (106) including progressive orders $`\varphi ϵ`$, $`\varphi ϵ^3`$, $`\varphi ϵ^5`$, $`\varphi ϵ^7`$, and $`\varphi ϵ^9`$ (thin dashed lines, from bottom to top) as a function of $`ϵ`$ for $`\varphi =0.15`$. The sequence of these thin dashed lines shows that this series converges in the $`ϵ`$–range of figure 6. The last thick full line in Figure 6 represents $`\chi (H=0)`$ including the contributions in order $`\varphi ^2ϵ^2`$. It shows that the latter are even for $`\varphi =0.15`$ not yet important.
#### 2 Behavior in second order of $`\varphi `$
Now we take a look at the functions $`L_{2,2}^{sphere}`$ (V Bb) and $`L_{2,2}^{iterative}`$ (V Bc) that add up to $`L_{2,2}`$ (V Ba) which enters in order $`\varphi ^2ϵ^2`$ into the magnetization (V Ba).
Figure 7 shows that the contributions $`L_{2,2}^{sphere}`$ and $`L_{2,2}^{iterative}`$ almost cancel each other at small $`\alpha `$. This is why the influence of the $`\varphi ^2ϵ^2`$–terms on the susceptibility in figure 6 is so small. However, at higher $`\alpha `$ the $`\varphi ^2ϵ^2L_{2,2}`$–term becomes important. Comparing the latter with the linear one, $`\varphi ϵL_{1,1}`$, one finds that they contribute for $`ϵ\varphi 0.5`$ equally at larger $`\alpha `$.
Except for very small $`\alpha `$ $`L_{2,2}`$ is negative, because it includes higher–order particle position correlations that result in a better modeling of the distance distribution due to the finite size of the particles. The mean distance is bigger in this approximation and the induced dipolar fields at the particle positions are therefore smaller.
The influence of the $`\varphi ^2ϵ^2L_{2,2}`$ contribution to the magnetization is shown in figure 8 for $`ϵ=2`$ and $`\varphi =0.05`$. For these parameters this term is already large enough to cancel almost exactly the sum of all contributions $`L_{1,n}\varphi ϵ^n`$, with $`n2`$ from the linear order in $`\varphi `$ at moderate $`\alpha `$. Figure 9 shows the susceptibility $`\chi (H)=\frac{M(H)}{H}`$ for the same parameters. At higher $`\alpha `$, the cancellation of the higher $`L_{1,n}`$–terms against the $`L_{2,2}`$–contribution can again be seen. At smaller $`\alpha `$, however, the behavior is different. There the contribution of the $`L_{1,n}`$–terms is much larger, whereas the $`L_{2,2}`$–contributions vanish.
### C Reliability of the $`O(\varphi ϵ)`$–approximation
We can determine the range of reliability of the simplest approximation
$`{\displaystyle \frac{M}{M_{sat}}}`$ $`=`$ $`L_{0,0}(\alpha )+L_{1,1}(\alpha )\varphi ϵ`$ (110)
$`=`$ $`(\alpha )+8(\alpha )^{}(\alpha )\varphi ϵ`$ (111)
to the magnetization that includes effects of dipolar interactions since we know the higher–order corrections in $`\varphi `$ as well as in $`ϵ`$. To that end we investigated the ratios
$$\left|\frac{O(\varphi ϵ^n)\text{–terms (}\text{n ¿ 1}\text{)}}{L_{0,0}(\alpha )+L_{1,1}(\alpha )\varphi ϵ}\right|,$$
(112)
and
$$\left|\frac{O(\varphi ^2)\text{–terms}}{L_{0,0}(\alpha )+L_{1,1}(\alpha )\varphi ϵ}\right|.$$
(113)
The first ratio assumes its maximum at $`\alpha =0`$, that means the initial susceptibility is most sensitive to higher order corrections in $`ϵ`$. The second ratio (113) assumes its maximum around $`\alpha =2`$, that is near the maximum of the absolute value of the numerator (as seen in figure 7).
In the $`ϵ`$$`\varphi `$ plane of figure 10(a) we show isolines of the maximal – with respect to $`\alpha `$ – ratio (112) and figure 10(b) shows the analogous isolines for the ratio (113). The comparison shows that the smallness of $`ϵ`$ is more important to keep the ratio (112) small, whereas in (113) the value of $`\varphi `$ is also important. As rules of thumb one can say that the approximation (111) is valid within about 1 – 2 percent if $`ϵ<1`$ and $`ϵ\varphi <0.04`$. If the first constraint is not fulfilled, higher orders in $`ϵ`$ have to be taken into account. Higher orders in $`\varphi `$ are needed if the second constraint is not fulfilled.
### D Comparison with experiments?
There are several papers that aim at investigating the influence of dipolar interactions on the magnetization by comparing theoretical models developed so far with experimental magnetizations of ferrofluids. The mean spherical model was reported to show good agreement with experiments. Pshenichnikov found also good agreement with the high temperature approximation , i. e., the approximation (111). But this ansatz failed in the magneto–granulometric analysis done in .
We do not present a comparison of our results with the experiments on the magnetization in the literature because of several problems: In our theory it is necessary to distinguish between the particle diameter $`D`$ and the magnetic core diameter $`D_{mag}`$ that is found in magneto–granulometric measurements. This problem does not arise in the mean spherical model or the high temperature approximation, where $`ϵ`$ and $`\varphi `$ enter only via the factor $`\varphi ϵ=\frac{Nm^2}{24V\mu _0kT}=\frac{M_{sat}m}{24kT}`$ that is independent of $`D`$. Also, corrections such as the temperature dependence of the saturation magnetization or the fluid density should be taken into account .
But the major problem in comparing directly with experiments is that our theory does not take into account the polydispersity of ferrofluids. The effect of polydispersity is significant already in the absence of any dipolar interaction. This can be inferred from the dashed and the dotted curves in Fig. 11 representing the reduced magnetization of noninteracting magnetic particles having a polydisperse and a monodisperse distribution of particle diameters, respectively. Here the common particle diameters of the latter is $`\overline{D^3}^{1/3}`$, where $`\overline{D^3}`$ is the third moment of the particle size distribution
$$P(D):=\frac{1}{\sqrt{2\pi }\sigma D_0e^{\sigma ^2/2}}e^{\frac{\mathrm{ln}^2D/D_0}{2\sigma ^2}}.$$
(114)
of the former. The mean magnetic moment $`\overline{m}`$ and the saturation magnetization of the two systems are the same. For comparison with the effect of dipolar interaction in monodisperse systems the full curve in Fig. 11 shows our result for $`M`$ (V Ba) including all terms $`\varphi ϵ^n`$ and the term $`\varphi ^2ϵ^2`$. Hence the effects of polydispersiveness alone, i. e., without interaction are comparable in size with the effect of dipolar interactions in monodisperse systems. Thus clearly an extension of the here presented Born–Mayer expansion method to the case of polydisperse interacting particles is desirable.
## VII Conclusion
We calculated the free energy and in particular the magnetization $`M`$ of a ferrofluid as a function of the macroscopic magnetic field $`H`$. To do so, we used the technique of the Born–Mayer expansion together with an expansion in terms of the dipolar coupling energy. The magnetic particles were assumed to be hard spheres with a common hard core diameter $`D`$ and magnetic moment $`m`$ that interact via long range dipolar interactions. This feature may result in a geometry dependence of thermodynamic properties. We treated this problem by dividing the dipolar field at some position $`𝐱_i`$ that is produced by the magnetic moments of the particles into a near–field and into a far–field part depending on whether the particle distance from $`𝐱_i`$ is larger than some radius $`R_s`$ or not. In this way every magnetic particle is imagined to be located in the center of a sphere of radius $`R_s`$. The far–field dipolar contribution from particles beyond $`R_s`$ is then replaced by a magnetic continuum with magnetization $`M`$ and magnetic field $`H`$. Here $`R_s`$ is chosen to be such that $`M`$ and $`H`$ are homogeneous on the scale of $`R_s`$. The magnetic continuum outside the sphere produces in the center of the sphere the magnetic field $`H_s=H+M/3`$. This field acts as an ”external” field on the particle in the center of the sphere. The near–field interaction of the latter with the other particles within the sphere being at a distance smaller than $`R_s`$ is treated explicitly. Thus in our statistical mechanical calculations there appear dipolar interactions only with interparticle distances less than $`R_s`$. However, since the cutoff dependence of the relevant expressions occurring in these calculations is negligible already beyond a radius of the order of $`10D100`$ nm we used $`R_s=\mathrm{}`$ in these expressions.
The expansion of the partition function for these interacting particles in terms of the volume ratio and the dipolar coupling strength $`ϵ`$ yields an expression for the magnetization
$$M_{sphere}=M_{sphere}(H+M/3).$$
(115)
as a function of the ”external” part of the field inside the sphere. The magnetization $`M_{sphere}`$ is then identified with the magnetization $`M(H)`$ inside the continuum so that a selfconsistent relation results. The aforementioned geometry dependence of $`M`$ in the general case is incorporated via $`H`$.
We presented two different expansions in $`ϵ`$ and $`\varphi `$, one containing only linear terms in $`\varphi `$, the other also second order $`\varphi `$ terms, but only up to $`O(ϵ^2)`$. We discussed the range of applicability in the $`\varphi `$$`ϵ`$ plane of their results for $`M(H)`$ and compared them to the most simple approximation to the magnetization that contains the dipolar effects only in linear order in $`ϵ`$ and $`\varphi `$. The selfconsistent relation for $`M(H)`$ that contains only up to second order terms in both parameters does not admit a ferromagnetic solution with spontaneous magnetization. Finally we showed that an extension to polydisperse interacting particles is desirable.
###### Acknowledgements.
This work was supported by the Deutsche Forschungsgemeinschaft (SFB 277).
## A The functions $`G_n`$ and $`K`$
The functions $`G_n^{}(x)`$ in (54) are related to $`G_n(x)`$ via (59):
$$G_n(x)=\frac{1}{16\pi ^3(n1)n!}\left(\frac{x}{\mathrm{sinh}x}\right)^2G_n^{}(x).$$
(A1)
The functions $`G_n(x)`$ introduced in (59) have the form
$`G_n(x)`$ $`=`$ $`G_n^{(0)}\left({\displaystyle \frac{1}{x}}\right)+G_n^{(1)}\left({\displaystyle \frac{1}{x}}\right)\mathrm{coth}x`$ (A3)
$`+G_n^{(2)}\left({\displaystyle \frac{1}{x}}\right)\mathrm{coth}^2x,`$
where the functions $`G_n^{(i)}(y)`$ are polynomials. The first four triple are given by
$`G_2^{(0)}(y)`$ $`=`$ $`{\displaystyle \frac{8}{5}}+{\displaystyle \frac{8}{5}}y^2+{\displaystyle \frac{12}{5}}y^4`$ (A5)
$`G_2^{(1)}(y)`$ $`=`$ $`{\displaystyle \frac{8}{5}}y{\displaystyle \frac{24}{5}}y^3`$ (A6)
$`G_2^{(2)}(y)`$ $`=`$ $`{\displaystyle \frac{12}{5}}y^2`$ (A7)
$`G_3^{(0)}(y)`$ $`=`$ $`{\displaystyle \frac{4}{35}}y^2{\displaystyle \frac{48}{35}}y^4{\displaystyle \frac{12}{7}}y^6`$ (A9)
$`G_3^{(1)}(y)`$ $`=`$ $`{\displaystyle \frac{8}{35}}y+{\displaystyle \frac{8}{5}}y^3+{\displaystyle \frac{24}{7}}y^5`$ (A10)
$`G_3^{(2)}(y)`$ $`=`$ $`{\displaystyle \frac{16}{105}}{\displaystyle \frac{8}{35}}y^2{\displaystyle \frac{12}{7}}y^4`$ (A11)
$`G_4^{(0)}(y)`$ $`=`$ $`{\displaystyle \frac{8}{105}}+{\displaystyle \frac{8}{35}}y^2+{\displaystyle \frac{92}{35}}y^4+{\displaystyle \frac{72}{7}}y^6+12y^8`$ (A13)
$`G_4^{(1)}(y)`$ $`=`$ $`{\displaystyle \frac{16}{105}}y{\displaystyle \frac{8}{5}}y^3{\displaystyle \frac{88}{7}}y^524y^7`$ (A14)
$`G_4^{(2)}(y)`$ $`=`$ $`{\displaystyle \frac{32}{105}}y^2+{\displaystyle \frac{16}{7}}y^4+12y^6`$ (A15)
$`G_5^{(0)}(y)`$ $`=`$ $`{\displaystyle \frac{12}{385}}y^2{\displaystyle \frac{208}{385}}y^4{\displaystyle \frac{852}{77}}y^6{\displaystyle \frac{480}{11}}y^8{\displaystyle \frac{540}{11}}y^{10}`$ (A17)
$`G_5^{(1)}(y)`$ $`=`$ $`{\displaystyle \frac{8}{231}}y+{\displaystyle \frac{16}{385}}y^3+{\displaystyle \frac{472}{77}}y^5+{\displaystyle \frac{600}{11}}y^7+{\displaystyle \frac{1080}{11}}y^9`$ (A18)
$`G_5^{(2)}(y)`$ $`=`$ $`{\displaystyle \frac{16}{1155}}+{\displaystyle \frac{16}{1155}}y^2{\displaystyle \frac{40}{77}}y^4{\displaystyle \frac{120}{11}}y^6{\displaystyle \frac{540}{11}}y^8.`$ (A19)
All functions $`G_n^{(i)}(x)`$ have a well defined limit for $`x0`$ although this is not obvious for the above explicit expressions. Their values at $`x=0`$ are closely related to the coefficients in the $`ϵ`$–expansion of the second virial coefficient for the system of dipolar hard spheres in the absence of a magnetic field. The calculation of this coefficient dates back to and can also be found in .
The function $`K`$ (B31) that appears in the $`O(\varphi ^2)`$–terms of the free energy is given by
$`K(x)`$ $`=`$ $`{\displaystyle \frac{6}{x}}\mathrm{coth}^3x+\left({\displaystyle \frac{18}{x^2}}+12\right)\mathrm{coth}^2x`$ (A20)
$``$ $`\left({\displaystyle \frac{18}{x^3}}+{\displaystyle \frac{24}{x}}\right)\mathrm{coth}x+{\displaystyle \frac{6}{x^4}}+{\displaystyle \frac{12}{x^2}}.`$ (A21)
## B Graphs in second order of $`\varphi `$
Here we determine the contribution to the canonical partition function from the graphs A–H shown in Fig. 3. There often appear hard core interaction terms that are just expressions of the requirement that two particles have to or must not overlap. We define two abbreviations:
$$e^{v_{ij}^{HC}}1=O_{ij},$$
(B1)
$$e^{v_{ij}^{HC}}=\overline{O}_{ij}.$$
(B2)
### 1 Graph A
The graph A represents $`f_{ij}^{(0)}f_{ik}^{(0)}`$. There are $`N^3`$ ways to choose the constituting particles. But because $`j`$ and $`k`$ are equivalent only $`N^3/2`$ distinctive graphs remain. Integration over all variables except the positions of particle $`j`$ and $`k`$ relative to $`i`$ yields
$`{\displaystyle \frac{N^3}{2}}{\displaystyle f_{12}^{(0)}f_{13}^{(0)}\underset{l}{}e^{v_l}d}\stackrel{}{𝐱}d\stackrel{}{\mathrm{\Omega }}`$ (B3)
$`=`$ $`{\displaystyle \frac{N^3}{2}}\left(4\pi {\displaystyle \frac{\mathrm{sinh}\alpha _s}{\alpha _s}}\right)^NV^{N2}{\displaystyle O_{12}O_{13}𝑑𝐫_{12}𝑑𝐫_{13}}.`$ (B4)
The remaining integral factorizes and we can make use of the results for $`A_0`$ (45). The contribution of graph A to the partition function is
$$Z_A=32NZ_0\varphi ^2$$
(B5)
where $`Z_0`$ is given by (39).
### 2 Graph B
The graph B represents $`f_{ij}^{(0)}f_{ik}^{(2)}`$. All three particles appear in different ways, thus there are $`N^3`$ different graphs. After integration over the degrees of freedom of the noninvolved particles and switching to relative coordinates with respect to particle $`i`$ the integral factorizes again and one can make use of the results for $`A_0`$ (45) and $`A_2`$ (55). We get
$$Z_B=16NZ_0\varphi ^2ϵ^2G_2(\alpha _s).$$
(B6)
### 3 Graph C
The graph C represents $`f_{ij}^{(0)}f_{kl}^{(0)}`$. Here we have also to include the next higher order term when we calculate the number of combinations to get the $`O(N)`$–terms in the final result: There are $`(N^46N^3)/8`$ different terms. The integral for graph C can be factorized so that
$$Z_C=8(N^26N)Z_0\varphi ^2.$$
(B7)
### 4 Graph D
The graph D represents $`f_{ij}^{(0)}f_{kl}^{(2)}`$. The calculation is similar to the calculation of graph C. Again, we need the next higher order term in $`N`$. There are $`(N^46N^3)/4`$ combinations, twice as much as for graph C because the pairs $`(i,j)`$ and $`(k,l)`$ are not identical. One gets
$$Z_D=4(N^26N)Z_0\varphi ^2ϵ^2G_2(\alpha _s).$$
(B8)
### 5 Graph E
The integral containing the term $`f_{ij}^{(0)}f_{jk}^{(0)}f_{ki}^{(0)}`$ is the first really new integral. It involves only hard core interactions and does not contribute to the final expression for the magnetization. But for completeness we will calculate it also. The trivial integrations yield
$$Z_E=\frac{N^3}{6}\left(4\pi \frac{\mathrm{sinh}\alpha _s}{\alpha _s}\right)^NV^{N2}O_{12}O_{13}O_{23}𝑑𝐫_{12}𝑑𝐫_{13}.$$
(B9)
We keep the distance $`𝐫_{12}`$ fixed. The center of particle 3 has then to be inside two spheres of radius $`D`$ around particle 1 and 2. Integrating over the position of particle 3 yields the overlap volume $`V_o`$ of the two spheres
$$V_o=\frac{4}{3}\pi D^3\left[1\frac{3}{4}\frac{r_{12}}{D}+\frac{1}{16}\left(\frac{r_{12}}{D}\right)^3\right].$$
(B10)
Therefore
$`Z_E`$ $`=`$ $`{\displaystyle \frac{N^3}{6}}{\displaystyle \frac{Z_0}{V^2}}{\displaystyle O_{12}V_o(r_{12})𝑑𝐫_{12}}`$ (B11)
$`=`$ $`{\displaystyle \frac{N^3}{6}}{\displaystyle \frac{Z_0}{V^2}}4\pi {\displaystyle _0^D}V_o(r_{12})r_{12}^2𝑑r_{12}.`$ (B12)
Performing the last integration results in
$`Z_E`$ $`=`$ $`5NZ_0\varphi ^2.`$ (B13)
### 6 Graph F
The graph F represents $`f_{ij}^{(0)}f_{ik}^{(0)}f_{jk}^{(1)}`$. As already stated this integral vanishes which can be seen as follows: Consider an arbitrary configuration belonging to some value of the integrand
$$e^{v_iv_jv_k}f_{ij}^{(0)}f_{ik}^{(0)}f_{jk}^{(1)}.$$
(B14)
While leaving the direction of the magnetic moments fixed the whole configuration can be freely rotated around particle $`j`$ changing only the $`f_{jk}^{(1)}`$ term. Integration over the resulting configurations involves again an averaging over a dipolar field on a spherical surface.
### 7 Graph G
The calculation of the $`N^3/2`$ integrals belonging to $`f_{ij}^{(0)}f_{ik}^{(0)}f_{jk}^{(2)}`$ is similar to the calculation for graph E. First we integrate over all degrees of freedom except the distance between particle $`j=1`$ and $`k=2`$ and the position of particle $`i=3`$:
$`Z_G`$ $`=`$ $`{\displaystyle \frac{N^3}{2}}\left(4\pi {\displaystyle \frac{\mathrm{sinh}\alpha _s}{\alpha _s}}\right)^NV^{N2}\pi G_2(\alpha _s)ϵ^2D^6`$ (B16)
$`\times {\displaystyle }O_{13}O_{23}\overline{O}_{12}r_{12}^4dr_{12}d𝐫_3.`$
Integrating over $`𝐫_3`$ results again in an overlap volume term:
$$Z_G=\frac{N^3}{2}\frac{Z_0}{V^2}\pi G_2(\alpha _s)ϵ^2D^6\overline{O}_{12}V_o(r_{12})r_{12}^4𝑑r_{12}.$$
(B17)
Here the lower integration boundary is $`r_{12}=D`$ because of the remaining hard core factor. The upper integration boundary is $`r_{12}=2D`$ because the possibility that particle $`3`$ overlaps with particle $`1`$ and $`2`$ is still required. The final result is
$$Z_G=\frac{1+6\mathrm{ln}2}{4}NZ_0\varphi ^2ϵ^2G_2(\alpha _s).$$
(B18)
### 8 Graph H
The last graph H is the most complicated one. It represents the term $`f_{ij}^{(0)}f_{ik}^{(1)}f_{jk}^{(1)}`$ that appears $`N^3/2`$ times. The problem here is to fulfill the requirement that particles $`j`$ and $`k`$ have to overlap in terms of properly chosen integration limits. We start with performing the trivial integrations:
$`Z_H`$ $`=`$ $`{\displaystyle \frac{1}{2}}N^3z_0^{N3}V`$ (B21)
$`\times {\displaystyle }e^{v_1v_2v_3}v_{12}^{DD}v_{13}^{DD}O_{23}\overline{O}_{12}\overline{O}_{13}`$
$`\times d𝐫_{12}d𝐫_{13}d\mathrm{\Omega }_1d\mathrm{\Omega }_2d\mathrm{\Omega }_3.`$
Whether the integrand vanishes due to the hard core factors depends only on the distances $`r_{12}`$, $`r_{13}`$, and the angle $`\vartheta _{23}`$ between $`𝐫_{12}`$ and $`𝐫_{13}`$. Consider a special orientation where
$`\widehat{𝐫}_{12}^0`$ $`=`$ $`(1,0,0)`$ (B22)
$`\widehat{𝐫}_{13}^0`$ $`=`$ $`(\mathrm{cos}\vartheta _2,0,\mathrm{sin}\vartheta _{23}),`$ (B23)
with $`0\vartheta _{23}\pi `$. A general configuration of the particles’ locations can be written as
$$\widehat{𝐫}_{1(2,3)}=_z(\phi )_y(\vartheta )_z(\psi )\widehat{𝐫}_{1(2,3)}^0,$$
(B24)
where $`_x`$, $`_y`$, and $`_z`$ are Eulerian rotation matrices for the angles $`\psi `$, $`\vartheta `$, and $`\phi `$. Using this form the integration over the factors that depend on these angles:
$$_0^{2\pi }_{\pi /2}^{\pi /2}_0^{2\pi }v_{12}^{DD}v_{13}^{DD}\mathrm{cos}\vartheta d\phi d\vartheta d\psi ,$$
(B25)
can easily be performed with mathematica. We call the result $`I(r_{12},r_{13},\vartheta _{23},𝐦_i)`$.
Next, we integrate over the orientations of the $`𝐦_𝐢`$:
$$e^{v_1v_2v_3}I(r_{12},r_{13},\vartheta _{23},𝐦_i)𝑑\mathrm{\Omega }_1𝑑\mathrm{\Omega }_2𝑑\mathrm{\Omega }_3.$$
(B26)
The result depends only on $`r_{12}`$, $`r_{13}`$, and $`\vartheta _{23}`$. Using it in (B21) yields
$`Z_H`$ $`=`$ $`{\displaystyle \frac{4}{3}}N^3{\displaystyle \frac{Z_0}{V^2}}K(\alpha _s)`$ (B29)
$`\times {\displaystyle }O_{23}\overline{O}_{12}\overline{O}_{13}{\displaystyle \frac{m^4\pi ^2(2\mathrm{cos}^2\vartheta _{23}\mathrm{sin}^2\vartheta _{23})}{10(4\pi \mu _0kT)^2r_{12}r_{13}}}`$
$`\times \mathrm{sin}\vartheta _{23}dr_{12}dr_{13}d\vartheta _{23},`$
where
$`K(\alpha _s)`$ $`=`$ $`{\displaystyle \frac{3}{8}}\left({\displaystyle \frac{\alpha _s}{\mathrm{sinh}\alpha _s}}\right)^3{\displaystyle _1^1}{\displaystyle _1^1}{\displaystyle _1^1}e^{\alpha _s(u_1+u_2+u_3)}`$ (B31)
$`\times (u_1^2+3)u_2u_3du_1du_2du_3.`$
The explicit expression for $`K(\alpha _s)`$ is given in Appendix A (eqn. A21).
Now we discuss the hard core terms. $`r_{12}`$ and $`r_{13}`$ have to be greater than $`D`$ to avoid the overlap with particle $`1`$. Furthermore $`|r_{12}r_{13}|<D`$ has to be fulfilled for particle 2 and 3 to overlap. As a last requirement, $`\vartheta _{23}`$ has to be smaller than some angle $`\vartheta _{23}^{\text{max}}`$ that depends on $`r_{12}`$ and $`r_{13}`$. Trigonometry shows that
$$\mathrm{cos}\vartheta _{23}^{\text{max}}=\frac{r_{12}^2+r_{13}^2D^2}{2r_{12}r_{13}}.$$
(B32)
In this configuration the distance between particle 2 and 3 is exactly $`D`$.
We perform the integration over $`\vartheta _{23}`$ from $`0`$ to $`\vartheta _{23}^{\text{max}}`$ in (B29), choose the correct limits for $`r_{12}`$ and $`r_{13}`$, and drop all hard core terms:
$`Z_H`$ $`=`$ $`{\displaystyle \frac{4}{3}}N^3{\displaystyle \frac{Z_0}{V^2}}K(\alpha _s)`$ (B35)
$`\times {\displaystyle _D^{\mathrm{}}}{\displaystyle _{\mathrm{min}(D,r_{13}D)}^{r_{13}+D}}{\displaystyle \frac{m^4\pi ^2}{10(4\pi \mu _0kT)^2r_{12}r_{13}}}`$
$`\times \left[{\displaystyle \frac{r_{12}^2+r_{13}^2D^2}{2r_{12}r_{13}}}\left({\displaystyle \frac{r_{12}^2+r_{13}^2D^2}{2r_{12}r_{13}}}\right)^3\right]dr_{12}dr_{13},`$
The result of the last integration is:
$`Z_H`$ $`=`$ $`{\displaystyle \frac{4}{3}}N^3{\displaystyle \frac{Z_0}{V^2}}K(\alpha _s){\displaystyle \frac{m^4\pi ^2}{48(4\pi \mu _0kT)^2}}`$ (B36)
$`=`$ $`NZ_0\varphi ^2ϵ^2K(\alpha _s).`$ (B37)
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# Effects of arbitrarily directed field on spin phase oscillations in biaxial molecular magnets
## I Introduction
In recent years, owing mainly to the rapid advances both in new technologies of miniaturization and in highly sensitive SQUID magnetometry, there have been considerable theoretical and experimental studies carried out on the nanometer-scale magnets which exhibit macroscopic quantum phenomena(MQP). A number of nanometer-scale particles in the superparamagnetic regime have been identified as candidates for the observation of MQP such as the tunneling of the magnetization (or the Néel vector) out of metastable potential minimum through the classically impenetrable barrier to a stable one, i.e., macroscopic quantum tunneling (MQT), or, more strikingly, macroscopic quantum coherence (MQC), where the magnetization (or the Néel vector) coherently oscillates between energetically degenerate easy directions over many periods. Up to now, molecular magnets have been the most promising candidates to observe MQP because they have well-defined structures and magnetic properties. The system that have recently attracted much attention are the molecules $`\left[\text{Mn}_{12}\text{O}_{12}\text{(CH}_3\text{COO)}_{16}\text{(H}_2\text{O)}_4\right]`$ (in short Mn<sub>12</sub>Ac) and $`\left[\text{(tacn)}_6\text{Fe}_8\text{O}_2\text{(OH)}_{12}\right]^{8+}`$ (in short Fe<sub>8</sub>), where tacn is a macrocyclic ligand triazacyclononance. The Mn<sub>12</sub>Ac molecule contains four Mn<sup>4+</sup> ($`S=3/2`$) ions in a central tretrahedron surrounded by eight Mn<sup>3+</sup> ($`S=2`$) ions. Oxygen bridges allow superexchange coupling among the Mn ions, and both high-field and ac susceptometry experiments indicate a $`S=10`$ ground state, resulting from the four inner Mn<sup>4+</sup> spins being paralleled to each other and the other eight Mn<sup>3+</sup> also parallel with the two groups antiparallel to each other. The magnetization relaxation experiments, the dynamic susceptibility measurements, and the dynamic hysteresis experiments all indicate that the thermally assisted, magnetic-field-tuned resonant coherently quantum tunneling of magnetization occurs between spin states in a large number of identical Mn<sub>12</sub>Ac molecules. Magnetic measurements have shown that Fe<sub>8</sub> also has a spin ground state $`S=10`$, which arises from competing antiferromagnetic interactions between the eight $`S=5/2`$ iron spins with six spins being parallel and antiparallel to the other two spins.
One of the most striking effects in the magnetic MQP is that for some spin systems with high symmetries, the tunneling behaviors of magnetization seem sensitive to the parity of total spin of the magnet. It has been theoretically demonstrated that the ground-state tunneling level splitting is completely suppressed to zero for the half-integer total spin ferromagnets with biaxial crystal symmetry in the absence of an external magnetic field, resulting from the destructive interference of the Berry phase or the Wess-Zumino, Chern-Simons term in the Euclidean action between the symmetry-related tunneling paths connecting two classically degenerate minima. Such destructive interference effect for half-integer spins is known as the topological quenching. But for the integer spins, the quantum interference between topologically different tunneling paths is constructive, and therefore the ground-state tunneling level splitting is nonzero. While such spin-parity effects are sometimes related to Kramers degeneracy, they typically go beyond this theorem in rather unexpected ways. One recent experimental method based on the Landau-Zener model was developed by Wernsdorfer and Sessoli to measure the very small tunnel splitting on the order of $`10^8`$K in molecular Fe<sub>8</sub> magnets. They observed a clear oscillation of the tunnel splitting as a function of the magnetic field applied along the hard anisotropy axis, which is direct evidence of the role of the topological spin phase (Berry phase) in the spin dynamics of these molecules.
Motivated by the experiment on topological phase interference or spin-parity effects in the molecule Fe<sub>8</sub>, in this paper we investigate the resonant quantum tunneling of the magnetization vector in molecular magnets with biaxial crystal symmetry in the presence of an external magnetic field at an arbitrarily directed angle in the ZY plane. By applying the instanton technique in the spin-coherent-state path-integral representation, we calculate both the Wentzel-Kramers-Brillouin (WKB) exponent and the preexponential factor in the ground-state tunnel splitting. Our results show that for the small angle $`\theta _H`$ of the applied magnetic field, the ground-state tunnel splitting oscillates with the field for both the integer and half-integer spins, and the oscillation behavior for integer spins is significantly different from that for half-integer spins. However, this oscillation is completely suppressed for the large angle region. The distinct angular dependence, together with the dependence of the ground-state tunnel splitting on the strength of the external applied magnetic field, may provide an independent experimental test for spin-parity effects in molecular magnets. It is noted that the instanton technique is semiclassical in nature, i.e., valid for large spins and in the continuum limit. Whether the instanton technique can be applied in studying the spin dynamics in the molecular magnet with $`S=10`$ (such as Fe<sub>8</sub>) is an open question. We study this problem with the help of exact diagonalization calculation. Our results show that the analytical calculation based on the instanton technique agrees excellently well with the exact diagonalization calculation, which strongly suggests that the molecular magnet with $`S=10`$ can be treated as a giant spin system.
## II Model and Method
The system of interest is a molecular magnet at a temperature well below its anisotropy gap, which has the following Euclidean action in the spin-coherent-state representation,
$$S_E(\theta ,\varphi )=𝑑\tau \left[iS\left(\frac{d\varphi }{d\tau }\right)iS\left(\frac{d\varphi }{d\tau }\right)\mathrm{cos}\theta +E(\theta ,\varphi )\right],$$
(1)
where $`S`$ is the total spin of the molecular magnet. The polar angle $`\theta `$ and the azimuthal angle $`\varphi `$ label the spin coherent state. $`E(\theta ,\varphi )`$ is the total energy of the molecular magnet which includes the magnetocrystalline anisotropy energy and the Zeeman energy when an external magnetic field is applied.
It is noted that the Euclidean action is written in the north-pole gauge, and the first two terms in Eq. (1) define the Wess-Zumino or Berry term which arises from the nonorthogonality of spin coherent states. The Wess-Zumino term has a simple geometrical or topological interpretation. For a closed path, this term equals $`iS`$ times the area swept out on the unit sphere between the path and the north pole. The first term in Eq. (1) is a total imaginary-time derivative, which has no effect on the classical equation of motion for the magnetization vector, but yields the boundary contribution to the Euclidean action. Loss et al. and von Delft and Henley studied the physical effect of this total derivative term, and they found that this term is crucial for the quantum properties of the magnetic particle and makes the tunneling behaviors of integer and half-integer spins strikingly different.
In the semiclassical limit, the dominant contribution to the Euclidean transition amplitude comes from finite action solutions of the classical equations of motion (instantons), which can be expressed as the following equations in the spherical coordinate system,
$`iS\left({\displaystyle \frac{d\overline{\theta }}{d\tau }}\right)\mathrm{sin}\overline{\theta }`$ $`=`$ $`{\displaystyle \frac{E}{\overline{\varphi }}},`$ (2)
$`iS\left({\displaystyle \frac{d\overline{\varphi }}{d\tau }}\right)\mathrm{sin}\overline{\theta }`$ $`=`$ $`{\displaystyle \frac{E}{\overline{\theta }}},`$ (3)
where $`\overline{\theta }`$ and $`\overline{\varphi }`$ denote the classical path. Note that the Euclidean action Eq. (1) describes the $`\left(11\right)`$dimensional dynamics in the Hamiltonian formulation with canonical variables $`\varphi `$ and $`P_\varphi =S(1`$cos$`\theta )`$.
According to the instanton technique in the spin-coherent-state path-integral representation, the instanton’s contribution to the tunneling rate $`\mathrm{\Gamma }`$ for MQT or the ground-state tunnel splitting $`\mathrm{\Delta }`$ for MQC (not including the geometric phase factor generated by the topological term in the Euclidean action) is given by
$$\mathrm{\Gamma }(or\mathrm{\Delta })=p_0\omega _p\left(\frac{S_{cl}}{2\pi }\right)^{1/2}e^{S_{cl}},$$
(4)
where $`\omega _p`$ is the small-angle precession or oscillation frequency in the well, and $`S_{cl}`$ is the classical action or the WKB exponent which minimizes the Euclidean action of Eq. (1). The preexponential factor $`p_0`$ originates from the quantum fluctuations about the classical path, which can be evaluated by expanding the Euclidean action to second order in the small fluctuations.
We describe the molecular magnet with biaxial crystal symmetry by the standard Hamiltonian,
$$=k_1\widehat{S}_z^2+k_2\widehat{S}_y^2,$$
(5)
where $`k_1>k_2>0`$ are proportional to the anisotropy coefficients, and we take the easy, medium, and hard axes as x, y, and z, respectively. If the magnetic field is applied in ZY plane, at an arbitrary angle $`0\theta _H90^o`$ with z, the Hamiltonian becomes
$$=k_1\widehat{S}_z^2+k_2\widehat{S}_y^2g\mu _BH_z\widehat{S}_zg\mu _BH_y\widehat{S}_y,$$
(6)
where $`g`$ is the landé factor, and $`\mu _B`$ is the Bohr magneton. The Zeeman energy term associated with the applied field $`\stackrel{}{H}(0,H_y,H_z)`$ $`(0,H\mathrm{sin}\theta _H,H\mathrm{cos}\theta _H)`$ is given in the last two terms of the Hamiltonian. If the field is below some critical value $`H_c\left(\theta _H\right)`$ (to be computed), the Hamiltonian Eq. (6) has two degenerate minima, and therefore the magnetization can resonate between these two directions, providing a case of MQC. Under the spin-coherent-state and the imaginary time representation, the $`E(\theta ,\varphi )`$ term in the Euclidean action is given by
$`E(\theta ,\varphi )`$ $`=`$ $`k_1S^2\mathrm{cos}^2\theta +k_2S^2\mathrm{sin}^2\theta \mathrm{sin}^2\varphi g\mu _BSH_z\mathrm{cos}\theta g\mu _BSH_y\mathrm{sin}\theta \mathrm{sin}\varphi `$ (7)
$`=`$ $`K_1\mathrm{cos}^2\theta +K_2\mathrm{sin}^2\theta \mathrm{sin}^2\varphi 2K_1(H_z/H_a)\mathrm{cos}\theta 2K_1(H_y/H_a)\mathrm{sin}\theta \mathrm{sin}\varphi ,`$ (8)
with $`K_1=k_1S^2`$ and $`K_2=k_2S^2`$ being the transverse and longitudinal anisotropy coefficients respectively, and $`H_a=2K_1/g\mu _BS`$ being the anisotropy field. Introducing $`\lambda =K_2/K_1`$, $`\mathrm{cos}\theta _0=H\mathrm{cos}\theta _H/H_a`$, and $`\mathrm{sin}\theta _0\mathrm{sin}\varphi _0=H\mathrm{sin}\theta _H/\lambda H_a`$, the $`E(\theta ,\varphi )`$ reduces to
$$E(\theta ,\varphi )=K_1\left(\mathrm{cos}\theta \mathrm{cos}\theta _0\right)^2+K_2\left(\mathrm{sin}\theta \mathrm{sin}\varphi \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\right)^2+E_0,$$
(9)
where $`E_0`$ is a constant that makes $`E(\theta ,\varphi )`$ zero at the initial state. It is clearly shown in Eq. (9) that the energy minima of system are at $`\theta _1=\theta _0`$, $`\varphi _1=\varphi _0`$ and $`\theta _2=\theta _0`$, $`\varphi _2=\pi \varphi _0`$, and therefore there are two different instanton paths of opposite windings around hard anisotropy axis. We denote them by instanton $`A:\varphi =\varphi _0\varphi =+\pi /2\varphi =\pi \varphi _0`$, and instanton $`B:`$ $`\varphi =\varphi _0\varphi =\pi /2\varphi =\pi \varphi _0`$.
The critical field at which the energy barrier disappears can be determined by,
$`\mathrm{cos}\theta _0`$ $`=`$ $`{\displaystyle \frac{H\mathrm{cos}\theta _H}{H_a}}1,`$ (10)
$`\mathrm{sin}\varphi _0`$ $`=`$ $`{\displaystyle \frac{H\mathrm{sin}\theta _H}{\lambda H_a\left(1\frac{H^2\mathrm{cos}^2\theta _H}{H_a^2}\right)^{1/2}}}1.`$ (11)
Taking into account Eqs. (10) and (11), we obtain
$$H_c=\frac{\lambda H_a}{\left(\mathrm{sin}^2\theta _H+\lambda ^2\mathrm{cos}^2\theta _H\right)^{1/2}}.$$
(12)
For the special cases $`\theta _H=0,\pi /2`$, we have $`H_c=H_a,\lambda H_a`$, respectively. The dependence of $`H_c/H_a`$ on $`\theta _H`$ is plotted in Fig. 1 for $`\lambda =0.71`$.
Now we investigate the tunneling behaviors of magnetization by applying the instanton technique in the spin-coherent-state path-integral representation. First of all, we must find out the classical path $`\overline{\theta }`$ and $`\overline{\varphi }`$ that satisfies the boundary condition. Along the classical path, $`E(\overline{\theta },\overline{\varphi })`$ is conserved, so that the relation between $`\overline{\theta }(\tau )`$ and $`\overline{\varphi }(\tau )`$ can be found purely by using energy conservation. After some algebra, we obtain
$$\mathrm{cos}\overline{\theta }=\frac{\left(\mathrm{cos}\theta _0\lambda ^{1/2}y\mathrm{sin}\overline{\varphi }\right)\pm i\lambda ^{1/2}\left(x\mathrm{sin}\overline{\varphi }\mathrm{sin}\theta _0\mathrm{sin}\varphi _0\right)}{1\lambda \mathrm{sin}^2\overline{\varphi }},$$
(13)
where $`x`$, $`y`$ are variables defined by
$$x=\left\{\frac{(a^2+b^2)^{1/2}+a}{2}\right\}^{1/2}\text{and}y=\left\{\frac{(a^2+b^2)^{1/2}a}{2}\right\}^{1/2},$$
(14)
with $`a=1\mathrm{cos}^2\theta _0\lambda \mathrm{sin}^2\overline{\varphi }+\lambda \mathrm{sin}^2\theta _0\mathrm{sin}^2\varphi _0`$ and $`b=2\lambda ^{1/2}\mathrm{sin}\theta _0\mathrm{cos}\theta _0\mathrm{sin}\varphi _0`$. Here we assume that the condition $`a`$ $`0`$ is always fulfilled for the small magnetic field $`H`$. The positive and negative sign in Eq. (13) are corresponding to instanton and anti-instanton solutions for path $`A`$ or $`B`$, respectively. For the low magnetic field , one can take $`\overline{\varphi }(\tau )`$ to be entirely real (see appendix A for detail). Then through the expression $`S_{cl}=iS_{\mathrm{}}^+\mathrm{}(1\mathrm{cos}(\overline{\theta }))(d\overline{\varphi }/d\tau )𝑑\tau `$, the desired Wentzel-Kramers-Brillouin (WKB) exponents (or classical actions) are found to be
$`Im(S_{cl}^A)`$ $`=`$ $`S(+\pi 2\varphi _0){\displaystyle \frac{2S\mathrm{cos}\theta _0}{(1\lambda )^{1/2}}}\left({\displaystyle \frac{\pi }{2}}\mathrm{arctan}\left((1\lambda )^{1/2}\mathrm{tan}\varphi _0\right)\right)`$ (16)
$`+2\lambda ^{1/2}S{\displaystyle _{\varphi _0}^{\pi /2}}{\displaystyle \frac{y\mathrm{sin}\varphi d\varphi }{1\lambda \mathrm{sin}^2\varphi }},`$
$`Re(S_{cl}^A)`$ $`=`$ $`2\lambda ^{1/2}S\{{\displaystyle \frac{\mathrm{sin}\theta _0\mathrm{sin}\varphi _0}{(1\lambda )^{1/2}}}({\displaystyle \frac{\pi }{2}}+\mathrm{arctan}\left((1\lambda )^{1/2}\mathrm{tan}\varphi _0\right))`$ (19)
$`+{\displaystyle _{\varphi _0}^{\pi /2}}{\displaystyle \frac{x\mathrm{sin}\varphi d\varphi }{1\lambda \mathrm{sin}^2\varphi }}\}`$
for instanton $`A`$, and
$`Im(S_{cl}^B)`$ $`=`$ $`S(\pi 2\varphi _0)+{\displaystyle \frac{2S\mathrm{cos}\theta _0}{(1\lambda )^{1/2}}}\left({\displaystyle \frac{\pi }{2}}+\mathrm{arctan}\left((1\lambda )^{1/2}\mathrm{tan}\varphi _0\right)\right)`$ (21)
$`+2\lambda ^{1/2}S{\displaystyle _{\varphi _0}^{\pi /2}}{\displaystyle \frac{y\mathrm{sin}\varphi d\varphi }{1\lambda \mathrm{sin}^2\varphi }},`$
$`Re(S_{cl}^B)`$ $`=`$ $`2\lambda ^{1/2}S\{{\displaystyle \frac{\mathrm{sin}\theta _0\mathrm{sin}\varphi _0}{(1\lambda )^{1/2}}}(+{\displaystyle \frac{\pi }{2}}+\mathrm{arctan}\left((1\lambda )^{1/2}\mathrm{tan}\varphi _0\right))`$ (24)
$`+{\displaystyle _{\varphi _0}^{\pi /2}}{\displaystyle \frac{x\mathrm{sin}\varphi d\varphi }{1\lambda \mathrm{sin}^2\varphi }}\}`$
for instanton $`B`$. Therefore, the difference of the action of two tunneling paths is
$`Im(\mathrm{\Delta }S_{cl})`$ $`=`$ $`Im(S_{cl}^B)Im(S_{cl}^A)=2S\pi \left(1{\displaystyle \frac{\mathrm{cos}\theta _0}{(1\lambda )^{1/2}}}\right)2\mathrm{\Phi }(H),`$ (25)
$`Re(\mathrm{\Delta }S_{cl})`$ $`=`$ $`Re(S_{cl}^B)Re(S_{cl}^A)=2S\pi \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\left({\displaystyle \frac{\lambda }{1\lambda }}\right)^{1/2}.`$ (26)
Note that $`Im(\mathrm{\Delta }S_{cl})`$ comes from the Berry phase $`iS(1\mathrm{cos}\overline{\theta })\left(d\overline{\varphi }/d\tau \right)`$, and leads to the oscillation of the ground-state tunnel splitting with the magnetic field.
The preexponential factor in the tunnel splitting can be evaluated by taking the asymptotic form of zero-mode $`\left(d\overline{\varphi }/d\tau \right),`$ which can be deduced from the classical equations of motion (see appendix A).
Finally, the ground-state tunnel splitting is found to be
$`\mathrm{}\mathrm{\Delta }`$ $`=`$ $`\mathrm{}\mathrm{\Delta }_0\left|1+e^{2i\mathrm{\Phi }(H)}e^{Re\mathrm{\Delta }S_{cl}}\right|`$ (27)
$`=`$ $`\mathrm{}\mathrm{\Delta }_0\left\{\left(1e^{Re\mathrm{\Delta }S_{cl}}\right)^2+4e^{Re\mathrm{\Delta }S_{cl}}\mathrm{cos}^2\mathrm{\Phi }(H)\right\}^{1/2},`$ (28)
with
$$\mathrm{}\mathrm{\Delta }_0=c\omega _0^{3/2}\frac{4\mathrm{sin}\theta _0}{\left(\mathrm{sin}^2\theta _0\lambda \right)^{1/2}}\left(\frac{S^2}{2\pi K_1}\right)^{1/2}\left(\frac{\mathrm{sin}^2\theta _0}{\mathrm{sin}^2\theta _0+\lambda \mathrm{cos}^2\theta _0\mathrm{sin}^2\varphi _0}\right)^{1/2}e^{Re(S_{cl}^A)},$$
(29)
where $`\omega _0=(2K_1/S)\lambda ^{1/2}\mathrm{sin}\theta _0\mathrm{cos}\varphi _0`$. The dimensionless prefactor $`c`$ can often be of order 1 or so, and is therefore relevant to the exact diagonalization calculation.
## III Results and Discussion
Before we discuss the Eq. (28), we note here that our model can be directly related to the model describing the molecular Fe<sub>8</sub> magnet,
$$=D\widehat{S}_z^2+E(\widehat{S}_x^2\widehat{S}_y^2)+g\mu _BH_x^{}\widehat{S}_x+g\mu _BH_y^{}\widehat{S}_y,$$
(30)
with $`K_1=(D+E)S^2`$, $`K_2=(DE)S^2`$ and $`H_z=H_x^{}`$, $`H_y=H_y^{}`$. According to the typical parameters of Fe<sub>8</sub>, $`D=0.275`$K, $`E=0.046`$K and $`g=2`$, we obtain that $`K_1=32.1`$K, $`K_2=22.9`$K, $`\lambda =0.71`$ and $`H_a=4.77`$T. These parameters are used throughout the whole calculation.
First, we discuss the effects of arbitrarily directed field on the spin phase oscillation. Our results show that the topological phase interference or spin-parity effects depend on the direction of the magnetic field significantly. From Eq (28), whether the ground-state tunnel splitting oscillates with the field is determined by two factors, $`\mathrm{\Phi }(H)`$ and $`Re(\mathrm{\Delta }S_{cl})`$. For a fixed field strength, i.e. $`HH_S\left(\sqrt{H_y^2+H_z^2}\right)`$, it brings two-fold effects to increase the field angle $`\theta _H`$: (i) the dependence of $`\mathrm{\Phi }(H)`$ on the field strength is reduced. As $`\theta _H`$ increases up to $`\frac{\pi }{2},`$ $`\mathrm{\Phi }(H)`$ $`\pi S\left(1\frac{H\mathrm{cos}\theta _H}{(1\lambda )^{1/2}H_a}\right)`$ gradually reduces to a constant $`\pi S.`$ (ii) The degree of oscillation $`\sigma `$ which is defined as a ratio of the oscillation part $`4e^{Re\mathrm{\Delta }S_{cl}}`$ to the non-oscillation part $`\left(1e^{Re\mathrm{\Delta }S_{cl}}\right)^2`$ in Eq. (28) decreases from $`+\mathrm{}`$ to a small value.
For $`\theta _H=0`$ (the field is along the hard anisotropy axis), the symmetry of two classical paths imposes $`Re(\mathrm{\Delta }S_{cl})=0,`$ then, the tunnel splitting $`\mathrm{\Delta }=2\mathrm{\Delta }_0\left|\mathrm{cos}\mathrm{\Phi }(H)\right|`$ oscillates with the field and thus vanish whenever
$$\frac{H}{H_a}=(1\lambda )^{1/2}\frac{(Sn1/2)}{S},$$
(31)
where $`n=0,1,2,\mathrm{}`$ is an integer. For $`\theta _H=\frac{\pi }{2}`$ (the field is along the medium anisotropy axis), on the other hand, $`\mathrm{\Phi }(H)`$ becomes a constant $`\pi S`$, and therefore the tunnel splitting increases monotonically with the field. When the field is applied in ZY plane with an arbitrarily angle $`\theta _H`$, one may expect to observe a crossover for a certain field angle $`\theta _H^c`$. Choosing the magnitude of $`\stackrel{}{H}`$ to be the first value that makes cos$`\mathrm{\Phi }(H)=0`$ and taking $`\sigma =1`$, we obtain
$$\theta _H^c\mathrm{arctan}\left(\frac{1.76\lambda ^{1/2}}{\pi }\right)$$
(32)
For the molecular Fe<sub>8</sub> magnet, $`\lambda =0.71`$, we find $`\theta _H^c25^o`$, which agrees well with the exact diagonalization calculation (shown in Fig. 2a) and the experimental results obtained by Wernsdorfer and Sessoli. It is noted that the value of $`\theta _H^c`$ is independent of the total spin $`S`$ of the molecular magnet, and depends on the parameter $`\lambda `$ only. For highly anisotropic materials, the typical values of the transverse and longitudinal anisotropy coefficients are $`K_110^7`$erg/cm<sup>3</sup> and $`K_210^5`$erg/cm<sup>3</sup>. Thus, $`\lambda =0.01`$, the value of $`\theta _H^c`$ is estimated to be $`3.2^o`$ and is much smaller than that for the molecular Fe<sub>8</sub> magnet, which means that even a very small misalignment of the field with the hard anisotropy axis can completely destroy the oscillation. Therefore, the molecular Fe<sub>8</sub> magnet is a better candidate for observing the oscillation of the ground-state tunnel splitting with the field compared with the highly anisotropic materials.
Next, we turn to the comparison between the analytical calculations and the numerical stimulations. It is noted that the instanton approach is semiclassical in nature, i.e., valid for large spins and in the continuum limit. Whether the instanton technique can be applied in studying the spin dynamics in molecular magnets with $`S=10`$ (such as Fe<sub>8</sub>) is an open question. We have performed the numerical diagonalization of the Hamiltonian Eq. (6) for the molecular Fe$`_{8\text{ }}`$magnet in the presence of an external magnetic field at an arbitrarily directed angle in ZY plane. As illustrated in Fig. 2a, for Fe<sub>8</sub>, the analytical results based on the instanton technique are in good agreement with the exact diagonalization results. In order to show the agreement more clearly, we have presented the numerical and analytical results in Table 1 for the molecular Fe<sub>8</sub> magnet in the magnetic field applied at angle $`\theta _H=0^o,15^o,30^o`$ and $`90^o`$. It is clearly shown that the accuracy of the semiclassical calculation is very high for the low magnetic field up to $`H=1.0`$T for the entire region of the angle $`0^o\theta _H90^o`$. As a result, we conclude that the molecular magnet with $`S=10`$ can be treated as a giant spin system. The numerical and analytical calculated tunnel splitting as a function of the field are shown in Fig. 2b for half-integer spin $`S=9.5`$, and the good agreement between numerical and analytical results is also found. From Fig. 2, it is obviously that the tunneling behavior of magnetization of integer spins is significantly different from that for half-integer spins. At the end of this section, we present the results with different parameter $`\lambda `$ for the field along $`𝐳`$ or $`𝐲`$ in Figs. (3a) and (3b). It is interesting to note that the tunnel splitting increases rapidly by lowing the parameter $`\lambda `$, which suggests that highly anisotropic materials are more suitable for observing MQP in experiments.
In conclusion, we have studied the ground-state tunnel splitting in the molecular magnets with biaxial crystal symmetry in the presence of an external magnetic field at an arbitrarily directed angle. The switching from oscillation to the monotonic growth of the ground-state tunnel splitting on the field angle has been shown in detail. Our results are suitable for a quantitative description of some aspects of the new experimental behavior on the molecular Fe<sub>8</sub> magnets.
## Acknowledgments
The authors are indebted to Professor W. Wernsdorfer and Professor R. Sessoli for providing their paper (Ref. 9). The financial support from NSF-China (Grant No.19974019) and China’s ”973” program is gratefully acknowledged.
Appendix A: Evaluate the ground-state tunnel splitting
In this appendix, the general scheme for calculation of the ground-state tunnel splitting is presented, the main assumptions and approximations are also outlined.
We start with the relation between $`\theta (\tau )`$ and $`\varphi (\tau )`$ obtained from the energy conversion (see Eq. (13)). Only one instanton, say the + one, need be considered explicitly (from now on, we drop all the bars and identify $`\theta =\overline{\theta }(\tau )`$ , $`\varphi =\overline{\varphi }(\tau )`$ for convenient):
$`\mathrm{cos}\theta `$ $`=`$ $`{\displaystyle \frac{\left(\mathrm{cos}\theta _0\lambda ^{1/2}y\mathrm{sin}\varphi \right)+i\lambda ^{1/2}\left(x\mathrm{sin}\varphi \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\right)}{1\lambda \mathrm{sin}^2\varphi }},`$ (33)
$`\mathrm{sin}\theta `$ $`=`$ $`{\displaystyle \frac{\left(x\lambda \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\mathrm{sin}\varphi \right)i\left(\lambda ^{1/2}\mathrm{cos}\theta _0\mathrm{sin}\varphi y\right)}{1\lambda \mathrm{sin}^2\varphi }}.`$ (35)
In order to determine the classical paths, we need another equation of motion(see Eq. (3)),
$`iS\mathrm{sin}\theta \dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{E(\theta ,\varphi )}{\theta }}`$ (36)
$`=`$ $`2K_1(\mathrm{cos}\theta \mathrm{cos}\theta _0)\mathrm{sin}\theta 2K_2(\mathrm{sin}\theta \mathrm{sin}\varphi \mathrm{sin}\theta _0\mathrm{sin}\varphi _0)\mathrm{cos}\theta \mathrm{sin}\varphi .`$ (37)
After dividing the factor $`2K_1\mathrm{sin}\theta `$ on both sides and substituting Eqs. (33) and (35), we rewrite Eq. (37) as
$`i{\displaystyle \frac{S}{2K_1}}\dot{\varphi }`$ $`=`$ $`\lambda ^{1/2}y\mathrm{sin}\varphi +i\lambda ^{1/2}\left(x\mathrm{sin}\varphi \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\right)+`$ (39)
$`\lambda \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\mathrm{sin}\varphi {\displaystyle \frac{\left(\mathrm{cos}\theta _0\lambda ^{1/2}y\mathrm{sin}\varphi \right)+i\lambda ^{1/2}\left(x\mathrm{sin}\varphi \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\right)}{\left(x\lambda \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\mathrm{sin}\varphi \right)i\left(\lambda ^{1/2}\mathrm{cos}\theta _0\mathrm{sin}\varphi y\right)}}`$
$`=`$ $`Ref(\varphi )+iImf(\varphi ),`$ (40)
where
$`Ref(\varphi )`$ $`=`$ $`\lambda ^{1/2}y\mathrm{sin}\varphi +\lambda \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\mathrm{sin}\varphi `$ (42)
$`\times {\displaystyle \frac{\left(1\lambda \mathrm{sin}^2\varphi \right)\left(x\mathrm{cos}\theta _0\lambda ^{1/2}\mathrm{sin}\theta _0\mathrm{sin}\varphi _0\right)}{\left(x\lambda \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\mathrm{sin}\varphi \right)^2+\left(\lambda ^{1/2}\mathrm{cos}\theta _0\mathrm{sin}\varphi y\right)^2}},`$
$`Imf(\varphi )`$ $`=`$ $`\lambda ^{1/2}\left(x\mathrm{sin}\varphi \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\right)+\lambda \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\mathrm{sin}\varphi `$ (47)
$`\times {\displaystyle \frac{\begin{array}{c}[(\mathrm{cos}\theta _0\lambda ^{1/2}y\mathrm{sin}\varphi )(\lambda ^{1/2}\mathrm{cos}\theta _0\mathrm{sin}\varphi y)+\lambda ^{1/2}\times \\ (x\mathrm{sin}\varphi \mathrm{sin}\theta _0\mathrm{sin}\varphi _0)(x\lambda \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\mathrm{sin}\varphi )]\end{array}}{\left(x\lambda \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\mathrm{sin}\varphi \right)^2+\left(\lambda ^{1/2}\mathrm{cos}\theta _0\mathrm{sin}\varphi y\right)^2}}.`$
As shown in Eq. (40) , if $`Ref(\varphi `$ $`)`$ is always zero for arbitrary $`\varphi `$, one can then take $`\varphi (\tau )`$ to be entirely real. Unfortunately, it is not the case, $`Ref(\varphi `$ $`)`$ is not always zero except $`H_z=0`$ or $`H_y=0`$. Checking the right-hand-side of Eq. (42) more carefully, we find that up to the leading term, Eq. (42) can be rewritten as
$`\left|Ref(\varphi )\right|`$ $``$ $`{\displaystyle \frac{\lambda ^2\mathrm{sin}^3\varphi }{\left(1\lambda \mathrm{sin}^2\varphi \right)^{1/2}}}\mathrm{cos}\theta _0\mathrm{sin}\varphi _0`$ (48)
$``$ $`{\displaystyle \frac{\lambda ^2}{\left(1\lambda \right)^{1/2}}}\mathrm{cos}\theta _0\mathrm{sin}\varphi _0`$ (49)
$``$ $`{\displaystyle \frac{\lambda }{\left(1\lambda \right)^{1/2}}}{\displaystyle \frac{H_zH_y}{H_a^2}}`$ (50)
$`=`$ $`{\displaystyle \frac{\lambda }{2\left(1\lambda \right)^{1/2}}}{\displaystyle \frac{H^2\mathrm{sin}2\theta _H}{H_a^2}}.`$ (51)
It is clearly shown from Eq. (51) that the $`Ref(\varphi )`$ is almost two orders of magnitude smaller than $`Imf(\varphi )`$ for the low magnetic field(i.e., for $`H/H_a1/5`$). So it is sufficient to drop the term $`Ref(\varphi )`$ in Eq. (40) for the first order approximation and treat $`\varphi (\tau )`$ as a real variable. Under this assumption, the classical Euclidean action can be evaluated easily without resolving the analytical solutions of $`\theta (\tau )`$ and $`\varphi (\tau )`$,
$`S_{cl}`$ $`=`$ $`iS{\displaystyle _{\mathrm{}}^+\mathrm{}}(1\mathrm{cos}(\theta (\tau )))(d\varphi /d\tau )𝑑\tau `$ (52)
$`=`$ $`iS{\displaystyle _{\varphi _i}^{\varphi _f}}(1\mathrm{cos}(\theta (\varphi )))𝑑\varphi `$ (53)
$`=`$ $`iS{\displaystyle _{\varphi _i}^{\varphi _f}}\left(1{\displaystyle \frac{\left(\mathrm{cos}\theta _0\lambda ^{1/2}y\mathrm{sin}\varphi \right)+i\lambda ^{1/2}\left(x\mathrm{sin}\varphi \mathrm{sin}\theta _0\mathrm{sin}\varphi _0\right)}{1\lambda \mathrm{sin}^2\varphi }}\right)𝑑\varphi ,`$ (54)
where $`\varphi _i=\varphi (\tau \mathrm{})`$, $`\varphi _f=\varphi (\tau +\mathrm{})`$. The results for instantons $`A`$ and $`B`$ are given in Eqs. (19) and (24).
Next let us study the preexponential factor at two different magnetic field directions of $`\theta _H=0`$ and $`\pi /2\theta _H>0.`$
I. $`\theta _H=0`$
The preexponential factor can be evaluated from zero-mode $`\left(d\varphi /d\tau \right)`$. It is obvious from symmetry that the preexponential factors along two instanton path $`A`$ and $`B`$ are equal.
For $`\theta _H=0`$, we have $`\varphi _0=0`$ and $`\varphi _i=0,`$ $`\varphi _f=\pi `$. Then, from Eq. (14) we get $`x=\left(1\mathrm{cos}^2\theta _0\lambda \mathrm{sin}^2\varphi \right)^{1/2}`$ and $`y=0`$. Therefore, Eq. (40) reduces to
$$\dot{\varphi }=\frac{2K_1}{S}\lambda ^{1/2}\left(1\mathrm{cos}^2\theta _0\lambda \mathrm{sin}^2\varphi \right)^{1/2}\mathrm{sin}\varphi .$$
(55)
This equation can be integrated easily. Defining $`\omega _0=\frac{2K_1}{S}\lambda ^{1/2}\mathrm{sin}\theta _0`$, we obtain
$$\mathrm{cos}\varphi =\left(1\mathrm{cos}^2\theta _0\lambda \right)^{1/2}\frac{\mathrm{tanh}(\omega _0\tau )}{\left(1\mathrm{cos}^2\theta _0\lambda \mathrm{tanh}^2\left(\omega _0\tau \right)\right)^{1/2}}.$$
(56)
It is easily verified that $`\varphi 0,\pi `$, as $`\tau \pm \mathrm{}`$.
Following the standard procedure of the Ref. 10, we write the final result as
$$\mathrm{}\mathrm{\Delta }=c\omega _0^{3/2}\frac{8\mathrm{sin}\theta _0}{\left(1\mathrm{cos}^2\theta _0\lambda \right)^{1/2}}\left(\frac{S^2}{2\pi K_1}\right)^{1/2}e^{S_{cl}}\left|\mathrm{cos}\mathrm{\Phi }(H)\right|,$$
(57)
where
$$S_{cl}=2S\left\{\frac{1}{2}\mathrm{ln}\left(\frac{1+\frac{\lambda ^{1/2}}{\mathrm{sin}\theta _0}}{1\frac{\lambda ^{1/2}}{\mathrm{sin}\theta _0}}\right)\frac{\mathrm{cos}\theta _0}{2\left(1\lambda \right)^{1/2}}\mathrm{ln}\left(\frac{1+\frac{\lambda ^{1/2}\mathrm{cos}\theta _0}{\left(1\lambda \right)^{1/2}\mathrm{sin}\theta _0}}{1\frac{\lambda ^{1/2}\mathrm{cos}\theta _0}{\left(1\lambda \right)^{1/2}\mathrm{sin}\theta _0}}\right)\right\}$$
(58)
and
$$\mathrm{\Phi }(H)=\pi S\left(1\frac{\mathrm{cos}\theta _0}{(1\lambda )^{1/2}}\right).$$
(59)
Eq. (57) can be also deduced from Eqs. (19), (25), (26) and (28) for the special case $`\varphi _0=0.`$ It is interesting to note that Eqs. (58) and (59) agree exactly with the previous result (Eqs. (3.10) and (3.11) in Ref. 7) found by Garg for the molecular Fe<sub>8</sub> magnet.
II. $`\pi /2\theta _H>0`$
For $`\pi /2\theta _H>0`$, it is difficult to integrate the equation of motion Eq. (40). So we have to take the asymptotic form of zero-mode $`\left(d\varphi /d\tau \right)`$. Notice that $`\varphi \pi \varphi _0`$ as $`\tau +\mathrm{}`$, we can expand $`Imf(\varphi )`$ according to the small parameter $`\alpha =`$ $`\varphi (\pi \varphi _0)`$. Up to the leading term, Eq. (40) has the form
$$\dot{\varphi }=\omega _0\left(\varphi (\pi \varphi _0)\right),$$
(60)
where $`\omega _0=\frac{2K_1}{S}\lambda ^{1/2}\mathrm{sin}\theta _0\mathrm{cos}\varphi _0.`$ We then integrate the Eq. (60)
$$\varphi =\pi \varphi _0c_0e^{\omega _0\tau },$$
(61)
where $`c_0`$ is an integration constant. Therefore, the asymptotic form of zero-mode $`\left(d\varphi /d\tau \right)`$ can be written as
$$d\varphi /d\tau =c_0\omega _0e^{\omega _0\tau }.$$
(62)
Straightforward, following Ref. 10, the tunnel splitting $`\mathrm{}\mathrm{\Delta }^A`$ and $`\mathrm{}\mathrm{\Delta }^B`$ corresponding to path $`A`$ and $`B`$ have the form
$`\mathrm{}\mathrm{\Delta }^A`$ $`=`$ $`c_A\omega _0^{3/2}\left({\displaystyle \frac{S^2}{2\pi K_1}}\right)^{1/2}\left({\displaystyle \frac{\mathrm{sin}^2\theta _0}{\mathrm{sin}^2\theta _0+\lambda \mathrm{cos}^2\theta _0\mathrm{sin}^2\varphi _0}}\right)^{1/2}e^{Re(S_{cl}^A)}e^{iIm(S_{cl}^A)}`$ (63)
$`\mathrm{}\mathrm{\Delta }^B`$ $`=`$ $`c_B\omega _0^{3/2}\left({\displaystyle \frac{S^2}{2\pi K_1}}\right)^{1/2}\left({\displaystyle \frac{\mathrm{sin}^2\theta _0}{\mathrm{sin}^2\theta _0+\lambda \mathrm{cos}^2\theta _0\mathrm{sin}^2\varphi _0}}\right)^{1/2}e^{Re(S_{cl}^B)}e^{iIm(S_{cl}^B)},`$ (65)
where $`c_A`$ and $`c_B`$ are the numerical factors in order of $`O(1)`$ that come from the integration constant mentioned above ($`c_0`$ in Eq. (62)). Because the existence of $`H_y`$ breaks the symmetry between the classical path $`A`$ and $`B`$, $`c_A`$ may not be equal to $`c_B`$. But for the small magnetic field, we assume theirs equivalence. Comparing Eqs. (63) and (65) with Eq. (57), we find
$$c_A=c_B=c\frac{4\mathrm{sin}\theta _0}{\left(1\mathrm{cos}^2\theta _0\lambda \right)^{1/2}}.$$
(66)
Now we turn to derive the total tunnel splitting $`\mathrm{}\mathrm{\Delta }`$ using a recently proposed effective Hamiltonian approach. For the present case, the effective Hamiltonian is found to be
$$_{eff}=\left[\begin{array}{c}0\\ \left(\mathrm{}\mathrm{\Delta }^A+\mathrm{}\mathrm{\Delta }^B\right)^{}\end{array}\begin{array}{c}\left(\mathrm{}\mathrm{\Delta }^A+\mathrm{}\mathrm{\Delta }^B\right)\\ 0\end{array}\right].$$
(67)
Diagonalizing the effective Hamiltonian, we obtain the desired result as shown in Eqs. (28) and (29).
| | | $`\theta _H=0^o`$ $`\mathrm{\Delta }_d\mathrm{\Delta }_i`$ | | $`\theta _H=15^o`$ $`\mathrm{\Delta }_d\mathrm{\Delta }_i`$ | | $`\theta _H=30^o`$ $`\mathrm{\Delta }_d\mathrm{\Delta }_i`$ | | $`\theta _H=90^o`$ $`\mathrm{\Delta }_d\mathrm{\Delta }_i`$ |
| --- | --- | --- | --- | --- | --- | --- | --- | --- |
| $`H_S=0.0`$T | | 0.00683 0.00683 | | 0.00683 0.00683 | | 0.00683 0.00683 | | 0.00683 0.00683 |
| $`H_S=0.2`$T | | 0.00556 0.00558 | | 0.00778 0.00780 | | 0.01445 0.01448 | | 0.05651 0.05658 |
| $`H_S=0.4`$T | | 0.00165 0.00167 | | 0.01714 0.01733 | | 0.06521 0.06591 | | 0.75350 0.75752 |
| $`H_S=0.6`$T | | 0.00502 0.00515 | | 0.04658 0.04782 | | 0.29419 0.30232 | | 8.20514 8.30553 |
| $`H_S=0.8`$T | | 0.01417 0.01485 | | 0.13101 0.13783 | | 1.34388 1.42229 | | 73.1429 74.7913 |
| $`H_S=1.0`$T | | 0.02362 0.02548 | | 0.39753 0.43399 | | 6.19230 6.88449 | | 537.366 557.203 |
| $`H_S=1.2`$T | | 0.02505 0.02808 | | 1.31263 1.51459 | | 28.5282 34.3140 | | 3277.23 3462.35 |
| $`H_S=1.4`$T | | 0.00799 0.00942 | | 4.70886 5.90363 | | 129.907 176.358 | | 16704.7 18092.3 |
Table 1. The ground-state tunnel splitting $`\mathrm{\Delta }_d`$ (calculated by diagonalization) and $`\mathrm{\Delta }_i`$ (calculated by instanton approach) with $`\theta _H=0^o`$, $`15^o`$, $`30^o`$ and $`90^o`$ are listed as a function of magnetic field $`H_S`$. The units for the tunnel splitting and magnetic field are $`10^7`$ Kelvin and Tesla, respectively. Here, $`S=10`$, $`K_1=32.1`$K, $`K_2=22.9`$K and $`\lambda =0.71`$ are used for the molecular Fe<sub>8</sub> magnet.
Figures Captions
Fig. 1. The critical magnetic field $`H_c`$ is plotted as a function of anlge $`\theta _H`$ at $`\lambda =0.71`$.
Fig. 2. The ground-state tunnel splitting $`\mathrm{\Delta }`$ are plotted as a function of the magnetic field $`H_S`$ at $`\theta _H=0^o,`$ $`3^o,`$ $`10^o,`$ $`30^o,`$ $`60^o,`$ and $`90^o`$ for (a) $`S=10`$ and (b) $`S=9.5`$. The other parameters are $`K_1=32.1`$K, $`K_2=22.9`$K and $`\lambda =0.71,`$ the typical values of the molecular Fe<sub>8</sub> magnet. The results of the instanton approach and the exact diagonalization are represented by the solid lines and square symbols, respectively.
Fig. 3. The ground-state tunnel splitting $`\mathrm{\Delta }`$ are plotted as a function of the magnetic field $`H_S`$ for $`\lambda =0.3,`$ $`0.5`$ and $`0.7`$ with $`S=10`$, $`K_1=32.1`$K at (a) $`\theta _H=0^o`$ and (b) $`\theta _H=90^o`$. The results of the instanton approach and the exact diagonalization are represented by the solid lines and square symbols, respectively.
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# 1 Introduction.
## 1 Introduction.
Physical theories formulated in different-than-usual spacetimes signatures have recently found increased attention. One of the reasons can be traced to the conjectured $`F`$-theory which supposedly lives in $`(2+10)`$ dimensions . The current interest in AdS theories motivated by the AdS/CFT correspondence furnishes another motivation. Two-time physics e.g. has started been explored by Bars and collaborators in a series of papers . From another point of view we can also recall that a fundamental theory is expected to explain not only the spacetime dimensionality, but even its signature (see ). Quite recently Hull and Hull-Khuri pointed out the existence of dualities relating different compactifications of theories formulated in different signatures. Such a result provides new insights to the whole question of spacetime signatures. In another context (see e.g. ) the existence of space-time dualities has also been remarked.
Majorana-Weyl spacetimes (i.e. those supporting Majorana-Weyl spinors) are at the very core of the present knowledge of the unification via supersymmetry, being at the basis of ten-dimensional superstrings, superYang-Mills and supergravity theories (and perhaps the already mentioned $`F`$-theory). A well-established feature of Majorana-Weyl spacetimes is that they are endorsed of a rich structure. A legitimate question that could be addressed is whether they are affected, and how, by space-time dualities. The answer is positive. Indeed all different Majorana-Weyl spacetimes which are possibly present in any given dimension are each-other related by duality transformations which are induced by the $`Spin(8)`$ triality automorphisms. The action of the triality automorphisms is quite non-trivial and has far richer consequences than the $`𝐙_2`$-duality (its most trivial representative) associated to the space-time $`(s,t)(t,s)`$ exchange discussed in . It corresponds to $`S_3`$, the six-element group of permutations of three letters, identified with the group of congruences of the triangle and generated by two reflections. The lowest dimension in which the triality action is non-trivial is $`8`$ (not quite a coincidence), where the spacetimes $`(8+0)(4+4)(0+8)`$ are all interrelated. They correspond to the transverse coordinates of the $`(9+1)(5+5)(1+9)`$ spacetimes respectively, where the triality action can also be lifted. Triality relates as well the $`12`$-dimensional Majorana-Weyl spacetimes $`(10+2)(6+6)(2+10)`$, i.e. the potentially interesting cases for the $`F`$-theory, and so on. Triality allows explaining the presence of points (read theory) in the brane-molecule table of ref. , corresponding to the different versions of e.g. superstrings, $`11`$-dimensional supermembranes, fivebranes.
As a consequence of triality, supersymmetric theories formulated with Majorana-Weyl spinors in a given dimension but with different signatures, are all dually mapped one into another. A three-language dictionary is here furnished with the exact translations among the different versions of such supersymmetric theories. It should be stressed that, unlike , the dualities here discussed are already present for the uncompactified theories and in this respect look more fundamental.
The reason why the triality of the $`d=8`$ dimension plays a role in dually relating higher dimensional spacetimes is a consequence of the fact that higher-dimensional Lorentz groups admit the $`8`$-dimensional Lorentz as a subgroup. This feature is neatly encoded in the representation properties of the higher dimensional Clifford $`\mathrm{\Gamma }`$-matrices. In fact, due to this argument, it can be shown that not only different-signature Majorana-Weyl spacetimes are duality-related, but also that different-signatures odd-dimensional Majorana spacetimes are all interconnected via triality. This is true in particular for the $`11`$-dimensional case relevant for the maximal supergravity and the $`M`$-theory.
Manifestations of triality are observed in different contexts. We will stress the fact that besides the original “Cartan” triality exhibited by vectors, chiral and antichiral spinors (and therefore also denoted as “VCA”-triality in the following) in each one of the three Majorana-Weyl spacetimes of signature $`(8+0)`$, $`(4+4)`$, $`(0+8)`$ respectively, consequences of triality are found at the Clifford’s $`\mathrm{\Gamma }`$ matrices representation level. In $`d=8`$ dimension this is seen by the fact that “Majorana type representations” for $`\mathrm{\Gamma }`$-matrices, i.e. such that all the $`\mathrm{\Gamma }`$’s have a definite (anti-)symmetry property, only exist for the $`4_S+4_A`$, $`8_S+0_A`$, $`0_S+8_A`$ cases.
Moreover it can be shown that the three Majorana-Weyl spacetimes of signatures $`(4+4)`$, $`(8+0)`$, $`(0+8)`$ are interrelated via the $`S_3`$ permutation group. We call this property the “signature-triality” or the “space-time triality”.
It is worth stressing the fact that the arising of the $`S_3`$ permutation group as a signature-duality group for Majorana-Weyl spacetimes in a given dimension is not a completely straightforward consequence of the existence of Majorana-Weyl spacetimes in three different signatures. Some extra-requirements have to be fulfilled in order this to be true. As an example we just mention that a necessary condition for the presence of $`S_3`$ requires that each given couple of the three different spacetimes must differ by an even number of signatures (in the text this point will be discussed in full detail); the flipping of an odd number of signatures, like in the Wick rotation from Minkowski to the Euclidean space, cannot be achieved with a $`𝐙_2`$ group when spinors are involved.
In the present paper various aspects of the triality property will be rather extensively discussed. This is due to the objective relevance of the transverse $`8`$-dimensional space of coordinates in the light-cone formulations of superstrings theories, branes, and the non-perturbative (M)atrix approach to $`M`$-theory. The eight dimensions are indeed the natural setting for the appearance of triality-related symmetries. Triality is of course a very well-known property which has been extensively investigated in the literature both for technical reasons (as an example it implies the equivalence of the NSR and GS formulation of superstrings, see e.g. ) as well in its more fundamental aspects. In this paper, besides discussing the signature triality, we express its action on the Majorana-type representations of the Clifford’s $`\mathrm{\Gamma }`$ matrices. Moreover, for completeness, we review and extend some of the original Cartan’s results on V-C-A triality.
He gave a concrete representation within a non-diagonal metric which, once diagonalized, shows the $`(4+4)`$-signature. We point out that the triality generators he produced do not close the $`S_3`$ permutation group. However, an alternative presentation is available which implies the existence of $`S_3`$. As a consequence the triality group $`𝒢_{Tr}`$ defined in section $`6`$ is given by the semidirect product of a linear subgroup $`𝒢`$ (investigated also in ) of $`24`$-dimensional matrices with such an $`S_3`$ permutation group. It is worth mentioning that in the formulas of paragraph 139 of Cartan’s book an obvious typographical mistake plus a sign error appear.
In the appendices $`2`$-$`4`$ we produce the triality generators closing the $`S_3`$ group for each one of the three signatures $`(4+4)`$, $`(8+0)`$, $`(0+8)`$ which carry the triality structure. All formulas are here expressed in a Majorana-Weyl basis.
Furthermore, the construction of bilinear and trilinear invariants under the $`S_3`$ permutation group of the three Majorana-Weyl spacetimes is performed. They can be possibly used to formulate supersymmetric Majorana-Weyl theories in a manifestly triality-invariant form which presents an explicit symmetry under exchange of space and time coordinates.
The scheme of this work is as follows. In the next section we recall, following and , the basic properties of $`\mathrm{\Gamma }`$-matrices and Majorana conditions needed for our construction. Majorana-type representations are analyzed in section $`3`$. We show there how to relate the Majorana-Weyl representations in $`d>8`$ to the $`8`$-dimensional Majorana-Weyl representations. In section $`4`$ we introduce, for $`d=8`$, the set of data necessary to define a supersymmetric Majorana-Weyl theory, i.e. the set of “words” of our three-languages dictionary. The Cartan’s triality among vectors, chiral and antichiral spinors is presented in section $`5`$. The main result is furnished in section $`6`$, where spacetime triality is discussed. In section $`7`$ the triality as a generator of a group of invariances is analyzed. In section $`8`$ an application of triality is made. It is shown how to connect via triality some points (theories) presented in the Blencowe-Duff “brane-molecule scan” of reference (further enlarged in the second paper referred in ). In the original paper only the presence of “mirror” theories connected by a $`𝐙_2`$ space-versus-time exchange was “explained”. In the Conclusions we furnish some comments and point out some perspectives of future works. In the appendices, besides the already mentioned results, some useful construction concerning $`\mathrm{\Gamma }`$ matrices representations in $`6`$ and $`8`$ dimensions is given.
## 2 Preliminary results.
In this section the basic ingredients needed for our construction and the conventions employed will be introduced. More detailed information can be found in and .
We denote as $`g_{mn}`$ the flat (pseudo-)euclidean metric of a $`(t+s)`$-spacetime. Time (space) directions in our conventions are associated to the $`+`$ (respectively $``$) sign.
The $`\mathrm{\Gamma }`$’s matrices are assumed to be unitary and without loss of generality a time-like $`\mathrm{\Gamma }`$ is normalized so that its square is $`+1`$. The three matrices $`𝒜`$, $``$, $`𝒞`$ are the generators of the three conjugation operations (hermitian, complex conjugation and transposition respectively) on the $`\mathrm{\Gamma }`$’s. In particular
$`𝒞\mathrm{\Gamma }^m𝒞^{}`$ $`=`$ $`\eta (1)^{t+1}\mathrm{\Gamma }_{}^{m}{}_{}{}^{T}`$ (1)
where $`\eta =\pm 1`$ is a sign. In even-dimensional spacetimes it labels the two inequivalent choices of the charge conjugation matrix $`𝒞`$.
A relation exists, given by the formula
$`𝒞`$ $`=`$ $`^T𝒜`$ (2)
expressing anyone of the three matrices $`𝒜,,𝒞`$ in terms of the two others.
Up to an inessential phase, $`𝒜`$ is specified by the product of all the time-like $`\mathrm{\Gamma }`$ matrices.
An unitary transformations $`U`$ applied on spinors act on $`\mathrm{\Gamma }^m`$, $`𝒜,,𝒞`$ according to
$`\mathrm{\Gamma }^m`$ $``$ $`U\mathrm{\Gamma }^mU^{}`$
$`𝒜`$ $``$ $`U𝒜U^{}`$
$``$ $``$ $`U^{}U^{}`$
$`𝒞`$ $``$ $`U^{}𝒞U^{}`$ (3)
A Majorana representation for the $`\mathrm{\Gamma }`$’s can be defined as the one in which $``$ is set equal to the identity. Spinors can be assumed real in this case.
In even dimensions we can also introduce the notion of Weyl representation, i.e. when the “generalized $`\mathrm{\Gamma }^5`$ matrix” is symmetric and block diagonal and with no loss of generality can be assumed to be the direct sum of the two equal-size blocks $`\mathrm{𝟏}(\mathrm{𝟏})`$. The compatibility of both Majorana and Weyl conditions constraints the spacetime $`(t+s)`$ to satisfy
$`st`$ $`=`$ $`0\mathrm{𝑚𝑜𝑑}8,\mathrm{𝑓𝑜𝑟}\mathrm{𝑏𝑜𝑡ℎ}\mathrm{𝑣𝑎𝑙𝑢𝑒𝑠}\eta =\pm 1`$ (4)
In even dimensions Majorana representations, but not of Weyl type, are also found for
$`st`$ $`=`$ $`2mod8for\eta =1;`$
$`st`$ $`=`$ $`6mod8for\eta =+1.`$ (5)
In odd dimensions Weyl spinors cannot be defined, while Majorana spinors exist for
$`st`$ $`=`$ $`1,7mod8`$ (6)
For $`d<8`$ the only spacetimes supporting Majorana-Weyl spinors have signatures $`(n+n)`$. At $`d=8`$ a new feature arises, Majorana-Weyl spinors can be found for different signatures.
Making explicit the relation between theories formulated in such different signatures is the main content of this paper.
## 3 Majorana-type representations.
It is convenient to introduce the notion of Majorana-type representation (or shortly MTR) of the Clifford’ s $`\mathrm{\Gamma }`$ matrices. It can be defined as a representation such that all the $`\mathrm{\Gamma }`$’s have a definite symmetry, being either symmetric or antisymmetric. In $`d`$ dimensions a MTR with $`p`$ symmetric and $`q`$ antisymmetric $`\mathrm{\Gamma }`$’s ($`p+q=d`$) will be denoted as $`(p_S+q_A)`$ in the following.
When specialized to such representations the $`𝒞`$ charge-conjugation matrix introduced in the previous section is given by either the product of all the symmetric $`\mathrm{\Gamma }`$ matrices, denoted as $`𝒞_S`$, or all the antisymmetric ones ($`𝒞_A`$)
$`𝒞_S`$ $`=`$ $`\mathrm{\Pi }_{i=1,\mathrm{},p}\mathrm{\Gamma }_{}^{i}{}_{S}{}^{}`$
$`𝒞_A`$ $`=`$ $`\mathrm{\Pi }_{i=1,\mathrm{},q}\mathrm{\Gamma }_{}^{i}{}_{A}{}^{}`$ (7)
Please notice that the index $`S,A`$ is not referred to the (anti-)symmetry property of the matrices $`C_{S,A}`$ themselves.
In even dimensions $`𝒞_S`$, $`𝒞_A`$ correspond to opposite values of $`\eta `$ in (1), while in odd dimensions, up to an inessential phase factor, the two definitions for the charge-conjugation matrix collapse into a single matrix. This is in agreement with the property that in odd dimensions, up to unitary conjugation, the $`𝒞`$-matrix is uniquely defined. The convenience of using MTR’s to discuss Wick rotations to and from the Euclidean space has been advocated in .
It can be easily recognized that a Majorana representation for Clifford’s $`\mathrm{\Gamma }`$ matrices in a given signature spacetime implies the $`\mathrm{\Gamma }`$’s belong to a MTR. Conversely, given a MTR with $`p_S,q_A`$ (anti-)symmetric $`\mathrm{\Gamma }`$ matrices, two spacetimes exist ($`t=p`$, $`s=q`$ with the choice $`𝒞𝒞_S`$, and respectively $`t=q`$, $`s=p`$ for $`𝒞𝒞_A`$) such that the representation is Majorana (i.e. $`=\mathrm{𝟏}`$). The admissible couples of $`(p_S,q_A)`$ values for a MTR can be immediately read from the Majorana tables given above (4), (5) and (6). The construction is such that $`𝒞`$ must correspond to the correct value of $`\eta `$ appearing in the tables.
The list of all possible MTR’s in any given dimension is therefore easily computed. In order to furnish an example we mention that in $`d=6`$ there exists a MTR (not of Weyl kind) with $`6`$ anticommuting $`\mathrm{\Gamma }`$ matrices plus an anticommuting $`\mathrm{\Gamma }^7`$ $`(0_S+6_A,\mathrm{\Gamma }_{}^{7}{}_{A}{}^{})`$. It provides a Majorana basis for an Euclidean $`6`$-dimensional space. A concrete realization of such a representation is presented in appendix $`1`$<sup>1</sup><sup>1</sup>1It is worth mentioning that such a representation can be written in terms of the octonionic structure constants (see ref. ). In the present work our results have been obtained making no explicit reference to octonions; investigating the connection with such a division algebra is outside the scope of this work..
For completeness we present a list of all MTR’s up to $`d=18`$. We obtain
$`\begin{array}{cccc}𝐝& 𝐖𝐑& & \mathrm{𝐍𝐖}𝐑\\ & & \multicolumn{-1}{c}{}& \\ 2+1& (1_S+1_A+\mathrm{\Gamma }_{}^{3}{}_{S}{}^{})& & (2_S+0_A+\mathrm{\Gamma }_{}^{3}{}_{A}{}^{})\\ & & \multicolumn{-1}{c}{}& \\ 4+1& (2_S+2_A+\mathrm{\Gamma }_{}^{5}{}_{S}{}^{})& & (3_S+1_A+\mathrm{\Gamma }_{}^{5}{}_{A}{}^{})\\ & & \multicolumn{-1}{c}{}& \\ 6+1& (3_S+3_A+\mathrm{\Gamma }_{}^{7}{}_{S}{}^{})& & (4_S+2_A+\mathrm{\Gamma }_{}^{7}{}_{A}{}^{})\\ 6+1& & & (0_S+6_A+\mathrm{\Gamma }_{}^{7}{}_{A}{}^{})\\ & & \multicolumn{-1}{c}{}& \\ 8+1& (8_S+0_A+\mathrm{\Gamma }_{}^{9}{}_{S}{}^{})& & \\ 8+1& (4_S+4_A+\mathrm{\Gamma }_{}^{9}{}_{S}{}^{})& & (5_S+3_A+\mathrm{\Gamma }_{}^{9}{}_{A}{}^{})\\ 8+1& (0_S+8_A+\mathrm{\Gamma }_{}^{9}{}_{S}{}^{})& & (1_S+7_A+\mathrm{\Gamma }_{}^{9}{}_{A}{}^{})\\ & & \multicolumn{-1}{c}{}& \\ 10+1& (9_S+1_A+\mathrm{\Gamma }_{}^{11}{}_{S}{}^{})& & (10_S+0_A+\mathrm{\Gamma }_{}^{11}{}_{A}{}^{})\\ 10+1& (5_S+5_A+\mathrm{\Gamma }_{}^{11}{}_{S}{}^{})& & (6_S+4_A+\mathrm{\Gamma }_{}^{11}{}_{A}{}^{})\\ 10+1& (1_S+9_A+\mathrm{\Gamma }_{}^{11}{}_{S}{}^{})& & (2_S+8_A+\mathrm{\Gamma }_{}^{11}{}_{A}{}^{})\\ & & \multicolumn{-1}{c}{}& \\ 12+1& (10_S+2_A+\mathrm{\Gamma }_{}^{13}{}_{S}{}^{})& & (11_S+1_A+\mathrm{\Gamma }_{}^{13}{}_{A}{}^{})\\ 12+1& (6_S+6_A+\mathrm{\Gamma }_{}^{13}{}_{S}{}^{})& & (7_S+5_A+\mathrm{\Gamma }_{}^{13}{}_{A}{}^{})\\ 12+1& (2_S+10_A+\mathrm{\Gamma }_{}^{13}{}_{S}{}^{})& & (3_S+9_A+\mathrm{\Gamma }_{}^{13}{}_{A}{}^{})\\ & & \multicolumn{-1}{c}{}& \\ 14+1& (11_S+3_A+\mathrm{\Gamma }_{}^{15}{}_{S}{}^{})& & (12_S+2_A+\mathrm{\Gamma }_{}^{15}{}_{A}{}^{})\\ 14+1& (7_S+7_A+\mathrm{\Gamma }_{}^{15}{}_{S}{}^{})& & (8_S+6_A+\mathrm{\Gamma }_{}^{15}{}_{A}{}^{})\\ 14+1& (3_S+11_A+\mathrm{\Gamma }_{}^{15}{}_{S}{}^{})& & (4_S+10_A+\mathrm{\Gamma }_{}^{15}{}_{A}{}^{})\\ 14+1& & & (0_S+14_A+\mathrm{\Gamma }_{}^{15}{}_{A}{}^{})\\ & & \multicolumn{-1}{c}{}& \\ 16+1& (16_S+0_A+\mathrm{\Gamma }_{}^{17}{}_{S}{}^{})& & \\ 16+1& (12_S+4_A+\mathrm{\Gamma }_{}^{17}{}_{S}{}^{})& & (13_S+3_A+\mathrm{\Gamma }_{}^{17}{}_{A}{}^{})\\ 16+1& (8_S+8_A+\mathrm{\Gamma }_{}^{17}{}_{S}{}^{})& & (9_S+7_A+\mathrm{\Gamma }_{}^{17}{}_{A}{}^{})\\ 16+1& (4_S+12_A+\mathrm{\Gamma }_{}^{17}{}_{S}{}^{})& & (5_S+11_A+\mathrm{\Gamma }_{}^{17}{}_{A}{}^{})\\ 16+1& (0_S+16_A+\mathrm{\Gamma }_{}^{17}{}_{S}{}^{})& & (1_S+15_A+\mathrm{\Gamma }_{}^{17}{}_{A}{}^{})\\ & & \multicolumn{-1}{c}{}& \\ 18+1& (17_S+1_A+\mathrm{\Gamma }_{}^{19}{}_{S}{}^{})& & (18_S+0_A+\mathrm{\Gamma }_{}^{19}{}_{A}{}^{})\\ 18+1& (13_S+5_A+\mathrm{\Gamma }_{}^{19}{}_{S}{}^{})& & (14_S+4_A+\mathrm{\Gamma }_{}^{19}{}_{A}{}^{})\\ 18+1& (9_S+9_A+\mathrm{\Gamma }_{}^{19}{}_{S}{}^{})& & (10_S+8_A+\mathrm{\Gamma }_{}^{19}{}_{A}{}^{})\\ 18+1& (5_S+13_A+\mathrm{\Gamma }_{}^{19}{}_{S}{}^{})& & (6_S+12_A+\mathrm{\Gamma }_{}^{19}{}_{A}{}^{})\\ 18+1& (1_S+17_A+\mathrm{\Gamma }_{}^{19}{}_{S}{}^{})& & (2_S+16_A+\mathrm{\Gamma }_{}^{19}{}_{A}{}^{})\\ & & \multicolumn{-1}{c}{}& \end{array}`$ (37)
Some comments are in order. In the second column we listed all MTR’s of Weyl type, in the last one the non-Weyl representations. An arrow connects two given Weyl and non-Weyl representations which are intertwined by an exchange of the “generalized $`\mathrm{\Gamma }^5`$” matrix with any other $`\mathrm{\Gamma }`$-matrix of opposite symmetry. Up to $`d=18`$ the only true genuine even-dimensional non-Weyl Majorana representations not having such a Weyl counterpart are the above-mentioned $`6`$-dimensional $`(0_S+6_A)`$ representation and the $`14`$-dimensional $`(0_S+14_A)`$, both appearing in the table. Their difference in dimensionality $`(=8)`$ is of course a consequence of the famous $`mod8`$ property of the $`\mathrm{\Gamma }`$ matrices, see e.g. .
Different MTR’s belong to different classes under similarity transformations of the $`\mathrm{\Gamma }`$’s representations. Indeed in, let’s say, an euclidean (all $`+`$ signs) space, the index
$`I=tr(\mathrm{\Gamma }^m\mathrm{\Gamma }_{m}^{}{}_{}{}^{T})=(p_Sq_A)tr\mathrm{𝟏}`$ (38)
takes different values for different MTR’s and is by construction invariant under the transformation $`\mathrm{\Gamma }O\mathrm{\Gamma }O^T`$ realized by orthogonal matrices $`O`$.
Up to $`d=8`$ (excluded) there exists a unique similarity class of MTR’s of Weyl type, so that Majorana-Weyl spinors can be defined in $`(n+n)`$ space-times only. A new feature arises starting from $`d8`$, Weyl and Majorana-type representations are compatible for different similarity classes. In $`d=8`$ the three solutions $`(8_S+0_A)`$, $`(4_S+4_A)`$, $`(0_S+8_A)`$, associated with the corresponding Majorana-Weyl spacetimes, are found.
An efficient tool to produce and analyze Weyl representations in any higher dimensions is furnished by the following algorithm and . It provides a recursive procedure to construct $`d`$-dimensional Weyl representations in terms of any given couple of $`r`$, $`s`$ lower-dimensional $`\mathrm{\Gamma }`$-matrices representations, where the even integers $`d,r,s`$ are constrained to satisfy
$`d`$ $`=`$ $`r+s+2`$ (39)
Moreover, if the $`r,s`$-dimensional representations are of Majorana-type, the $`d`$-dimensional one is Majorana-Weyl.
The algorithm can be expressed through the formula
$`\mathrm{\Gamma }_{d}^{}{}_{}{}^{i=1,\mathrm{},s+1}`$ $`=`$ $`\sigma _x\mathrm{𝟏}_L\mathrm{\Gamma }_{s}^{}{}_{}{}^{i=1,\mathrm{},s+1}`$
$`\mathrm{\Gamma }_{d}^{}{}_{}{}^{s+1+j=s+2,\mathrm{},d}`$ $`=`$ $`\sigma _y\mathrm{\Gamma }_{r}^{}{}_{}{}^{j=1,\mathrm{},r+1}\mathrm{𝟏}_R`$ (40)
where $`\mathrm{𝟏}_{L,R}`$ are the unit-matrices in the respective spaces, while $`\sigma _x=e_{12}+e_{21}`$ and $`\sigma _y=ie_{12}+ie_{21}`$ are the $`2`$-dimensional Pauli matrices. $`\mathrm{\Gamma }_{r}^{}{}_{}{}^{r+1}`$ corresponds to the “generalized $`\mathrm{\Gamma }^5`$-matrix” in $`r+1`$ dimensions. In the above formula the values $`r,s=0`$ are allowed. The corresponding $`\mathrm{\Gamma }_{0}^{}{}_{}{}^{1}`$ is just $`1`$.
By iteratively applying the above algorithm starting from $`1`$ (that is $`r,s=0`$), we obtain as a first step the $`2`$-dimensional Pauli matrices and next, from $`r=0`$, $`s=2`$, the $`4`$-dimensional MW representation as a second step. The $`6`$-dimensional MW representation ($`3_S+3_A`$) is obtained as a further step. It can be produced either from $`r=0`$, $`s=4`$, or $`r=2`$, $`s=2`$. The non-Weyl ($`0_S+6_A`$) representation is constructed with the method explained in appendix $`1`$.
The higher-dimensional MW representations as well are obtained via the algorithm. An explicit way of constructing them is, e.g., in accordance with the table
$`\begin{array}{ccccccc}r=6& (0_S+6_A)& & s=0& (1)& & (8_S+0_A)\\ r=6& (3_S+3_A)& & s=0& (1)& & (4_S+4_A)\\ r=0& (1)& & s=0& (0_S+6_A)& & (0_S+8_A)\\ & & & & & & \\ r=0& (1)& & s=8& (8_S+0_A)& & (9_S+1_A)\\ r=8& (4_S+4_A)& & s=0& (1)& & (5_S+5_A)\\ r=8& (8_S+0_A)& & s=0& (1)& & (1_S+9_A)\\ & & & & & & \\ r=0& (1)& & s=10& (9_S+1_A)& & (10_S+2_A)\\ r=0& (1)& & s=10& (5_S+5_A)& & (6_S+6_A)\\ r=0& (1)& & s=10& (1_S+9_A)& & (2_S+10_A)\\ & & & & & & \\ r=0& (1)& & s=12& (10_S+2_A)& & (11_S+3_A)\\ r=0& (1)& & s=12& (6_S+6_A)& & (7_S+7_A)\\ r=0& (1)& & s=12& (2_S+10_A)& & (3_S+11_A)\\ & & & & & & \\ r=6& (0_S+6_A)& & s=8& (8_S+0_A)& & (16_S+0_A)\\ r=0& (1)& & s=14& (11_S+3_A)& & (12_S+4_A)\\ r=0& (1)& & s=14& (7_S+7_A)& & (8_S+8_A)\\ r=0& (1)& & s=14& (3_S+11_A)& & (4_S+12_A)\\ r=8& (8_S+0_A)& & s=6& (0_S+6_A)& & (0_S+16_A)\\ & & & & & & \\ r=0& (1)& & s=16& (16_S+0_A)& & (17_S+1_A)\\ r=0& (1)& & s=16& (12_S+4_A)& & (13_S+5_A)\\ r=0& (1)& & s=16& (8_S+8_A)& & (9_S+9_A)\\ r=0& (1)& & s=16& (4_S+12_A)& & (5_S+13_A)\\ r=0& (1)& & s=16& (0_S+16_A)& & (1_S+17_A)\end{array}`$ (63)
Notice that in order to explicitly construct all MW-representations up to $`d=18`$ with the help of formula (40), the only extra-knowledge of the $`(0_S+6_A`$) non-Weyl Majorana is required.
In the following we will refer to “space-time triality” as the $`8`$-dimensional property that the Majorana-Weyl condition is satisfied in three different signatures and will show its connection to the usual triality property of the $`8`$ dimensions, as well as the $`S_3`$ permutation group.
Due to the “lifting” formula (40), this feature is extended to higher-dimensional MW-spacetimes. Solutions in three different signatures arise as well in $`d=10`$, $`d=12`$ and $`d=14`$. Such solutions are a direct consequence of the embedding of the eight-dimensional Lorentz algebra into higher dimensions. In dimensions higher than $`8`$ the property that Majorana-Weyl conditions (or simply Majorana conditions in odd dimensions) are consistent in different signatures can therefore be regarded as a “derived” property which is fundamentally rooted in the $`8`$-dimensions. This argument holds even in the case (for $`d16`$), where solutions to the MW-constraints in more than three different signatures are obtained. As an example the $`d=18`$ case can be produced with the help of (40) for the values $`r=s=8`$. Therefore the $`5`$ different $`18`$-dimensional Majorana-Weyl representations are obtained from tensoring two $`8`$-dimensional Majorana-Weyl representations.
On a physical ground and not just for purely mathematical purposes, at the present state of the art we do not need getting involved into such complications since the most promising dimensions where the ultimate candidate theories for unification are expected to live correspond to $`d=10,11,12`$.
As a final comment in this section we emphasize once more that all data needed to define theories in such dimensions can be recovered, through a set of reconstruction formulas based on (40), from the $`8`$-dimensional data. In particular all “space-time trialities” in $`d>8`$ are encoded in the $`8`$-dimensional “space-time triality”. For this reason in the following we can concentrate ourselves in investigating the $`8`$-dimensional case.
## 4 The set of data for Majorana-Weyl supersymmetric theories.
In this section we present the set of data needed to specify a supersymmetric theory involving Majorana-Weyl spinors. The most suitable basis one can use in this case is the Majorana-Weyl basis previously discussed, where all spinors are either real or imaginary. In such a representation the following set of data underlines any given theory:
i) the vector fields (or, in the string/brane picture, the bosonic coordinates of the target $`x_m`$), specified by a vector index denoted by $`m`$;
ii) the spinor fields (or, in the string/brane picture, the fermionic coordinates of the target $`\psi _a`$, $`\chi _{\dot{a}}`$), specified by chiral and antichiral indices $`a`$, $`\dot{a}`$ respectively;
iii) the diagonal (pseudo-)orthogonal spacetime metric $`(g^1)^{mn}`$, $`g_{mn}`$ which we will assume to be flat;
iv) the $`𝒜`$ matrix introduced in section $`2`$, used to define barred spinors, coinciding with the $`\mathrm{\Gamma }^0`$-matrix in the Minkowski case; in a MWR is decomposed in an equal-size block diagonal form such as $`𝒜=A\stackrel{~}{A}`$, with structure of indices $`(A)_{a}^{}{}_{}{}^{b}`$ and $`(\stackrel{~}{A})_{\dot{a}}^{}{}_{}{}^{\dot{b}}`$ respectively;
v) the charge-conjugation matrix $`𝒞`$ which also appears in an equal-size block diagonal form $`𝒞=C^1\stackrel{~}{C}^1`$. It is invariant under bispinorial transformations and it can be promoted to be a metric in the space of chiral (and respectively antichiral) spinors, used to raise and lower spinorial indices. Indeed we can set $`(C^1)^{ab}`$, $`(C)_{ab}`$, and $`(\stackrel{~}{C}^1)^{\dot{a}\dot{b}}`$, $`(\stackrel{~}{C})_{\dot{a}\dot{b}}`$;
vi) the $`\mathrm{\Gamma }`$-matrices, which are decomposed in equal-size blocks as in (210), where the $`\sigma ^m`$’s are upper-right blocks and the $`\stackrel{~}{\sigma }^m`$’s lower-left blocks having structure of indices $`(\sigma ^m)_{a}^{}{}_{}{}^{\dot{b}}`$ and $`(\stackrel{~}{\sigma }^m)_{\dot{a}}^{}{}_{}{}^{b}`$ respectively;
vii) the $`\eta =\pm 1`$ sign, labeling the two inequivalent choices for $`𝒞`$.
We recall that by definition the $``$ matrix is automatically set to be the identity ($`=\mathrm{𝟏}`$) in a Majorana-Weyl representation.
The above structures are common in any theory involving Majorana-Weyl spinors. In the following we will furnish a dictionary relating Majorana-Weyl spacetimes with the same dimensionality, but with different signatures. The structures i)-vii) will be related via triality transformations which close the $`S_3`$ permutation group. They constitute the “words” in a three-language dictionary. According to the discussion in the previous section without loss of generality we can limit ourselves in analyzing the $`d=8`$-dimensional case. In this particular dimension the three indices, vector ($`m`$), chiral ($`a`$) and antichiral ($`\dot{a}`$) take values $`m,a,\dot{a}\{1,\mathrm{},8\}`$.
We mention that in the $`(4+4)`$-signature the $`(4_S+4_A)`$-representation of the $`\mathrm{\Gamma }`$-matrices has to be employed for both values of $`\eta `$ in order to provide a Majorana-Weyl basis. In the ($`t=8`$, $`s=0`$) signature the $`(8_S+0_A)`$-representation offers a MW basis for $`\eta =+1`$, while the $`(0_S+8_A)`$ offers it for $`\eta =1`$. The converse is true in the ($`t=0`$, $`s=8`$)-signature.
The three $`8`$-dimensional MW-type representations are explicitly constructed in the appendices $`2`$-$`4`$, together with their charge-conjugation matrices $`𝒞`$’s.
## 5 The Cartan’s V-C-A triality.
In this section we review the basic features of the Cartan’s triality involving vectors, chiral and antichiral spinors of $`SO(8)`$ and $`SO(4,4)`$.
The fundamental reason behind triality is the peculiar property of the $`D_4`$ Lie algebra, the only one admitting a group of symmetry for the corresponding Dynkin diagram other than the identity or $`𝐙_2`$. Its group of symmetry is the six-elements non-abelian group of permutations $`S_3`$ which, as well-known, corresponds to the outer automorphisms ($`OutAut/Int`$) of $`D_4`$.
The groups $`SO(8)`$ and $`SO(4,4)`$ are obtained by exponentiating different real forms of the $`D_4`$ Lie algebra.
For such groups and the corresponding metrics, the euclidean one with either all $`+`$ or all $``$ signs and the pseudoeuclidean metric $`(++++)`$, the spinor representations are Majorana-Weyl and satisfy the properties discussed in the previous section. In particular chiral and antichiral spinors can be consistently defined.
A unique and “miraculous” feature of the above spacetimes, not shared by any other case, consists in the fact that vectors, chiral and antichiral spinors (in short V-C-A’s) have the same dimensionality, being all eight-components.
Transformations exchanging V-C-A’s are found and, as we discuss later, they can all be identified. Such a property can be visualized with the following triangle diagram whose vertices are the interrelated V-C-A’s:
$`\begin{array}{ccccc}& & V& & \\ & & & & \\ C& & & & A\end{array}`$ (67)
Vectors ($`V_m`$), chiral ($`\psi _a`$) and antichiral ($`\chi _{\dot{a}}`$) spinors can be conveniently arranged into a single $`24`$-component “triality vector” $`T`$
$`T_M`$ $`=`$ $`\left(\begin{array}{c}V_m\\ \psi _a\\ \chi _{\dot{a}}\end{array}\right)`$ (71)
whose Lorentz-transformation properties are given by
$`T`$ $``$ $`T^{}=e^{\frac{1}{2}\omega _{mn}\mathrm{\Sigma }^{mn}}T`$ (72)
where
$`\mathrm{\Sigma }^{mn}`$ $`=`$ $`\mathrm{\Sigma }_{V}^{}{}_{}{}^{mn}\mathrm{\Sigma }_{C}^{}{}_{}{}^{mn}\mathrm{\Sigma }_{A}^{}{}_{}{}^{mn}`$ (73)
while $`(\mathrm{\Sigma }_{V,C,A})^{mn}`$ are the Lorentz-generators for vectors, chiral and antichiral spinors respectively:
$`(\mathrm{\Sigma }_{V}^{}{}_{}{}^{mn})_{k}^{}{}_{}{}^{l}`$ $`=`$ $`\delta _{}^{m}{}_{k}{}^{}(g^1)^{nl}(g^1)^{ml}\delta _{}^{n}{}_{k}{}^{}`$
$`\mathrm{\Sigma }_{C}^{}{}_{}{}^{mn}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(\sigma ^m\stackrel{~}{\sigma }^n\sigma ^n\stackrel{~}{\sigma }^m)`$
$`\mathrm{\Sigma }_{A}^{}{}_{}{}^{mn}`$ $`=`$ $`{\displaystyle \frac{1}{4}}(\stackrel{~}{\sigma }^m\sigma ^n\stackrel{~}{\sigma }^n\sigma ^m)`$ (74)
As far as the Lorentz transformation properties alone are concerned, there is no need to discuss (and even introduce) the character, commuting or anticommuting, of spinors. However, the triality admits to be interpreted as an invariance property which allows to introduce the $`𝒢_{Tr}`$ triality group of symmetry (to be defined below). For such interpretation we need to specify whether the spinors are assumed commuting or anticommuting. Following Cartan we discuss here the case of commuting spinors. We will comment the modifications occurring when anticommuting spinors are taken into account<sup>2</sup><sup>2</sup>2Commuting spinors too are relevant in the physical literature, not only as supersymmetric ghosts in a BRST quantization scheme; they also appear, e.g., in the super-twistors quantization approach, see . The fudamental reference on twistors is ..
For both values of $`\eta =\pm 1`$, the following bilinear Lorentz-invariants can be introduced
$`_V`$ $`=`$ $`V_{}^{T}{}_{m}{}^{}(g^1)^{mn}V_n`$
$`_C`$ $`=`$ $`\psi ^TC^1\psi `$
$`_A`$ $`=`$ $`\chi ^T\stackrel{~}{C}^1\chi `$ (75)
(we use the conventions discussed in the previous section), while the trilinear Lorentz-invariant
$`𝒯`$ $`=`$ $`\mathrm{\Psi }^T𝒞\mathrm{\Gamma }^m\mathrm{\Psi }V_m=2(\psi ^TC^1\sigma ^m\chi V_m),\mathrm{\Psi }\left(\begin{array}{c}\psi _a\\ \chi _{\dot{a}}\end{array}\right)`$ (78)
is non-vanishing for $`\eta =1`$. For anticommuting spinors the bilinears $`_𝒞`$, $`_𝒜`$ are identically vanishing, while $`𝒯`$ is non-vanishing for $`\eta =+1`$.
The group of invariances $`𝒢`$ is introduced as the group of linear homogeneous transformations acting on the $`8\times 3=24`$ dimensional space of “triality-vectors” leaving invariant, separately, $`_𝒱`$, $`_𝒞`$, $`_𝒜`$ and $`𝒯`$.
The group of triality $`𝒢_{Tr}`$ is defined by relaxing one condition, as the group of $`24`$-dimensional homogeneous linear transformations leaving invariant $`𝒯`$ and the total bilinear $`_{Sum}`$
$`_{Sum}`$ $`=`$ $`_𝒱+_𝒞+_𝒜`$ (79)
It can be proven that $`𝒢_{Tr}`$ is given by the semidirect product of $`𝒢`$ and the finite group $`S_3`$
$`𝒢_{Tr}`$ $`=`$ $`𝒢_SS_3`$ (80)
This result directly follows from the Cartan’s approach, even if it is not explicitly stated in the Cartan’s book<sup>3</sup><sup>3</sup>3He introduced the analogs of the $`𝒫`$, $``$ transformations of formula (81) which however, in his case, do not satisfy the relations (89) and cannot therefore be taken as the generators of the $`S_3`$ group.. The transformations in $`S_3`$ are obtained from the generators $`𝒫`$, $``$ whose action, symbolically, is given by
$`𝒫`$ $`:`$ $`VV,CA`$
$``$ $`:`$ $`VC,AA`$ (81)
As a consequence $`𝒫,`$ can be decomposed into eight-dimensional blocks matrices according to
$`𝒫=\left(\begin{array}{ccc}P_1& & \\ & & P_2\\ & P_3& \end{array}\right),`$ $`=\left(\begin{array}{ccc}& R_1& \\ R_2& & \\ & & R_3\end{array}\right)`$ (88)
$`𝒫`$, $``$ can be carefully chosen in such a way to satisfy the relations
$`𝒫^2=^2=\mathrm{𝟏},`$ $`(𝒫)^3=\mathrm{𝟏}`$ (89)
showing their nature of $`S_3`$ generators.
There is an arbitrariness in the choice of $`𝒫,`$ since any couple of gauge-transformed generators $`\widehat{𝒫}=g𝒫g^1`$, $`\widehat{}=gg^1`$, with $`g𝒢_{Tr}`$, satisfy the same properties and can be equally taken for generating $`S_3`$. It should be mentioned that, under a Lorentz transformation realized by $`e^{\omega \mathrm{\Sigma }}`$, the generators $`𝒫`$ ($``$) are mapped into $`𝒫^{}=e^{\omega \mathrm{\Sigma }}𝒫e_{}^{\omega \mathrm{\Sigma }}{}_{}{}^{1}`$ (and respectively $`^{}=e^{\omega \mathrm{\Sigma }}e_{}^{\omega \mathrm{\Sigma }}{}_{}{}^{1}`$). However, since we are free to introduce a gauge-compensating transformation, we can make invariant the choice of $`𝒫,`$ in any Lorentzian system of reference. As a consequence the constants $`𝒫,`$ matrices can be introduced in manifestly Lorentz-invariant actions if needed.
It is not easy, due to the intertwined nature of their relations (89), to determine a concrete realization for $`𝒫,`$. We arrived at it with a trial-and-error procedure. The final result is furnished in the appendices $`2`$-$`4`$ for each one of the three possible metrics.
If we disregard $`_{Sum}`$ and just demand the invariance of $`𝒯`$, in this case both $`𝒫,`$ can always be assumed real-valued. This is no longer true when the invariance of $`_{Sum}`$ is required, due to the fact that the three metrics on V-C-A’s can differ by an overall $``$ sign. This case is easily managed by noticing that the transformations $`P_2iP_2`$, $`P_3iP_3`$, $`P_1`$ unchanged (and $`R_1iR_1`$, $`R_2iR_2`$, $`R_3`$ unchanged) solves the problem without altering neither the (89) relations nor the $`𝒯`$-invariance requirement.
We finally comment that vectors, chiral and antichiral spinors can be identified under triality. If a basis in these three different spaces has already be chosen, the most natural way of identifying them ($`VX_1C`$, $`VY_1A`$) make use of the eight-dimensional invertible operators $`X_1`$, $`Y_1`$ entering
$`𝒫=\left(\begin{array}{ccc}& X_1& \\ & & X_2\\ X_3& & \end{array}\right),`$ $`𝒫=\left(\begin{array}{ccc}& & Y_1\\ Y_2& & \\ & Y_3& \end{array}\right)`$ (96)
The reason to use them is because $`X_1`$, $`Y_1`$ map into vectors respectively chiral and antichiral spinors with transformations which correspond to even permutations of V,C,A.
## 6 The signature triality.
In this section we discuss other consequences of the $`S_3`$ automorphisms of $`D_4`$. Besides being responsible for the Cartan’s V-C-A triality in fact, triality properties are associated with other structures. For purpose of clarity it will be convenient to represent them symbolically with triangle diagrams as the one shown in (67).
An extra-consequence of triality appears at the level of Majorana-Weyl representations for Clifford’s $`\mathrm{\Gamma }`$-matrices, see section $`3`$. The three different eight-dimensional representations can in fact be placed into the diagram
$`\begin{array}{ccccc}& & (4_S+4_A)& & \\ & & & & \\ (8_S+0_A)& & & & (0_S+8_A)\end{array}`$ (100)
which exhibits the triality operating at the level of the $`\mathrm{\Gamma }`$-matrices.
We have recalled that such MW-representations are associated with the space-time signature, and therefore triality can also be regarded as operating on space-times according to
$`\left(\begin{array}{ccccc}& & (5+5)& & \\ & & & & \\ (9+1)& & & & (1+9)\end{array}\right)`$ $``$ $`\left(\begin{array}{ccccc}& & (4+4)& & \\ & & & & \\ (8+0)& & & & (0+8)\end{array}\right)`$ (107)
The arrow has been inserted to recall that such triality can be lifted to higher dimensions or, conversely, that the $`8`$-dimensional spacetimes arise as transverse-coordinates spaces in physical theories.
The triality operating at the level of spacetime signatures is the one visualized by the previous diagram and will be described in this section.
Cartan’s V-C-A triality and signature triality can also be combined and symbolically represented by a sort of fractal-like double-triality diagram as follows
$`\begin{array}{ccccc}& & \begin{array}{ccccc}& & V& & \\ & & \left(\mathrm{𝟒}+\mathrm{𝟒}\right)& & \\ C& & & & A\end{array}& & \\ & & & & \\ & & & & \\ & & & & \\ & & & & \\ & & & & \\ \begin{array}{ccccc}& & V& & \\ & & \left(\mathrm{𝟖}+\mathrm{𝟎}\right)& & \\ C& & & & A\end{array}& & & & \begin{array}{ccccc}& & V& & \\ & & \left(\mathrm{𝟎}+\mathrm{𝟖}\right)& & \\ C& & & & A\end{array}\end{array}`$ (124)
The bigger triangle illustrates the signature triality, while the smaller triangles visualize the trialities for vectors, chiral and antichiral spinors living in each space-time.
We give now the explicit expression of the duality transformations relating theories formulated in the three above spacetimes or, in other words, the “translation rules” for the set of data discussed in section $`4`$.
To be definite we discuss the $`\eta =1`$ case (we recall that the $`\eta `$-sign has been introduced in section $`2`$); the modifications to be introduced in the $`\eta =+1`$ case are immediate.
Since we are working in a Majorana-Weyl basis it is always true that $`=\mathrm{𝟏}`$. The data defining our theories are therefore specified by the metric in the vectors’ space $`g^1`$, as well as the charge-conjugation matrix $`𝒞`$ which contains both the metric for chiral ($`C^1`$) and antichiral $`(\stackrel{~}{C}^1`$) spinors. In a Majorana-Weyl basis the $`𝒜`$ matrix is identified with $`𝒞`$ via the equation (2).
It is convenient to collectively denote as $`g_{}^{}{}_{}{}^{1}`$ (where $`V,C,A`$) the three metrics associated to, respectively, vectors and chiral and antichiral spinors in the $`(4+4)`$ spacetime. The analogous metrics when associated to the $`(8+0)`$ spacetime will be denoted with a tilde ($`\stackrel{~}{g}_{}^{}{}_{}{}^{1}`$), while the hat will denote the three metrics associated to the $`(0+8)`$ spacetime ($`\widehat{g}_{}^{}{}_{}{}^{1}`$).
Working in the $`\eta =1`$ case with the representations given in appendices $`24`$ we have
$`g_{V}^{}{}_{}{}^{1}=\mathrm{𝟏}_\mathrm{𝟒}\mathrm{𝟏}_\mathrm{𝟒},g_{C}^{}{}_{}{}^{1}=\mathrm{𝟏}_\mathrm{𝟒}\mathrm{𝟏}_\mathrm{𝟒},g_{A}^{}{}_{}{}^{1}=\mathrm{𝟏}_\mathrm{𝟒}\mathrm{𝟏}_\mathrm{𝟒};`$
$`\stackrel{~}{g}_{V}^{}{}_{}{}^{1}=\mathrm{𝟏}_\mathrm{𝟖},\stackrel{~}{g}_{C}^{}{}_{}{}^{1}=\mathrm{𝟏}_\mathrm{𝟖},\stackrel{~}{g}_{A}^{}{}_{}{}^{1}=\mathrm{𝟏}_\mathrm{𝟖};`$
$`\widehat{g}_{V}^{}{}_{}{}^{1}=\mathrm{𝟏}_\mathrm{𝟖},\widehat{g}_{C}^{}{}_{}{}^{1}=\mathrm{𝟏}_\mathrm{𝟖},\widehat{g}_{A}^{}{}_{}{}^{1}=\mathrm{𝟏}_\mathrm{𝟖};`$ (125)
The duality transformations can therefore be expressed by similarity transformations realized by non-orthogonal bridge matrices connecting the different metrics. We can introduce indeed the eight-dimensional bridge matrices $`\stackrel{~}{K}_{}`$, $`\widehat{K}_{}`$ such that
$`\stackrel{~}{g}_{}^{}{}_{}{}^{1}`$ $`=`$ $`\stackrel{~}{K}_{}g_{}^{}{}_{}{}^{1}\stackrel{~}{K}_{}^{}{}_{}{}^{T}`$
$`\widehat{g}_{}^{}{}_{}{}^{1}`$ $`=`$ $`\widehat{K}_{}g_{}^{}{}_{}{}^{1}\widehat{K}_{}^{}{}_{}{}^{T}`$ (126)
(as before “$``$” assumes the values $`V`$, $`C`$, $`A`$).
Vectors, chiral and antichiral spinors (V-C-A’s, collectively denoted as $`\phi _{}`$) in the $`(4+4)`$ spacetime are transformed into the $`(8+0)`$-signature V-C-A’s $`\stackrel{~}{\phi }_{}`$ according to
$`\stackrel{~}{\phi }_{}`$ $``$ $`(\stackrel{~}{K}_{}^{}{}_{}{}^{T})^1\phi _{}`$ (127)
An analogous transformation maps them into the ($`0+8)`$-signature V-C-A’s $`\widehat{\phi }_{}`$.
The $`16`$-dimensional matrices $`\stackrel{~}{H}`$, $`\widehat{H}`$ are constructed with the help of $`\stackrel{~}{K}_C,\stackrel{~}{K}_A`$ (and respectively $`\widehat{K}_C,\widehat{K}_A`$) according to
$`\stackrel{~}{H}`$ $`=`$ $`\stackrel{~}{K}_C\stackrel{~}{K}_A`$
$`\widehat{H}`$ $`=`$ $`\widehat{K}_C\widehat{K}_A`$ (128)
They are used to express the ($`8+0`$), and respectively ($`0+8`$), charge-conjugation matrices $`\stackrel{~}{𝒞}`$ ($`\widehat{𝒞}`$) in terms of the ($`4+4`$) charge-conjugation matrix $`𝒞`$ via the similarity transformation
$`\stackrel{~}{𝒞}`$ $`=`$ $`\stackrel{~}{H}𝒞\stackrel{~}{H}^T`$ (129)
and the corresponding equation obtained by replacing “$`\stackrel{~}{}`$” with “$`\widehat{}`$”.
For what concerns the Clifford’s $`\mathrm{\Gamma }`$ matrices in the ($`4+4`$)-signature, they are mapped into the $`(8+0)`$ $`\stackrel{~}{\mathrm{\Gamma }}`$’s according to
$`\mathrm{\Gamma }^m`$ $``$ $`\stackrel{~}{\mathrm{\Gamma }}^{\stackrel{~}{m}}=(\stackrel{~}{H}^T)^1\mathrm{\Gamma }^m\stackrel{~}{H}^T(\stackrel{~}{K}_{V}^{}{}_{}{}^{T})_{m}^{}{}_{}{}^{\stackrel{~}{m}}`$ (130)
As usual, an analogous relation maps them into the $`(0+8)`$-signature $`\widehat{\mathrm{\Gamma }}`$’s.
It is furthermore convenient to introduce the set of $`K_{}`$ matrices, defined through
$`K_{}`$ $`=`$ $`\stackrel{~}{K}_{}\widehat{K}_{},`$ (131)
connecting the $`(8+0)`$ with the $`(0+8)`$ signatures.
Let us denote with $`𝒲`$ (and respectively $`𝒵`$) the transformations mapping the different signatures set of data according to the following symbolic actions:
$`𝒲`$ $`:`$ $`(4+4)(4+4),(8+0)(0+8)`$
$`𝒵`$ $`:`$ $`(4+4)(8+0),(0+8)(0+8)`$ (132)
Such transformations are explicitly realized by $`24`$-dimensional matrices $`W_{}`$ (and respectively $`Z_{}`$) acting on column vectors of the kind $`(\varphi _{}\stackrel{~}{\varphi }_{}\widehat{\varphi }_{})`$. They are expressed in terms of the $`8`$-dimensional block matrices $`K_{}`$, $`\stackrel{~}{K}_{}`$. Indeed we can put them into the form
$`W_{}`$ $`=`$ $`\left(\begin{array}{ccc}\mathrm{𝟏}_\mathrm{𝟖}& 0& 0\\ 0& 0& K_{}\\ 0& K_{}& 0\end{array}\right)`$ (136)
$`Z_{}`$ $`=`$ $`\left(\begin{array}{ccc}0& \stackrel{~}{K}_{}& 0\\ \stackrel{~}{K}_{}& 0& 0\\ 0& 0& \mathrm{𝟏}_\mathrm{𝟖}\end{array}\right)`$ (140)
The matrices $`\stackrel{~}{K}_{}`$, $`K_{}`$ can be carefully chosen in such a way that the $`𝒲`$, $`𝒵`$ transformations can be regarded as the generators of the $`S_3`$ group (i.e. the signature triality group). This implies of course that the set of relations
$`𝒲^2=𝒵^2`$ $`=`$ $`\mathrm{𝟏}`$
$`(𝒲𝒵)^3`$ $`=`$ $`\mathrm{𝟏}`$ (141)
must be satisfied.
We mention here that when an odd number of signatures is flipped, as it happens for the standard Wick rotation from the Minkowski into the Euclidean space, the $`𝐙_2`$ group can not be realized with an action on the $`\mathrm{\Gamma }`$’s matrices. Indeed when, let’s say, the $`m=1`$ direction is flipped, the corresponding Clifford matrix $`\mathrm{\Gamma }^1`$ is mapped into $`i\mathrm{\Gamma }^1`$, which leads to a $`𝐙_4`$ group. A necessary condition to realize a $`𝐙_2`$ group on Clifford’s $`\mathrm{\Gamma }`$ matrices is that an even number of signatures has to be flipped (as it happens when changing e.g. $`(++)()`$). This is due to the fact that the $`\sigma _y=ie_{12}+ie_{21}`$ Pauli matrix, satisfying
$`\sigma _{y}^{}{}_{}{}^{2}`$ $`=`$ $`\mathrm{𝟏}`$ (142)
can be employed, allowing via a similarity transformation, to switch the sign
$`\sigma _y\mathrm{𝟏}_2\sigma _{y}^{}{}_{}{}^{T}`$ $`=`$ $`\mathrm{𝟏}_2`$ (143)
In the case here considered a consistent choice for $`K_{}`$, $`\stackrel{~}{K}_{}`$ is given by
$`K_V`$ $`=`$ $`\sigma _y\sigma _y\sigma _y\sigma _y`$
$`K_C`$ $`=`$ $`\mathrm{𝟏}_\mathrm{𝟖}`$
$`K_A`$ $`=`$ $`\sigma _y\sigma _y\sigma _y\sigma _y`$
$`\stackrel{~}{K}_V`$ $`=`$ $`\mathrm{𝟏}_\mathrm{𝟒}\sigma _y\sigma _y`$
$`\stackrel{~}{K}_C`$ $`=`$ $`\mathrm{𝟏}_\mathrm{𝟒}\sigma _y\sigma _y`$
$`\stackrel{~}{K}_A`$ $`=`$ $`\sigma _y\sigma _y\mathrm{𝟏}_\mathrm{𝟒}`$ (144)
It is a straightforward exercise to verify the consistency of the full set of relations (141) with the above choice of matrices.
For what concerns $`\widehat{K}_{}`$, they are immediately read from (131).
The above construction shows, as promised, that theories involving $`8`$-dimensional Majorana-Weyl spinors in different signatures can all be dually related in such a way to close the $`S_3`$ group of permutations.
## 7 An application: the brane-molecule scan.
We present in this section an application of the signature-dualities properties induced by triality. It concerns the so-called brane-molecule scan of ref. , originally appeared in the Blencowe-Duff paper in NPB, and later revisited in the Duff’s Tasi Lecture Notes. In these works the conditions for the existence of classical supersymmetric branes of arbitrary signatures embedded in flat target spaces whose signatures too are left arbitrary, are analyzed in details. Such conditions involve of course the property of spinors (Majorana, Weyl or Majorana-Weyl), plus extra identities involving the Clifford’s $`\mathrm{\Gamma }`$-matrices, which are needed in order to implement the kappa-symmetry.
The general result being presented in , here we just recall the allowed branes for $`d=10,11,12`$ dimensional targets. We have
i) the superstrings $`(1+1)\{(1+9),(5+5)or(9+1)\}`$,
ii) the five-branes $`\{(1+5)\{(1+9)or(5+5)\}`$, as well as the mirror copies
$`(5+1)\{(9+1)or(5+5)\}`$, obtained by exchanging space with time,
iii) the membranes $`(1+2)\{(9+2),(6+5),(5+6),(2+9),or(1+10)\}`$, as well as
the mirror copies obtained by space-time exchange,
iv) finally, the non-minkowskian $`(2+2)\{(2+10),(6+6)or(10+2)\}`$ branes.
In all the above cases the target space-times can be recovered from the eight-dimensional spacetimes, according to the discussion in section $`3`$. It is therefore clear that the “translation rules” expressing the various versions of the above theories in different signatures can be expressed in terms of the previous section results. In particular for the superstring, as well as for the $`(2+2)`$ brane, the three different versions are dually mapped in such a way to close the $`S_3`$ group.
In the original papers it was remarked the presence of a trivial $`𝐙_2`$ symmetry obtained by space-time exchange, connecting e.g. the $`(1+9)`$ to the $`(9+1)`$ target spacetimes solutions of the superstring. However, the presence of the extra solution $`(5+5)`$ was left “unexplained”. It turns out to be connected with, let’s say, the $`(1+9)`$ solution via another, non-trivial, generator of the $`S_3`$ triality group.
The above triality property holds for other classes of supersymmetric theories, like the $`10`$-dimensional superYang-Mills theories, which are allowed in different signatures (as before in the $`(1+9),(5+5),(9+1)`$ cases).
## 8 Triality as an invariance.
In this section we wish to discuss another possible application of triality related to the Duff’s viewpoint that a fundamental theory of everything should explain not only the dimensionality, but as well the signature of the spacetime. According to the Hull and Hull-Khuri results , the different versions in each given signature of a given theory are dually mapped one into another and therefore all equivalent. This is also the main content of the previous sections analysis, where we pointed out the role played by the $`S_3`$ triality group in such a context.
However our own results admit another possible interpretation. Having identified $`S_3`$ as the duality group, a question that can be naturally raised is whether this group can be assumed as a group of invariance, providing a signature-independent framework for the description of our theories. All this is much in the same spirit as general relativity providing a coordinate independent scheme.
Such a question deserves of course a careful investigation. It should be mentioned that it is not so difficult to realize an $`S_3`$-triality invariant formulation for strings, branes, etc. Other questions however, like the possibility that the $`S_3`$ group could be spontaneously broken, have no answer at present and need further investigation. A possible mechanism for producing a spontaneous breaking could use potential terms induced by the trilinear term given in formula (78).
The most natural setting to investigate such questions seems to be the eight-dimensional light-cone formulations of superstrings and branes, since the triality can be manifestly realized in this dimensionality.
We conclude this section pointing out the property that triality induced by supersymmetry strongly constraints the number of finite groups allowing a space-versus-time exchange symmetry. In bosonic theories, since all signatures are consistent and therefore allowed, the number of such groups is huge. In the supersymmetric case things are different. Let us discuss the $`10`$-dimensional superstring theories to be definite. The allowed signatures of the target spacetimes are $`(1+9),(5+5),(9+1)`$. If a space-versus-time coordinate exchange symmetry is required, as a consequence only three groups arise, namely
i) the identity group $`\mathrm{𝟏}`$, corresponding to a theory formulated in the single spacetime $`(5+5)`$,
ii) the $`𝐙_2`$ group of symmetry, for a theory formulated by using two spacetimes copies $`(1+9)`$,
$`(9+1)`$,
iii) finally, the $`S_3`$ group, which underlines a unified “space-time” theory requiring the whole
set of $`10`$-dimensional Majorana-Weyl spacetimes $`(1+9),(5+5),(9+1)`$.
## 9 Conclusions.
In this paper we have investigated various consequences of the triality automorphisms of $`Spin(8)`$. We have discussed in a suitable Majorana-Weyl basis (and partially extended) the original Cartan’s results concerning the transformation properties of $`8`$-dimensional vectors, chiral and antichiral spinors. Moreover, we have shown how triality affects the representation properties of the Clifford $`\mathrm{\Gamma }`$ matrices (the so-called Majorana-type representations). We have pointed out that triality naturally encodes the property that supersymmetry in higher-dimensions are consistently formulated in different-signature spacetimes. The $`S_3`$ triality group quite naturally arises in such a context. Indeed, we have been able to prove that different signature-versions of a given supersymmetric theory such, to give an example, as superstrings, can be dually transformed one into another with trasformations, conveniently chosen, which allow closing the $`S_3`$ group.
The possible role of $`S_3`$ as a symmetry group for a space-versus-time exchange invariance has also been mentioned.
One of the main motivations to investigate such properties concerned the still highly speculative but quite fascinating topics that a supersymmetric fundamental theory could shed light on the nature of the time. This point of view has been advocated by Duff in a series of papers and recently gained increased attention due to the results of Hull and collaborators .
The fact that supersymmetric theories seemingly related with the supposed $`F`$-theory find a natural formulation in a two-time physical world poses of course a strong challenge. Bars and collaborators , e.g., are exploring the possibility that the minkowksian time arises as a gauge-fixing of such a two-time physical world.
The present paper fits into this line of research. It is worth stressing however that, no matter which was the original motivation, the present paper contains a detailed account of mathematical results and property of triality which can be possibly used in technical analysis in different and more down-to-earth contexts. Just to mention an example which is currently under investigation, the here presented “translation rules” among different signatures can be used to define, let’s say, superYang-Mills theories in a $`(5+5)`$ signature. Such theories can be dimensionally reduced to an AdS $`(++)`$ signature. Standard minkowskian $`10`$-dimensional superYang-Mills theories do not admit of course such a dimensional reduction to AdS.
It is furthermore worth mentioning that some of the mathematical results here presented seems new, we have been unable to find them in the existing mathematical literature, at least not in such an explicit form as here given.
Acknowledgments
We are pleased to acknowledge J. A. Helayël-Neto and L.P. Colatto for both encouragement and helpful discussions. We are grateful to DCP-CBPF for the kind hospitality. F.T. acknowledges financial support from CNPq.
Appendix $`1`$:
The $`(0_S+6_A)`$-representation for the $`6`$-dimensional $`\mathrm{\Gamma }`$’s.
For completeness we furnish here a representation of the $`6`$-dimensional Clifford’ s $`\mathrm{\Gamma }`$-matrices realized in terms of antisymmetric $`\mathrm{\Gamma }`$’s only. As discussed in the text, this realization allows to construct all Majorana-type representations for Clifford’s $`\mathrm{\Gamma }`$-matrices up to dimension $`d=12`$ with the help of the recursive formula (40).
The representation given below is the correct one to be used in order to express a $`6`$-dimensional Euclidean space in the Majorana basis, which is consistent for $`t=6`$, $`s=0`$ when $`\eta =1`$, as well as for $`t=0`$, $`s=6`$ and $`\eta =+1`$.
We constructed the $`(0_S+6_A)`$-representation out of the $`(3_S+3_A)`$-representation (which is easily obtained through (40)) after computing in such a case, for one of the above Euclidean spaces, the value of the $``$-matrix and later finding a transformation $`U`$ mapping $`U^{}U^{}=\mathrm{𝟏}`$. As a consequence we obtain $`\mathrm{\Gamma }_{}^{i}{}_{(0S+6A)}{}^{}=U\mathrm{\Gamma }_{}^{i}{}_{(3S+3A)}{}^{}U^{}`$ according to (3). The extra $`\mathrm{\Gamma }^7`$ matrix is also antisymmetric.
The final result, here presented for all $`\mathrm{\Gamma }`$’s real (i.e. $`\mathrm{\Gamma }_{}^{i}{}_{}{}^{2}=\mathrm{𝟏}`$), is given by:
$`\mathrm{\Gamma }^1=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\end{array}\right)`$ $`\mathrm{\Gamma }^2=\left(\begin{array}{cccccccc}0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\end{array}\right)`$ (161)
$`\mathrm{\Gamma }^3=\left(\begin{array}{cccccccc}0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\end{array}\right)`$ $`\mathrm{\Gamma }^4=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\end{array}\right)`$ (178)
$`\mathrm{\Gamma }^5=\left(\begin{array}{cccccccc}0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\end{array}\right)`$ $`\mathrm{\Gamma }^6=\left(\begin{array}{cccccccc}0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\end{array}\right)`$ (195)
$`\mathrm{\Gamma }^7=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\end{array}\right)`$ (204)
Appendix $`2`$:
The ($`4_S+4_A`$)-representation for the $`8`$-dimensional $`\mathrm{\Gamma }`$’s.
In the following three appendices we explicitly furnish the three Majorana-Weyl representations for the $`8`$-dimensional $`\mathrm{\Gamma }`$-matrices. Moreover in each case a specific realization of the $`𝒫`$, $``$ generators of the $`S_3`$ permutation group, leaving invariant the trilinear term (78) for the choice $`\eta =1`$, i.e. for commuting spinors, is given (in the opposite case, $`\eta =+1`$, the trilinear term is automatically vanishing for commuting spinors).
The $`\mathrm{\Gamma }`$’s, as well as the generators $`𝒫,`$, are presented in real form (confront the discussion in section $`5`$).
Each one of the three representations, being MW, admits a $`\mathrm{\Gamma }^9`$ matrix of the kind
$`\mathrm{\Gamma }^9`$ $`=`$ $`\left(\begin{array}{cc}\mathrm{𝟏}_\mathrm{𝟖}& 0\\ 0& \mathrm{𝟏}_\mathrm{𝟖}\end{array}\right)`$ (207)
while the $`\mathrm{\Gamma }^i`$ for $`i=1,2,\mathrm{},8`$ are decomposed according to
$`\mathrm{\Gamma }^i`$ $`=`$ $`\left(\begin{array}{cc}0& \sigma ^i\\ \stackrel{~}{\sigma }^i& 0\end{array}\right)`$ (210)
We point out that the $`8`$-dimensional $`\mathrm{\Gamma }`$-matrices introduced in this one and the two following appendices are not directly obtainable through the recursion formula (40). Rather, they have been conveniently chosen in order to provide a diagonal metric for both chiral and antichiral spinors which coincides, up to an overall sign, with the given diagonal metric for vectors.
In this appendix we present the results for the ($`4_S+4_A`$)-representation, the one which has to be used in order to introduce MW-spinors in a $`t=4`$, $`s=4`$ spacetime.
We have
$`\sigma ^1=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\end{array}\right)`$ $`\sigma ^2=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\end{array}\right)`$ (227)
$`\sigma ^3=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\end{array}\right)`$ $`\sigma ^4=\left(\begin{array}{cccccccc}0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\end{array}\right)`$ (244)
$`\sigma ^5=\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)`$ $`\sigma ^6=\left(\begin{array}{cccccccc}0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\end{array}\right)`$ (261)
$`\sigma ^7=\left(\begin{array}{cccccccc}0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\end{array}\right)`$ $`\sigma ^8=\left(\begin{array}{cccccccc}0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\end{array}\right)`$ (278)
while $`\stackrel{~}{\sigma }_i=\sigma _{i}^{}{}_{}{}^{T}`$ for $`i=1,2,3,4`$ and $`\stackrel{~}{\sigma }_i=\sigma _{i}^{}{}_{}{}^{T}`$ for $`i=5,6,7,8`$, (i.e. the first four $`\mathrm{\Gamma }`$-matrices are antisymmetric, the last four ones symmetric).
The diagonal charge-conjugation matrices are given by
$`C^1`$ $`=`$ $`\mathrm{𝟏}_\mathrm{𝟒}\mathrm{𝟏}_\mathrm{𝟒}`$
$`\stackrel{~}{C}^1`$ $`=`$ $`\eta C^1`$ (279)
A convenient basis for $`𝒫,`$ is
$`𝒫_{VV}𝒫_1`$ $`=`$ $`\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)`$ (288)
$`𝒫_{AC}𝒫_2`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\end{array}\right)`$ (297)
$`𝒫_{CA}𝒫_3`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\end{array}\right)`$ (306)
and
$`_{CV}_1`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\end{array}\right)`$ (315)
$`_{VC}_2`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\end{array}\right)`$ (324)
$`_{AA}_3`$ $`=`$ $`\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)`$ (333)
Appendix $`3`$:
The ($`0_S+8_A`$)-representation for the $`8`$-dimensional $`\mathrm{\Gamma }`$’s.
This representation is the one to be used in order to introduce a MW-basis for an Euclidean $`t=8`$, $`s=0`$ space when $`\eta =1`$. In this case all $`\mathrm{\Gamma }`$’s have to be assumed imaginary (in the following formulas we drop, as usual, the $`i`$).
The following formulas are directly obtainable from those presented in appendix 2 after applying a special case of the transformation (3) discussed in section $`2`$. They are here presented for completeness. We have
$`\sigma ^1=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\end{array}\right)`$ $`\sigma ^2=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\end{array}\right)`$ (350)
$`\sigma ^3=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\end{array}\right)`$ $`\sigma ^4=\left(\begin{array}{cccccccc}0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\end{array}\right)`$ (367)
$`\sigma ^5=\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)`$ $`\sigma ^6=\left(\begin{array}{cccccccc}0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\end{array}\right)`$ (384)
$`\sigma ^7=\left(\begin{array}{cccccccc}0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\end{array}\right)`$ $`\sigma ^8=\left(\begin{array}{cccccccc}0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\end{array}\right)`$ (401)
Since the $`\mathrm{\Gamma }^i`$ are antisymmetric it follows that $`\stackrel{~}{\sigma }^i=\sigma _{}^{i}{}_{}{}^{T}`$ for any $`i=1,2,\mathrm{},8`$.
The two diagonal charge-conjugation matrices are
$`C^1`$ $`=`$ $`\mathrm{𝟏}_\mathrm{𝟖}`$
$`\stackrel{~}{C}^1`$ $`=`$ $`\eta \mathrm{𝟏}_\mathrm{𝟖}`$ (402)
A real basis for $`𝒫,`$ is obtained by transforming the corresponding mappings in the $`(4_S+4_A)`$ case. It is given by
$`𝒫_{VV}𝒫_1`$ $`=`$ $`\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)`$ (411)
$`𝒫_{AC}𝒫_2`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\end{array}\right)`$ (420)
$`𝒫_{CA}𝒫_3`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\end{array}\right)`$ (429)
and
$`_{CV}_1`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\end{array}\right)`$ (438)
$`_{VC}_2`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\end{array}\right)`$ (447)
$`_{AA}_3`$ $`=`$ $`\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)`$ (456)
Appendix $`4`$:
The ($`8_S+0_A`$)-representation for the $`8`$-dimensional $`\mathrm{\Gamma }`$’s.
This representation is the one to be used in order to introduce a MW-basis for an Euclidean $`t=0`$, $`s=8`$ space when $`\eta =1`$. In this case as well as all $`\mathrm{\Gamma }`$’s have to be assumed imaginary (in the following formulas we drop, as usual, the $`i`$).
As before the following formulas are obtained after transforming the corresponding formulas in the ($`4_S+4_A`$) case. We have
$`\sigma ^1=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\end{array}\right)`$ $`\sigma ^2=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\end{array}\right)`$ (473)
$`\sigma ^3=\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\end{array}\right)`$ $`\sigma ^4=\left(\begin{array}{cccccccc}0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\end{array}\right)`$ (490)
$`\sigma ^5=\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)`$ $`\sigma ^6=\left(\begin{array}{cccccccc}0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\end{array}\right)`$ (507)
$`\sigma ^7=\left(\begin{array}{cccccccc}0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\end{array}\right)`$ $`\sigma ^8=\left(\begin{array}{cccccccc}0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\end{array}\right)`$ (524)
while now $`\stackrel{~}{\sigma }^i=\sigma _{}^{i}{}_{}{}^{T}`$ for any $`i=1,2,\mathrm{},8`$.
The two diagonal charge conjugation matrices are given by
$`C^1`$ $`=`$ $`\mathrm{𝟏}_\mathrm{𝟖}`$
$`\stackrel{~}{C}^1`$ $`=`$ $`\eta \mathrm{𝟏}_\mathrm{𝟖}`$ (525)
A real basis for $`𝒫,`$ is given by
$`𝒫_{VV}𝒫_1`$ $`=`$ $`\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)`$ (534)
$`𝒫_{AC}𝒫_2`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\end{array}\right)`$ (543)
$`𝒫_{CA}𝒫_3`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\end{array}\right)`$ (552)
and
$`_{CV}_1`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\end{array}\right)`$ (561)
$`_{VC}_2`$ $`=`$ $`\left(\begin{array}{cccccccc}0& 0& 0& 0& 0& 0& 0& 1\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\end{array}\right)`$ (570)
$`_{AA}_3`$ $`=`$ $`\left(\begin{array}{cccccccc}1& 0& 0& 0& 0& 0& 0& 0\\ 0& 1& 0& 0& 0& 0& 0& 0\\ 0& 0& 1& 0& 0& 0& 0& 0\\ 0& 0& 0& 1& 0& 0& 0& 0\\ 0& 0& 0& 0& 1& 0& 0& 0\\ 0& 0& 0& 0& 0& 1& 0& 0\\ 0& 0& 0& 0& 0& 0& 1& 0\\ 0& 0& 0& 0& 0& 0& 0& 1\end{array}\right)`$ (579)
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# On 𝑡-structures and Torsion Theories Induced by Compact Objects
## 0. Introduction
In the representation theory of finite dimensional algebras, torsion theories were studied by several authors in connection with classical tilting modules. For these torsion theories, there are equivalences between torsion (resp., torsionfree) classes and torsionfree (resp., torsion) classes, which is known as Theorem of Brenner and Butler (\[HR\]). One of the authors gave one to one correspondence between classical tilting modules and torsion theories with certain conditions (\[Ho1\], \[Ho2\]). But in the case of a self-injective algebra $`A`$, tilting modules are essentially isomorphic to $`A`$. In \[Ri\], Rickard introduced the notion of tilting complexes as a generalization of tilting modules, and showed that these complexes induce equivalences between derived categories of module categories. Tilting complexes of term length two are often studied in the case of self-injective algebras (e.g. \[Hl\], \[HK\]). On the other hand, for triangulated categories, Beilinson, Bernstein and Deligne introduced the notions of $`t`$-structures and admissible abelian subcategories, and studied relationships between them (\[BBD\]). In this paper, first, we deal with a compact object $`C`$ in a triangulated category, and study a $`t`$-structure induced by $`C`$. Second, in an abelian category $`𝒜`$ we deal with a complex $`P^{}`$ of small projective objects of term length two and study a torsion theory induced by $`P^{}`$.
In Section 1, we show that a compact object $`C`$ in a triangulated category $`𝒯`$, which satisfies suitable conditions, induces a $`t`$-structure $`(𝒯^0(C),𝒯^0(C))`$, and its core $`𝒯^0(C)`$ is equivalent to the category $`\mathrm{𝖬𝗈𝖽}B`$ of left $`B`$-modules, where $`B=\mathrm{End}_𝒯(C)^{\mathrm{op}}`$ (Theorem 1.3). In Section 2, we define subcategories $`𝒳(P^{})`$, $`𝒴(P^{})`$ of an abelian category $`𝒜`$ satisfying the condition Ab4, and show when $`(𝒳(P^{}),𝒴(P^{}))`$ is a torsion theory (Theorem 2.10). Furthermore, we show that if $`P^{}`$ induces a torsion theory $`(𝒳(P^{}),𝒴(P^{}))`$ for $`𝒜`$, then the core $`𝖣(𝒜)^0(P^{})`$ is admissible abelian, and then there is a torsion theory $`(𝒴(P^{})[1],𝒳(P^{}))`$ for $`𝖣(𝒜)^0(P^{})`$ (Theorem 2.15). In Section 3, we apply results of Section 2 to module categories. We characterize a torsion theory for the category $`\mathrm{𝖬𝗈𝖽}A`$ of left $`A`$-modules, and for its core $`𝖣(\mathrm{𝖬𝗈𝖽}A)^0(P^{})`$ (Theorems 3.5 and 3.8). Furthermore, using a torsion theory, we give equivalent conditions for $`P^{}`$ to be a tilting complex (Corollary 3.6). In Section 4, We show that, if $`P^{}`$ is a tilting complex, then it induces equivalences between torsion theories for $`\mathrm{𝖬𝗈𝖽}A`$ and for $`\mathrm{𝖬𝗈𝖽}B`$, where $`B=\mathrm{End}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(P^{})^{\mathrm{op}}`$ (Theorem 4.4). In Section 5, in the case of artin algebras, if a torsion theory $`(𝒳,𝒴)`$ satisfies certain conditions, then there exists a tilting complex $`P^{}`$ of term length two such that a torsion theory $`(𝒳,𝒴)`$ coincides with $`(𝒳(P^{}),𝒴(P^{}))`$ (Theorem 5.8). As a consequence, we have one to one correspondence between tilting complexes of term length two and torsion theories with certain conditions (Corollary 3.7, Propositions 5.5, 5.7 and Theorem 5.8).
## 1. $`t`$-structures Induced by Compact Objects
In this section, we deal with a triangulated category $`𝒯`$ and its full subcategory $`𝒞`$. We will call $`𝒞`$ admissible abelian provided that $`\mathrm{Hom}_𝒯(𝒞,𝒞[n])=0`$ for $`n<0`$, and that all morphisms in $`𝒞`$ are $`𝒞`$-admissible in the sense of \[BBD\], 1.2.3. In this case, according to \[BBD\], Proposition 1.2.4, $`𝒞`$ is an abelian category. A triangulated category $`𝒯`$ is said to contain direct sums if direct sums of objects indexed by any set exist in $`𝒯`$. An object $`C`$ of $`𝒯`$ is called compact if $`\mathrm{Hom}_𝒯(C,)`$ commutes with direct sums. Furthermore, a collection $`𝒮`$ of compact objects of $`𝒯`$ is called a generating set provided that $`X=0`$ whenever $`\mathrm{Hom}_𝒯(𝒮,X)=0`$, and that $`𝒮`$ is stable under suspension (see \[Ne\] for details). For an object $`C𝒯`$ and an integer $`n`$, we denote by $`𝒯^n(C)`$ (resp., $`𝒯^n(C)`$) the full subcategory of $`𝒯`$ consisting of $`X𝒯`$ with $`\mathrm{Hom}_𝒯(C,X[i])=0`$ for $`i<n`$ (resp., $`i>n`$), and set $`𝒯^0(C)=𝒯^0(C)𝒯^0(C)`$.
For an abelian category $`𝒜`$, we denote by $`𝖢(𝒜)`$ the category of complexes of $`𝒜`$, and denote by $`𝖣(𝒜)`$ (resp., $`𝖣^+(𝒜)`$, $`𝖣^{}(𝒜)`$, $`𝖣^\mathrm{b}(𝒜)`$) the derived category of complexes of $`𝒜`$ (resp., complexes of $`𝒜`$ with bounded below homologies, complexes of $`𝒜`$ with bounded above homologies, complexes of $`𝒜`$ with bounded homologies). For an additive category $``$, we denote by $`𝖪()`$ (resp., $`𝖪^{}()`$, $`𝖪^\mathrm{b}()`$) the homotopy category of complexes of $``$ (resp., bounded above complexes of $``$, bounded complexes of $``$) (see \[RD\] for details).
###### Proposition 1.1.
Let $`𝒯`$ be a triangulated category which contains direct sums, $`C`$ a compact object satisfying $`\mathrm{Hom}_𝒯(C,C[n])=0`$ for $`n>0`$. Then for any $`r`$ and any object $`X𝒯`$, there exist an object $`X^r𝒯^r(C)`$ and a morphism $`\alpha ^r:XX^r`$ in $`𝒯`$ such that
1. for any $`ir`$, $`\mathrm{Hom}_𝒯(C,\alpha ^r[i])`$ is an isomorphism,
2. for every object $`Y𝒯^r(C)`$, $`\mathrm{Hom}_𝒯(\alpha ^r,Y)`$ is an isomorphism.
###### Proof.
Let $`X_0=X`$. For $`n1`$, by induction we construct a distinguished triangle
$$C_n[nr]\stackrel{g_n}{}X_{n1}\stackrel{h_n}{}X_n$$
as follows. If $`\mathrm{Hom}_𝒯(C,X_{n1}[rn])=0`$, then we set $`C_n=0`$. Otherwise, we take a direct sum $`C_n`$ of copies of $`C`$ and a morphism $`g_n^{}:C_nX_{n1}[rn]`$ such that $`\mathrm{Hom}_𝒯(C,g_n^{})`$ is an epimorphism, and let $`g_n=g_n^{}[nr]`$. Then, by easy calculation, we have the following:
$`(a)`$ $`\mathrm{Hom}_𝒯(C,X_n[i])=0\text{for}rni<r,`$
$`(b)`$ $`\mathrm{Hom}_𝒯(C,h_n[i])\text{is an isomorphism for any}n\text{and}ir.`$
Let $`X^r`$ be a homotopy colimit $`\underset{}{\mathrm{𝗁𝗈𝖼𝗈𝗅𝗂𝗆}}X_n`$ and $`\alpha ^r:XX^r`$ a structural morphism $`X_0\underset{}{\mathrm{𝗁𝗈𝖼𝗈𝗅𝗂𝗆}}X_n`$. According to \[Ne\], Lemma 2.8, the conditions ($`a`$), ($`b`$) imply that $`X^r`$ belongs to $`𝒯^r(C)`$ and satisfies the statement (i). For an object $`Y𝒯^r(C)`$, we have an exact sequence
$`\mathrm{Hom}_𝒯(C_n[nr],Y[j1])\mathrm{Hom}_𝒯(X_n,Y[j])`$ $``$
$`\mathrm{Hom}_𝒯(X_{n1},Y[j])`$ $`\mathrm{Hom}_𝒯(C_n[nr],Y[j]).`$
Since $`\mathrm{Hom}_𝒯(C[i],Y[j])=0`$ for $`ji<r`$, $`\mathrm{Hom}_𝒯(h_n,Y[j])`$ is an isomorphism for any $`n1`$ and $`j0`$. Then, we have an epimorphism
$$\underset{n}{}\mathrm{Hom}_𝒯(X_n,Y[j])\stackrel{1\text{shift}}{}\underset{n}{}\mathrm{Hom}_𝒯(X_n,Y[j])$$
for any $`j0`$. Therefore, we have an exact sequence
$$0\mathrm{Hom}_𝒯(X^r,Y)\underset{n}{}\mathrm{Hom}_𝒯(X_n,Y)\stackrel{1\text{shift}}{}\underset{n}{}\mathrm{Hom}_𝒯(X_n,Y)0.$$
Hence we have
$`\mathrm{Hom}_𝒯(X^r,Y)`$ $`\underset{}{\mathrm{lim}}\mathrm{Hom}_𝒯(X_n,Y)`$
$`\mathrm{Hom}_𝒯(X,Y).`$
###### Definition 1.2 (\[BBD\]).
Let $`𝒯`$ be a triangulated category. For full subcategories $`𝒯^0`$ and $`𝒯^0`$, $`(𝒯^0,𝒯^0)`$ is called a t-structure on $`𝒯`$ provided that
1. $`\mathrm{Hom}_𝒯(𝒯^0,𝒯^1)=0`$;
2. $`𝒯^0𝒯^1`$ and $`𝒯^0𝒯^1`$;
3. for any $`X𝒯`$, there exists a distinguished triangle
$$X^{}XX^{\prime \prime }$$
with $`X^{}𝒯^0`$ and $`X^{\prime \prime }𝒯^1`$,
where $`𝒯^n=𝒯^0[n]`$ and $`𝒯^n=𝒯^0[n]`$.
A t-structure $`(𝒯^0,𝒯^0)`$ on $`𝒯`$ is called non-degenerate if $`_n𝒯^n=_n𝒯^n=\{0\}`$.
###### Theorem 1.3.
Let $`𝒯`$ be a triangulated category which contains direct sums, $`C`$ a compact object satisfying $`\mathrm{Hom}_𝒯(C,C[n])=0`$ for $`n>0`$, and $`B=\mathrm{End}_𝒯(C)^{\mathrm{op}}`$. If $`\{C[i]:i\}`$ is a generating set, then the following hold.
1. $`(𝒯^0(C),𝒯^0(C))`$ is a non-degenerate t-structure on $`𝒯`$.
2. $`𝒯^0(C)`$ is admissible abelian.
3. The functor
$$\mathrm{Hom}_𝒯(C,):𝒯^0(C)\mathrm{𝖬𝗈𝖽}B$$
is an equivalence.
###### Proof.
(1) For any object $`X𝒯^0(C)`$, we take an object $`X^1𝒯^1(C)`$ and a morphism $`\alpha ^1:XX^1`$ satisfying the conditions of Proposition 1.1. Then for any $`Y𝒯^1(C)`$, by Proposition 1.1 (ii), we have
$$\mathrm{Hom}_𝒯(X^1,Y)\mathrm{Hom}_𝒯(X,Y).$$
By Proposition 1.1 (i), $`X𝒯^0(C)`$ implies that $`\mathrm{Hom}_𝒯(C,X^1[i])=0`$ for all $`i`$. Since $`\{C[i]:i\}`$ is a generating set, we have $`X^1=0`$, and hence $`\mathrm{Hom}_𝒯(X,Y)=0`$. It is easy to see that $`𝒯^0(C)𝒯^1(C)`$ and $`𝒯^0(C)𝒯^1(C)`$. For any object $`Z𝒯`$, we take an object $`Z^1𝒯^1(C)`$ and a morphism $`\alpha ^1:ZZ^1`$ satisfying the conditions of Proposition 1.1, and embed $`\alpha ^1`$ in a distinguished triangle
$$Z^{}ZZ^1.$$
Applying $`\mathrm{Hom}_𝒯(C,)`$ to the above triangle, by Proposition 1.1 (i), we have $`Z^{}𝒯^0(C)`$. Since $`\{C[i]:i\}`$ is a generating set, it is easy to see that $`(𝒯^0(C),𝒯^0(C))`$ is non-degenerate.
(2) Since $`𝒯^0(C)`$ is the core of the $`t`$-structure $`(𝒯^0(C),𝒯^0(C))`$, the assertion follows by \[BBD\], Theorem 1.3.6.
(3) Step 1: According to \[BBD\], Proposition 1.2.2, the short exact sequences in $`𝒯^0(C)`$ are just the distinguished triangles
$$XYZ$$
with $`X,Y`$ and $`Z`$ belonging to $`𝒯^0(C)`$. It follows that $`\mathrm{Hom}_𝒯(C,):𝒯^0(C)\mathrm{𝖬𝗈𝖽}B`$ is exact. Let $`M\mathrm{𝖬𝗈𝖽}B`$ and take a free presentation $`P_1P_0M0`$. We take $`C^{}=C^0𝒯^0(C)`$ and $`\alpha =\alpha ^0:CC^{}`$ satisfying the conditions of Proposition 1.1. Then there exist sets $`I,J`$ and a collection of morphisms $`h_{ij}:C^{}C^{}`$ such that
$$\begin{array}{ccc}P_1& & P_0\\ & & & & \\ \mathrm{Hom}_𝒯(C,C^{})^{(J)}& \stackrel{_{ij}\mathrm{Hom}(C,h_{ij})}{}& \mathrm{Hom}_𝒯(C,C^{})^{(I)}\end{array}$$
is commutative, where the vertical arrows are isomorphisms. We take an exact sequence in $`𝒯^0(C)`$
$$C_{}^{}{}_{}{}^{(J)}\stackrel{_{ij}h_{ij}}{}C_{}^{}{}_{}{}^{(I)}X0.$$
Since $`C`$ is compact, by the exactness of $`\mathrm{Hom}_𝒯(C,)`$, we have $`\mathrm{Hom}_𝒯(C,X)M`$.
Step 2: We show that $`\mathrm{Hom}_𝒯(C,)`$ reflects isomorphisms. Let
$$X\stackrel{𝑢}{}YZ$$
be a distinguished triangle in $`𝒯`$ with $`X,Y𝒯^0(C)`$ and with $`\mathrm{Hom}_𝒯(C,u)`$ an isomorphism. Then, by applying $`\mathrm{Hom}_𝒯(C,)`$, we get $`\mathrm{Hom}_𝒯(C,Z[n])=0`$ for all $`n`$, and hence $`Z=0`$. It follows that $`u`$ is an isomorphism.
Next, we show that $`\mathrm{Hom}_𝒯(C,)`$ is faithful. Let $`v:XY`$ be a morphism in $`𝒯^0(C)`$ with $`\mathrm{Hom}_𝒯(C,v)=0`$. By the exactness of $`\mathrm{Hom}_𝒯(C,)`$, $`\mathrm{Hom}_𝒯(C,\mathrm{Im}v)\mathrm{Im}\mathrm{Hom}_𝒯(C,v)=0`$. Since $`\mathrm{Im}v𝒯^0(C)`$, we have $`\mathrm{Hom}_𝒯(C,\mathrm{Im}v[n])=0`$ for all $`n`$, and hence $`\mathrm{Im}v=0`$ and $`v=0`$.
Let $``$ be the full subcategory of $`𝒯^0(C)`$ consisting of objects $`X`$ such that there exists an exact sequence $`C_1C_0X0`$ in $`𝒯^0(C)`$, where $`C_0,C_1`$ are direct sums of copies of $`C^{}`$. Since $`\mathrm{Hom}_𝒯(C,)`$ is faithful, by the same technique as in (1), it is not hard to see that $`\mathrm{Hom}_𝒯(C,)|_{}`$ is full dense, and hence an equivalence. It remains to show that $`=𝒯^0(C)`$. For an object $`X𝒯^0(C)`$, we have a commutative diagram
$$\begin{array}{ccccccc}\mathrm{Hom}_𝒯(C,C^{(J)})& \stackrel{\mathrm{Hom}(C,f)}{}& \mathrm{Hom}_𝒯(C,C^{(I)})& \stackrel{\mathrm{Hom}(C,g)}{}& \mathrm{Hom}_𝒯(C,X)& & 0\\ \mathrm{Hom}_𝒯(C,\alpha _1)& & \mathrm{Hom}_𝒯(C,\alpha _0)& & \\ \mathrm{Hom}_𝒯(C,C_{}^{}{}_{}{}^{(J)})& \stackrel{\mathrm{Hom}(C,f^{})}{}& \mathrm{Hom}_𝒯(C,C_{}^{}{}_{}{}^{(I)})\end{array}$$
with the top row being exact and with the vertical arrows being isomorphisms. And we have a commutative diagram in $`𝒯`$
$$\begin{array}{ccccc}C^{(J)}& \stackrel{f}{}& C^{(I)}& \stackrel{g}{}& X\\ \alpha _1& & \alpha _0& & \\ C_{}^{}{}_{}{}^{(J)}& \stackrel{f^{}}{}& C_{}^{}{}_{}{}^{(I)}\end{array}$$
with $`gf=0`$. By Proposition 1.1(ii), there exists $`h:C_{}^{}{}_{}{}^{(I)}X`$ such that $`g=h\alpha _0`$. Since $`\mathrm{Hom}_𝒯(C,hf^{})=0`$, we have $`hf^{}=0`$. Then there exists $`w:\mathrm{Cok}f^{}X`$ such that $`g=wg^{}\alpha _0`$, where $`g^{}:C_{}^{}{}_{}{}^{(I)}\mathrm{Cok}f^{}`$ is a canonical morphism. Then $`\mathrm{Hom}_𝒯(C,w)`$ is an isomorphism, and hence $`w`$ is an isomorphism and $`X\mathrm{Cok}f^{}`$. ∎
###### Remark 1.4.
Under the condition of Theorem 1.3, according to \[BBD\], Proposition 1.3.3, there exists a functor $`()^n:𝒯𝒯^n(C)`$ (resp., $`()^n:𝒯𝒯^n(C)`$) which is the right (resp., left) adjoint of the natural embedding functor $`𝒯^n(C)𝒯`$ (resp., $`𝒯^n(C)𝒯`$).
For an object $`C`$ in a triangulated category $`𝒯`$ and integers $`st`$, let $`𝒯^{[s]}(C)=𝒯^0(C)[s]`$, $`𝒯^{[s,t]}(C)=𝒯^t(C)𝒯^s(C)`$, and $`𝒯^\mathrm{b}(C)=(_n𝒯^n(C))(_n𝒯^n(C))`$. An object $`M`$ of an abelian category $`𝒜`$ is called small provided that $`\mathrm{Hom}_𝒜(M,)`$ commutes with direct sums in $`𝒜`$.
###### Corollary 1.5.
Let $`𝒜`$ be an abelian category satisfying the condition Ab4 (i.e. direct sums of exact sequences are exact) and $`T^{}`$ a bounded complex of small projective objects of $`𝒜`$ satisfying
1. $`\{T^{}[i]:i\}`$ is a generating set for $`𝖣(𝒜)`$,
2. $`\mathrm{Hom}_{𝖣(𝒜)}(T^{},T^{}[i])=0`$ for $`i0`$.
If either of the following conditions (1) or (2) is satisfied, then we have an equivalence of triangulated categories
$$𝖣(𝒜)^\mathrm{b}(T^{})𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}B),$$
where $`B=\mathrm{End}_{𝖣(𝒜)}(T^{})^{\mathrm{op}}`$.
1. $`𝒜`$ has enough projectives.
2. $`𝒜`$ has enough injectives and $`𝖣(𝒜)^0(T^{})𝖣^+(𝒜)`$.
Moreover, if $`𝖣(𝒜)^0(T^{})𝖣^\mathrm{b}(𝒜)`$, then we have an equivalence
$$𝖣^\mathrm{b}(𝒜)𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}B).$$
###### Proof.
According to \[BN\], Corollary 1.7, $`𝖣(𝒜)`$ contains direct sums. Since $`T^{}`$ is a bounded complex of small projective objects of $`𝒜`$, $`T^{}`$ is a compact object in $`𝖣(𝒜)`$. By Theorem 1.3 $`𝖣(𝒜)`$ has a $`t`$-structure $`(𝖣(𝒜)^0(T^{}),𝖣(𝒜)^0(T^{}))`$, and $`\mathrm{Hom}_{𝖣(𝒜)}(T^{},):𝖣(𝒜)^0(T^{})\mathrm{𝖬𝗈𝖽}B`$ is an equivalence.
(1) By the construction of $`X^r`$ in Proposition 1.1, $`𝖣^{}(𝒜)`$ also has a $`t`$-structure $`(𝖣^{}(𝒜)^0(T^{}),𝖣^{}(𝒜)^0(T^{}))`$ and hence by Theorem 1.3 (3) we have $`𝖣^{}(𝒜)^0(T^{})=𝖣(𝒜)^0(T^{})`$. According to \[Ri\], Proposition 10.1, we have a fully faithful $``$-functor $`F^{}:𝖣^{}(\mathrm{𝖬𝗈𝖽}B)𝖣^{}(𝒜)`$. Also, since $`F^{}(B)=T^{}`$, $`F^{}`$ sends $`B`$-modules to objects in $`𝖣(𝒜)^0(T^{})`$. Then we have a fully faithful $``$-functor
$$F:𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}B)𝖣(𝒜),$$
which sends $`B`$-modules to objects in $`𝖣(𝒜)^0(T^{})`$. For any $`X𝖣(𝒜)^\mathrm{b}(T^{})`$, there exist integers $`mn`$ such that $`X𝖣(𝒜)^{[m,n]}(T^{})`$. Let $`l=nm`$. If $`l=0`$, then there exist obviously a $`B`$-module $`M`$ and an integer $`s`$ such that $`XF(M[s])`$. If $`l>0`$, then we have a distinguished triangle
$$X^{n1}XX^n$$
with $`X^n𝖣(𝒜)^{[n]}(T^{})`$ and $`X^{n1}𝖣(𝒜)^{[m,n1]}(T^{})`$. Since $`F`$ is full, by induction on $`l`$, there exists $`U^{}𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}B)`$ such that $`XF(U^{})`$.
(2) By the assumption, $`𝖣^+(𝒜)`$ also has a $`t`$-structure $`(𝖣^+(𝒜)^0(T^{}),𝖣^+(𝒜)^0(T^{}))`$. Thus $`𝖣^+(𝒜)^\mathrm{b}(T^{})=𝖣(𝒜)^\mathrm{b}(T^{})`$, and hence $`𝖣^+(𝒜)^0(T^{})=𝖣(𝒜)^0(T^{})`$. By \[BBD\], Section 3, we have a $``$-functor $`\mathrm{real}:𝖣^\mathrm{b}(𝖣(𝒜)^0(T^{}))𝖣^+(𝒜)`$, and then we have a $``$-functor
$$F:𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}B)𝖣(𝒜),$$
which sends $`B`$-modules to objects in $`𝖣(𝒜)^0(T^{})`$. Let $`f\mathrm{Hom}_{𝖣(𝒜)}(X^{},Y^{}[n])`$ with $`X^{},Y^{}𝖣(𝒜)^0(T^{})`$ and $`n>0`$. Take a distinguished triangle in $`𝖣^+(𝒜)`$
$$X_1^{}V^{}\stackrel{𝑡}{}X^{}$$
such that $`V^{}`$ is a direct sum of copies of $`T^{}`$ and $`\mathrm{Hom}_{𝖣(𝒜)}(T^{},t)`$ is an epimorphism. By easy calculation, $`X_1^{}𝖣(𝒜)^0(T^{})`$, and hence we get an exact sequence in $`𝖣(𝒜)^0(T^{})`$
$$0X_1^{}V^{}\stackrel{𝑡}{}X^{}0.$$
Since $`\mathrm{Hom}_{𝖣(𝒜)}(T^{},Y^{}[n])=0`$, we have $`ft=0`$, i.e. $`t`$ effaces $`f`$. Thus the epimorphic version of effacibility in \[BBD\], Proposition 3.1.16 can be applied.
Finally, it is easy to see that $`𝖣(𝒜)^0(T^{})𝖣^\mathrm{b}(𝒜)`$ implies $`𝖣^\mathrm{b}(𝒜)=𝖣(𝒜)^\mathrm{b}(T^{})`$. ∎
## 2. Torsion Theories for Abelian Categories
Throughout this section, we fix the following notation. Let $`𝒜`$ be an abelian category satisfying the condition Ab4, and let $`d_P^1:P^1P^0`$ be a morphism in $`𝒜`$ with the $`P^i`$ being small projective objects of $`𝒜`$, and denote by $`P^{}`$ the mapping cone of $`d_P^1`$. We set $`𝒞(P^{})=𝖣(𝒜)^0(P^{})`$, $`B=\mathrm{Hom}_{𝖣(𝒜)}(P^{})^{\mathrm{op}}`$, and define a pair of full subcategories of $`𝒜`$
$`𝒳(P^{})`$ $`=\{X𝒜:\mathrm{Hom}_{𝖣(𝒜)}(P^{},X[1])=0\},`$
$`𝒴(P^{})`$ $`=\{X𝒜:\mathrm{Hom}_{𝖣(𝒜)}(P^{},X)=0\}.`$
For any $`X𝒜`$, we define a subobject of $`X`$
$$\tau (X)=\underset{f\mathrm{Hom}_𝒜(\mathrm{H}^0(P^{}),X)}{}\mathrm{Im}f$$
and an exact sequence in $`𝒜`$
$$(e_X):0\tau (X)\stackrel{j_X}{}X\pi (X)0.$$
###### Remark 2.1.
It is easy to see that $`P^{}`$ is a compact object of $`𝖣(𝒜)`$, and we have $`𝒳(P^{})=𝖣(𝒜)^0(P^{})𝒜`$ and $`𝒴(P^{})=𝖣(𝒜)^1(P^{})𝒜`$.
###### Lemma 2.2.
For any $`X𝒜`$, the following hold.
1. $`\mathrm{Ker}(\mathrm{Hom}_𝒜(d_P^1,X))\mathrm{Hom}_{𝖣(𝒜)}(P^{},X)`$.
2. $`\mathrm{Cok}(\mathrm{Hom}_𝒜(d_P^1,X))\mathrm{Hom}_{𝖣(𝒜)}(P^{},X[1])`$.
###### Lemma 2.3.
For any $`X𝒜`$, the following hold.
1. $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},X[n])=0`$ for $`n>1`$ and $`n<0`$.
2. $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},X)\mathrm{Hom}_𝒜(\mathrm{H}^0(P^{}),X)`$.
###### Lemma 2.4.
The following hold.
1. $`𝒳(P^{})`$ is closed under factor objects and direct sums.
2. $`𝒴(P^{})`$ is closed under subobjects.
3. For any $`X𝒜`$, $`\mathrm{Hom}_A(\mathrm{H}^0(P^{}),j_X)`$ is an isomorphism.
###### Lemma 2.5.
For any $`X^{}𝖣(𝒜)`$ and $`n`$, we have a functorial exact sequence
$$0\mathrm{Hom}_{𝖣\left(𝒜\right)}(P^{},\mathrm{H}^{n1}\left(X^{}\right)\left[1\right])\mathrm{Hom}_{𝖣\left(𝒜\right)}(P^{},X^{}\left[n\right])\mathrm{Hom}_{𝖣\left(𝒜\right)}(P^{},\mathrm{H}^n\left(X^{}\right))0.$$
Moreover, the above short exact sequence commutes with direct sums.
###### Proof.
For $`X^{}[n]𝖣(𝒜)`$, applying $`\mathrm{Hom}_{𝖣(𝒜)}(,X^{}[n])`$ to a distinguished triangle
$$P^1\stackrel{d_P^1}{}P^0P^{},$$
we have a short exact sequence
$`0\mathrm{Cok}(\mathrm{Hom}_{𝖣(𝒜)}(d_P^1,X^{}[n1]))`$ $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},X^{}[n])`$
$`\mathrm{Ker}(\mathrm{Hom}_{𝖣(𝒜)}(d_P^1,X^{}[n]))0.`$
Also, by Lemma 2.2 we get
$`\mathrm{Ker}(\mathrm{Hom}_{𝖣(𝒜)}(d_P^1,X^{}[n]))`$ $`\mathrm{Ker}(\mathrm{Hom}_𝒜(d_P^1,\mathrm{H}^n(X^{})))`$
$`\mathrm{Hom}_{𝖣(𝒜)}(P^{},\mathrm{H}^n(X^{})),`$
$`\mathrm{Cok}(\mathrm{Hom}_{𝖣(𝒜)}(d_P^1,X^{}[n1]))`$ $`\mathrm{Cok}(\mathrm{Hom}_𝒜(d_P^1,\mathrm{H}^{n1}(X^{})))`$
$`\mathrm{Hom}_{𝖣(𝒜)}(P^{},\mathrm{H}^{n1}(X^{})[1]).`$
Since the $`P^i`$ are small objects, the above short exact sequence commutes with direct sums. ∎
###### Lemma 2.6.
The following are equivalent.
1. $`\{P^{}[i]:i\}`$ is a generating set for $`𝖣(𝒜)`$.
2. $`𝒳(P^{})𝒴(P^{})=\{0\}`$.
###### Proof.
(1) $``$ (2). For any $`X𝒳(P^{})𝒴(P^{})`$, by Lemma 2.3 (1), $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},X[n])=0`$ for all $`n`$ and hence $`X=0`$.
(2) $``$ (1). Let $`X^{}𝖣(𝒜)`$ with $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},X^{}[n])=0`$ for all $`n`$. Then by Lemma 2.5, $`\mathrm{H}^n(X^{})𝒳(P^{})𝒴(P^{})=\{0\}`$. ∎
###### Lemma 2.7.
The following hold.
1. $`\mathrm{H}^0(P^{})𝒳(P^{})`$ if and only if $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},P^{}[i])=0`$ for all $`i>0`$.
2. $`\mathrm{H}^1(P^{})𝒴(P^{})`$ if and only if $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},P^{}[i])=0`$ for all $`i<0`$.
###### Proof.
By Lemma 2.5. ∎
###### Definition 2.8.
A pair $`(𝒳,𝒴)`$ of full subcategories $`𝒳,𝒴`$ in an abelian category $`𝒜`$ is called a torsion theory for $`𝒜`$ provided that the following conditions are satisfied (see e.g. \[Di\] for details):
1. $`𝒳𝒴=\{0\}`$;
2. $`𝒳`$ is closed under factor objects;
3. $`𝒴`$ is closed under subobjects;
4. for any object $`X`$ of $`𝒜`$, there exists an exact sequence $`0X^{}XX^{\prime \prime }0`$ in $`𝒜`$ with $`X^{}𝒳`$ and $`X^{\prime \prime }𝒴`$.
###### Remark 2.9.
Let $`𝒜`$ be an abelian category and $`(𝒳,𝒴)`$ a torsion theory for $`𝒜`$. Then for any $`Z𝒜`$, the following hold.
1. $`Z𝒳`$ if and only if $`\mathrm{Hom}_𝒜(Z,𝒴)=0`$.
2. $`Z𝒴`$ if and only if $`\mathrm{Hom}_𝒜(𝒳,Z)=0`$.
###### Theorem 2.10.
The following are equivalent.
1. $`\{P^{}[i]:i\}`$ is a generating set for $`𝖣(𝒜)`$ and $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},P^{}[i])=0`$ for all $`i>0`$.
2. $`𝒳(P^{})𝒴(P^{})=\{0\}`$ and $`\mathrm{H}^0(P^{})𝒳(P^{})`$.
3. $`𝒳(P^{})𝒴(P^{})=\{0\}`$ and $`\tau (X)𝒳(P^{})`$, $`\pi (X)𝒴(P^{})`$ for all $`X𝒜`$.
4. $`(𝒳(P^{}),𝒴(P^{}))`$ is a torsion theory for $`𝒜`$.
###### Proof.
(1) $``$ (2). By Lemmas 2.6 and 2.7 (1).
(2) $``$ (3). Let $`X𝒜`$. Since $`\mathrm{H}^0(P^{})𝒳(P^{})`$, it follows that $`\tau (X)𝒳(P^{})`$. Next, apply $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},)`$ to the canonical exact sequence $`(e_X)`$. It then follows by Lemmas 2.3 (2) and 2.4 (3) that $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},j_X)`$ is an isomorphism. Thus $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},\pi (X))=0`$ and hence $`\pi (X)𝒴(P^{})`$.
(3) $``$ (4). Obvious.
(4) $``$ (1). By Lemmas 2.3 (2), 2.6, 2.7 (1) and Remark 2.9 (1). ∎
###### Definition 2.11.
For a complex $`X^{}=(X^i,d^i)`$, we define the following truncations:
$`\sigma _{>n}(X^{})`$ $`:\mathrm{}0\mathrm{Im}d^nX^{n+1}X^{n+2}\mathrm{},`$
$`\sigma _n(X^{})`$ $`:\mathrm{}X^{n2}X^{n1}\mathrm{Ker}d^n0\mathrm{},`$
$`\sigma _{}^{}{}_{n}{}^{}(X^{})`$ $`:\mathrm{}0\mathrm{Cok}d^{n1}X^{n+1}X^{n+2}\mathrm{},`$
$`\sigma _{}^{}{}_{<n}{}^{}(X^{})`$ $`:\mathrm{}X^{n2}X^{n1}\mathrm{Im}d^{n1}0\mathrm{}.`$
###### Lemma 2.12.
For any $`X^{}𝖣(𝒜)`$ with $`\mathrm{H}^n(X^{})=0`$ for $`n>0`$ and $`n<1`$, there exists a distinguished triangle in $`𝖣(𝒜)`$ of the form
$$\mathrm{H}^1(X^{})[1]X^{}\mathrm{H}^0(X^{}).$$
###### Proof.
We have exact sequences in $`𝖢(𝒜)`$
$$0\sigma _1(X^{})X^{}\sigma _{>1}(X^{})0,$$
$$0\sigma _{}^{}{}_{<0}{}^{}(\sigma _{>1}(X^{}))\sigma _{>1}(X^{})\sigma _{}^{}{}_{0}{}^{}(X^{})0.$$
Also, $`\sigma _1(X^{})\mathrm{H}^1(X^{})[1]`$, $`\sigma _{}^{}{}_{<0}{}^{}(\sigma _{>1}(X^{}))0`$ and $`\sigma _{}^{}{}_{0}{}^{}(X^{})\mathrm{H}^0(X^{})`$ in $`𝖣(𝒜)`$. Thus we get a desired distinguished triangle in $`𝖣(𝒜)`$. ∎
###### Lemma 2.13.
Assume $`𝒳(P^{})𝒴(P^{})=\{0\}`$. Then for any $`X^{}𝖣(𝒜)`$, the following are equivalent.
1. $`X^{}𝒞(P^{})`$.
2. $`\mathrm{H}^n(X^{})=0`$ for $`n>0`$ and $`n<1`$, $`\mathrm{H}^0(X^{})𝒳(P^{})`$ and $`\mathrm{H}^1(X^{})𝒴(P^{})`$.
###### Proof.
By Lemma 2.5. ∎
###### Remark 2.14.
Let $`𝒜`$ be an abelian category and $`𝒳,𝒴`$ full subcategories of $`𝒜`$. Then the pair $`(𝒳,𝒴)`$ is a torsion theory for $`𝒜`$ if and only if the following two conditions are satisfied:
1. $`\mathrm{Hom}_𝒜(𝒳,𝒴)=0`$;
2. for any object $`X`$ in $`𝒜`$, there exists an exact sequence $`0X^{}XX^{\prime \prime }0`$ in $`𝒜`$ with $`X^{}𝒳`$ and $`X^{\prime \prime }𝒴`$.
###### Theorem 2.15.
Assume $`𝒳(P^{})𝒴(P^{})=\{0\}`$ and $`\mathrm{H}^0(P^{})𝒳(P^{})`$. Then the following hold.
1. $`𝒞(P^{})`$ is admissible abelian.
2. The functor
$$\mathrm{Hom}_{𝖣(𝒜)}(P^{},):𝒞(P^{})\mathrm{𝖬𝗈𝖽}B$$
is an equivalence.
3. $`(𝒴(P^{})[1],𝒳(P^{}))`$ is a torsion theory for $`𝒞(P^{})`$.
###### Proof.
(1) and (2) According to Theorem 2.10, Theorem 1.3 can be applied.
(3) Note first that by Lemma 2.13 we have $`𝒳(P^{})𝒞(P^{})`$ and $`𝒴(P^{})[1]𝒞(P^{})`$. Also, it is trivial that $`\mathrm{Hom}_{𝖣(𝒜)}(𝒴(P^{})[1],𝒳(P^{}))=0`$. Let $`X^{}𝒞(P^{})`$. Then by Lemmas 2.12 and 2.13 we have a distinguished triangle in $`𝖣(𝒜)`$ of the form
$$\mathrm{H}^1(X^{})[1]X^{}\mathrm{H}^0(X^{}).$$
It follows that the sequence in $`𝒞(P^{})`$
$$0\mathrm{H}^1(X^{})[1]X^{}\mathrm{H}^0(X^{})0$$
is exact. Thus by Remark 2.14 $`(𝒴(P^{})[1],𝒳(P^{}))`$ is a torsion theory for $`𝒞(P^{})`$. ∎
###### Proposition 2.16.
Assume $`P^{}`$ satisfies the conditions
1. $`\{P^{}[i]:i\}`$ is a generating set for $`𝖣(𝒜)`$,
2. $`\mathrm{Hom}_{𝖣(𝒜)}(P^{},P^{}[i])=0`$ for $`i0`$.
If $`𝒜`$ has either enough projectives or enough injectives, then we have an equivalence of triangulated categories
$$𝖣^\mathrm{b}(𝒜)𝖣^\mathrm{b}(\mathrm{𝖬𝗈𝖽}B).$$
###### Proof.
Let $`X^{}𝖣(𝒜)`$. According to Lemma 2.5 and Theorem 2.10, it is easy to see that if $`X^{}`$ belongs to $`𝖣(𝒜)^0(P^{})`$ (resp., $`𝒞(P^{})`$), then $`\mathrm{H}^n(X^{})=0`$ for $`n<1`$ (resp., $`n<1`$ and $`n>0`$). Thus we have
$$𝖣(𝒜)^0(P^{})𝖣^+(𝒜)\text{and}𝒞(P^{})𝖣^\mathrm{b}(𝒜),$$
so that Corollary 1.5 can be applied. ∎
## 3. Torsion Theories for Module Categories
In this section, we apply results of Section 2 to the case of module categories. In and after this section, $`R`$ is a commutative ring and $`I`$ is an injective cogenerator in the category of $`R`$-modules. We set $`D=\mathrm{Hom}_R(,I)`$. Let $`A`$ be an $`R`$-algebra and denote by $`\mathrm{𝖯𝗋𝗈𝗃}A`$ (resp., $`\mathrm{𝗉𝗋𝗈𝗃}A`$) the full additive subcategory of $`\mathrm{𝖬𝗈𝖽}A`$ consisting of projective (resp., finitely generated projective) modules. We denote by $`A^{\mathrm{op}}`$ the opposite ring of $`A`$ and consider right $`A`$-modules as left $`A^{\mathrm{op}}`$-modules. Also, we denote by $`()^{}`$ both the $`A`$-dual functors $`\mathrm{Hom}_A(,A)`$ and set $`\nu =D()^{}`$.
It is well known that, in a module category, the small projective objects are just the finitely generated projective modules. In the following, we deal with the case where $`𝒜=\mathrm{𝖬𝗈𝖽}A`$ and use the same notation as in Section 2.
###### Lemma 3.1.
For any $`X\mathrm{𝖬𝗈𝖽}A`$, we have
$$\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(P^{},X[1])\mathrm{H}^1((P^{})^{})_AX.$$
###### Proof.
We have
$`\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(P^{},X[1])`$ $`\mathrm{Hom}_{𝖪(\mathrm{𝖬𝗈𝖽}A)}(P^{},X[1])`$
$`\mathrm{H}^1(\mathrm{Hom}_A^{}(P^{},X))`$
$`\mathrm{H}^1((P^{})^{}_A^{}X)`$
$`\mathrm{H}^1((P^{})^{})_AX.`$
###### Lemma 3.2.
The following hold.
1. $`𝒳(P^{})=\mathrm{Ker}(\mathrm{H}^1((P^{})^{})_A)`$.
2. $`𝒴(P^{})=\mathrm{Ker}(\mathrm{Hom}_A(\mathrm{H}^0(P^{}),))`$.
###### Proof.
By Lemmas 2.3 (2) and 3.1. ∎
###### Lemma 3.3.
The following hold.
1. $`D(\mathrm{H}^1((P^{})^{}))\mathrm{H}^1(\nu (P^{}))`$.
2. $`𝒳(P^{})=\mathrm{Ker}(\mathrm{Hom}_A(,\mathrm{H}^1(\nu (P^{}))))`$ and hence $`\mathrm{H}^0(P^{})𝒳(P^{})`$ if and only if $`\mathrm{H}^1(\nu (P^{}))𝒴(P^{})`$.
3. $`\mathrm{Ker}(\mathrm{Tor}_1^A(\mathrm{H}^1((P^{})^{}),))=\mathrm{Ker}(\mathrm{Ext}_A^1(,\mathrm{H}^1(\nu (P^{}))))`$.
###### Proof.
We have $`D(\mathrm{H}^1((P^{})^{}))\mathrm{H}^1(D((P^{})^{}))=\mathrm{H}^1(\nu (P^{}))`$ and for any $`X\mathrm{𝖬𝗈𝖽}A`$ we have
$`D(\mathrm{H}^1((P^{})^{})_AX)`$ $`\mathrm{Hom}_A(X,\mathrm{H}^1(\nu (P^{}))),`$
$`D(\mathrm{Tor}_1^A((\mathrm{H}^1((P^{})^{}),X)))`$ $`\mathrm{Ext}_A^1(X,\mathrm{H}^1(\nu (P^{}))).`$
###### Lemma 3.4.
The following hold.
1. $`𝒳(P^{})\mathrm{Ker}(\mathrm{Ext}_A^1(\mathrm{H}^0(P^{}),))`$.
2. $`𝒴(P^{})\mathrm{Ker}(\mathrm{Tor}_1^A(\mathrm{H}^1((P^{})^{}),))`$.
###### Proof.
This is due essentially to Auslander \[Au\]. We have an exact sequence in $`\mathrm{𝖬𝗈𝖽}A`$
$$0\mathrm{H}^1(P^{})P^1P^0\mathrm{H}^0(P^{})0$$
with the $`P^i`$ finitely generated projective, and an exact sequence in $`\mathrm{𝖬𝗈𝖽}A^{\mathrm{op}}`$
$$0\mathrm{H}^0(P^{})^{}P^0P^1\mathrm{H}^1((P^{})^{})0$$
with the $`P^i`$ finitely generated projective.
(1) Let $`X\mathrm{𝖬𝗈𝖽}A`$. For any $`M\mathrm{𝖬𝗈𝖽}A^{\mathrm{op}}`$, we have a functorial homomorphism
$$\theta _M:M_AX\mathrm{Hom}_A(M^{},X),mx(hh(m)x)$$
which is an isomorphism if $`M`$ is finitely generated projective. Since the $`P^i`$ are reflexive, we have $`\mathrm{H}^0(P^{})\mathrm{H}^0((P^{})^{})`$ and $`\mathrm{H}^1(P^{})\mathrm{H}^1((P^{})^{})^{}`$. We have a commutative diagram
$$\begin{array}{ccccccc}P^0_AX& & P^1_AX& & \mathrm{H}^1((P^{})^{})_AX& & 0\\ \theta _{P^0}& & \theta _{P^1}& & & & \\ \mathrm{Hom}_A(P^0,X)& & \mathrm{Hom}_A(P^1,X)& & \mathrm{Hom}_A(\mathrm{H}^1(P^{}),X)\end{array}$$
with the top row exact. Since the $`\theta _{P^i}`$ are isomorphisms, $`\mathrm{Ext}_A^1(\mathrm{H}^0(P^{}),X)`$ is embedded in $`\mathrm{H}^1((P^{})^{})_AX`$. The assertion follows by Lemma 3.2.
(2) Let $`X\mathrm{𝖬𝗈𝖽}A`$. For any $`Y\mathrm{𝖬𝗈𝖽}A`$, we have a functorial homomorphism
$$\eta _Y:Y^{}_AX\mathrm{Hom}_A(Y,X),hx(yh(y)x)$$
which is an isomorphism if $`Y`$ is finitely generated projective. We have a commutative diagram
$$\begin{array}{ccccccc}& & \mathrm{H}^0(P^{})^{}_AX& & P^0_AX& & P^1_AX\\ & & & & \eta _{P^0}& & \eta _{P^1}& & \\ 0& & \mathrm{Hom}_A(\mathrm{H}^0(P^{}),X)& & \mathrm{Hom}_A(P^0,X)& & \mathrm{Hom}_A(P^1,X)\end{array}$$
with the bottom row exact. Since the $`\eta _{P^i}`$ are isomorphisms, $`\mathrm{Tor}_1^A(\mathrm{H}^1((P^{})^{}),X)`$ is a homomorphic image of $`\mathrm{Hom}_A(\mathrm{H}^0(P^{}),X)`$. The assertion follows by Lemma 3.2. ∎
###### Theorem 3.5.
The following are equivalent.
1. $`𝒳(P^{})𝒴(P^{})=\{0\}`$ and $`\mathrm{H}^0(P^{})𝒳(P^{})`$.
2. $`𝒳(P^{})𝒴(P^{})=\{0\}`$ and $`\tau (X)𝒳(P^{})`$, $`\pi (X)𝒴(P^{})`$ for all $`X\mathrm{𝖬𝗈𝖽}A`$.
3. $`(𝒳(P^{}),𝒴(P^{}))`$ is a torsion theory for $`\mathrm{𝖬𝗈𝖽}A`$.
4. $`𝒳(P^{})`$ consists of the modules generated by $`\mathrm{H}^0(P^{})`$ and $`𝒴(P^{})`$ consists of the modules cogenerated by $`\mathrm{H}^1(\nu (P^{}))`$.
###### Proof.
(1) $``$ (2) $``$ (3). By Theorem 2.10.
(3) $``$ (4). Since $`\mathrm{Hom}_A(\mathrm{H}^0(P^{}),)`$ vanishes on $`𝒴(P^{})`$, $`\mathrm{H}^0(P^{})𝒳(P^{})`$. Thus $`𝒳(P^{})`$ contains the modules generated by $`\mathrm{H}^0(P^{})`$. Conversely, let $`X𝒳(P^{})`$. Then, since (1) implies (2), $`\pi (X)𝒴(P^{})`$ and hence $`\mathrm{Hom}_A(X,\pi (X))=0`$. Thus $`X=\tau (X)`$, which is generated by $`\mathrm{H}^0(P^{})`$. Next, since by Lemma 3.3 (2) $`\mathrm{H}^1(\nu (P^{}))𝒴(P^{})`$, $`𝒴(P^{})`$ contains the modules cogenerated by $`\mathrm{H}^1(\nu (P^{}))`$. Conversely, let $`X𝒴(P^{})`$. Take a set of generators $`\{f_\lambda \}_{\lambda \mathrm{\Lambda }}`$ for an $`R`$-module $`\mathrm{Hom}_A(X,\mathrm{H}^1(\nu (P^{})))`$ and set
$$f:X\mathrm{H}^1(\nu (P^{}))^\mathrm{\Lambda },x(f_\lambda (x))_{\lambda \mathrm{\Lambda }}.$$
It is obvious that $`\mathrm{Hom}_A(f,\mathrm{H}^1(\nu (P^{})))`$ is surjective. Also, by Lemmas 3.3 (3) and 3.4(2) we have $`\mathrm{Ext}_A^1(\mathrm{Im}f,\mathrm{H}^1(\nu (P^{})))=0`$. Applying $`\mathrm{Hom}_A(,\mathrm{H}^1(\nu (P^{})))`$ to the canonical exact sequence
$$0\mathrm{Ker}fX\mathrm{Im}f0,$$
we get $`\mathrm{Hom}_A(\mathrm{Ker}f,\mathrm{H}^1(\nu (P^{})))=0`$. Thus $`\mathrm{Ker}f𝒳(P^{})𝒴(P^{})`$ and hence $`\mathrm{Ker}f=0`$.
(4) $``$ (1). By Lemma 3.3 (2). ∎
###### Corollary 3.6.
The following are equivalent.
1. $`P^{}`$ is a tilting complex.
2. $`𝒳(P^{})𝒴(P^{})=\{0\},\mathrm{H}^0(P^{})𝒳(P^{})`$ and $`\mathrm{H}^1(P^{})𝒴(P^{})`$.
3. $`(𝒳(P^{}),𝒴(P^{}))`$ is a torsion theory for $`\mathrm{𝖬𝗈𝖽}A`$ and $`\mathrm{H}^1(P^{})𝒴(P^{})`$.
###### Proof.
By Lemmas 2.6, 2.7 and Theorem 3.5. ∎
For an object $`X`$ in an additive category $``$, we denote by $`\mathrm{𝖺𝖽𝖽}(X)`$ the full subcategory of $``$ consisting of objects which are direct summands of finite direct sums of copies of $`X`$.
###### Corollary 3.7.
For any tilting complexes $`P_1^{}:P_1^1P_1^0`$, $`P_2^{}:P_2^1P_2^0`$ for $`A`$ of term length two, the following are equivalent.
1. $`(𝒳(P_1^{}),𝒴(P_1^{}))=(𝒳(P_2^{}),𝒴(P_2^{}))`$.
2. $`\mathrm{𝖺𝖽𝖽}(P_1^{})=\mathrm{𝖺𝖽𝖽}(P_2^{})`$ in $`𝖪^\mathrm{b}(\mathrm{𝖯𝗋𝗈𝗃}A)`$.
###### Proof.
(1) $``$ (2). It follows by Corollary 3.6 that $`Q^{}=P_1^{}P_2^{}`$ is a tilting complex such that $`(𝒳(Q^{}),𝒴(Q^{}))=(𝒳(P_i^{}),𝒴(P_i^{}))`$ ($`i=1,2`$). Let $`B=\mathrm{End}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(Q^{})^{\mathrm{op}}`$ and for $`i=1,2`$ denote by $`e_i`$ the composite of canonical homomorphisms $`Q^{}P_i^{}Q^{}`$ . Then for $`i=1,2`$ we have an equivalence $`𝖣^{}(\mathrm{𝖬𝗈𝖽}B)𝖣^{}(\mathrm{𝖬𝗈𝖽}e_iBe_i)`$ which sends $`Be_i`$ to $`e_iBe_i`$, so that the $`Be_i`$ are tilting complexes for $`B`$, i.e. projective generators for $`\mathrm{𝖬𝗈𝖽}B`$. It follows by Morita Theory that $`\mathrm{𝖺𝖽𝖽}B=\mathrm{𝖺𝖽𝖽}Be_i`$ in $`\mathrm{𝖬𝗈𝖽}B`$. Thus $`\mathrm{𝖺𝖽𝖽}(P_1^{})=\mathrm{𝖺𝖽𝖽}(P_2^{})`$ in $`𝖪^\mathrm{b}(\mathrm{𝖯𝗋𝗈𝗃}A)`$.
(2) $``$ (1). It is obviously deduced that $`\mathrm{𝖺𝖽𝖽}(\mathrm{H}^1(\nu (P_1^{})))=\mathrm{𝖺𝖽𝖽}(\mathrm{H}^1(\nu (P_2^{})))`$ and $`\mathrm{𝖺𝖽𝖽}(\mathrm{H}^0(P_1^{}))=\mathrm{𝖺𝖽𝖽}(\mathrm{H}^0(P_2^{}))`$. ∎
###### Theorem 3.8.
Assume $`𝒳(P^{})𝒴(P^{})=\{0\}`$ and $`\mathrm{H}^0(P^{})𝒳(P^{})`$. Then the following hold.
1. $`\{P^{}[i]:i\}`$ is a generating set for $`𝖣(\mathrm{𝖬𝗈𝖽}A)`$.
2. $`𝒞(P^{})`$ is admissible abelian.
3. $`(𝒴(P^{})[1],𝒳(P^{}))`$ is a torsion theory for $`𝒞(P^{})`$.
4. The functor
$$\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(P^{},):𝒞(P^{})\mathrm{𝖬𝗈𝖽}B$$
is an equivalence.
###### Proof.
By Lemma 2.6 and Theorem 2.15. ∎
###### Remark 3.9.
The following are equivalent.
1. $`P^{}`$ is a tilting complex.
2. $`𝒳(P^{})𝒴(P^{})=\{0\}`$ and $`P^{}𝒞(P^{})`$.
###### Example 3.10 (cf. \[HK\]).
Let $`A`$ be a finite dimensional algebra over a field $`k`$ given by a quiver
$$\begin{array}{ccc}1& \stackrel{\alpha }{}& 2\\ \delta & & \beta & & \\ 4& \underset{\gamma }{}& 3\end{array}$$
with relations $`\beta \alpha =\gamma \beta =\delta \gamma =\alpha \delta =0`$. For each vertex $`i`$, we denote by $`S(i),P(i)`$ the corresponding simple and indecomposable projective left $`A`$-modules, respectively. Define a complex $`P^{}`$ as the mapping cone of the homomorphism
$$d_P^1=\left[\begin{array}{cccc}f& 0& 0& 0\\ 0& 0& g& 0\end{array}\right]:P(2)^2P(4)^2P(1)P(3),$$
where $`f`$ and $`g`$ denote the right multiplications of $`\alpha `$ and $`\gamma `$, respectively. Then $`P^{}`$ is not a tilting complex. However, $`P^{}`$ satisfies the assumption of Theorem 3.8 and hence we have an equivalence of abelian categories
$$\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(P^{},):𝒞(P^{})\mathrm{𝖬𝗈𝖽}B,$$
where $`B=\mathrm{End}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(P^{})^{\mathrm{op}}`$ is a finite dimensional $`k`$-algebra given by a quiver
$$1234.$$
There exist exact sequences in $`𝒞(P^{})`$ of the form
$$0S(1)S(2)[1]P(1)[1]0,0S(3)S(4)[1]P(3)[1]0,$$
and these objects and morphisms generate $`𝒞(P^{})`$.
## 4. Equivalences between Torsion Theories
Throughout this section, $`P^{}`$ is assumed to be a tilting complex. Then there exists an equivalence of triangulated categories
$$F:𝖣^{}(\mathrm{𝖬𝗈𝖽}B)𝖣^{}(\mathrm{𝖬𝗈𝖽}A)$$
such that $`F(B)=P^{}`$. Let $`G:𝖣^{}(\mathrm{𝖬𝗈𝖽}A)𝖣^{}(\mathrm{𝖬𝗈𝖽}B)`$ be a quasi-inverse of $`F`$. For any $`n`$, we have ring homomorphisms
$$B\mathrm{End}_A(\mathrm{H}^n(P^{}))^{\mathrm{op}}\text{and}B\mathrm{End}_A(\mathrm{H}^n((P^{})^{})).$$
In particular, $`\mathrm{H}^0(P^{})`$ is an $`A`$-$`B`$-bimodule and $`\mathrm{H}^1((P^{})^{})`$ is a $`B`$-$`A`$-bimodule.
###### Lemma 4.1.
The following hold.
1. For any $`X^{}𝒞(P^{})`$, we have $`G(X^{})\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(P^{},X^{})`$.
2. We have an equivalence
$$\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(P^{},):𝒞(P^{})\mathrm{𝖬𝗈𝖽}B$$
whose quasi-inverse is given by the restriction of $`F`$ to $`\mathrm{𝖬𝗈𝖽}B`$.
###### Proof.
See \[Ri\], Section 4. ∎
###### Lemma 4.2.
There exists a tilting complex $`Q^{}𝖪^\mathrm{b}(\mathrm{𝗉𝗋𝗈𝗃}B)`$ such that
1. $`Q^{}G(A)`$,
2. $`Q^i=0`$ for $`i>1`$ and $`i<0`$,
3. $`\mathrm{H}^i(Q^{})\mathrm{H}^i((P^{})^{})`$ for $`0i1`$,
4. $`\mathrm{H}^i(\mathrm{Hom}_B^{}(Q^{},B))\mathrm{H}^i(P^{})`$ for $`1i0`$.
###### Proof.
By \[Ri\], Proposition 6.3, there exists $`Q^{}𝖪^\mathrm{b}(\mathrm{𝗉𝗋𝗈𝗃}B)`$ satisfying $`Q^{}G(A)`$. Since
$`\mathrm{H}^i(Q^{})`$ $`\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}B)}(B,Q^{}[i])`$
$`\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(P^{},A[i])`$
$`\mathrm{H}^i((P^{})^{}),`$
we have $`Q^{}\sigma _1(Q^{})`$ in $`𝖪^\mathrm{b}(\mathrm{𝗉𝗋𝗈𝗃}B)`$. Also, since
$`\mathrm{H}^i(\mathrm{Hom}_B^{}(Q^{},B))`$ $`\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}B)}(Q^{},B[i])`$
$`\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(A,P^{}[i])`$
$`\mathrm{H}^i(P^{}),`$
we have $`\mathrm{Hom}_B^{}(Q^{},B)\sigma _0(\mathrm{Hom}_B^{}(Q^{},B))`$ in $`𝖪^\mathrm{b}(\mathrm{𝗉𝗋𝗈𝗃}B^{\mathrm{op}})`$ and $`Q^{}\sigma _{}^{}{}_{0}{}^{}(Q^{})`$ in $`𝖪^\mathrm{b}(\mathrm{𝗉𝗋𝗈𝗃}B)`$. Thus, we can assume $`Q^i=0`$ for $`i>1`$ and $`i<0`$. ∎
###### Lemma 4.3.
For any $`M\mathrm{𝖬𝗈𝖽}B`$, the following hold.
1. $`\mathrm{H}^i(F(M))=0`$ for $`i>0`$ and $`i<1`$.
2. $`\mathrm{H}^0(F(M))\mathrm{H}^0(P^{})_BM`$.
3. $`\mathrm{H}^1(F(M))\mathrm{Hom}_B(\mathrm{H}^1((P^{})^{}),M)`$.
###### Proof.
For any $`i`$, we have
$`\mathrm{H}^i(F(M))`$ $`\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}A)}(A,F(M)[i])`$
$`\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}B)}(Q^{},M[i]).`$
Thus $`\mathrm{H}^i(F(M))=0`$ for $`i>0`$ and $`i<1`$. Also,
$`\mathrm{H}^0(F(M))`$ $`\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}B)}(Q^{},M)`$
$`\mathrm{H}^0(\mathrm{Hom}_B^{}(Q^{},M))`$
$`\mathrm{H}^0(\mathrm{Hom}_B^{}(Q^{},B)_BM)`$
$`\mathrm{H}^0(\mathrm{Hom}_B^{}(Q^{},B))_BM`$
$`\mathrm{H}^0(P^{})_BM,`$
$`\mathrm{H}^1(F(M))`$ $`\mathrm{Hom}_{𝖣(\mathrm{𝖬𝗈𝖽}B)}(Q^{},M[1])`$
$`\mathrm{H}^1(\mathrm{Hom}_B^{}(Q^{},M))`$
$`\mathrm{Hom}_B(\mathrm{H}^1(Q^{}),M)`$
$`\mathrm{Hom}_B(H^1((P^{})^{}),M).`$
###### Theorem 4.4.
Define a pair of full subcategories of $`\mathrm{𝖬𝗈𝖽}B`$
$$𝒰(P^{})=\mathrm{Ker}(\mathrm{H}^0(P^{})_B),𝒱(P^{})=\mathrm{Ker}(\mathrm{Hom}_B(\mathrm{H}^1((P^{})^{}),)).$$
Then the following hold.
1. $`(𝒰(P^{}),𝒱(P^{}))`$ is a torsion theory for $`\mathrm{𝖬𝗈𝖽}B`$.
2. We have a pair of functors
$$\mathrm{Hom}_A(\mathrm{H}^0(P^{}),):𝒳(P^{})𝒱(P^{}),\mathrm{H}^0(P^{})_B:𝒱(P^{})𝒳(P^{})$$
which define an equivalence.
3. We have a pair of functors
$$\mathrm{H}^1\left(\left(P^{}\right)^{}\right)_A:𝒴\left(P^{}\right)𝒰\left(P^{}\right),\mathrm{Hom}_B(\mathrm{H}^1\left(\left(P^{}\right)^{}\right),):𝒰\left(P^{}\right)𝒴\left(P^{}\right)$$
which define an equivalence.
###### Proof.
(1) According to Lemmas 3.2 and 4.2, we can apply Corollary 3.6 for a tilting complex $`Q^{}`$ to conclude that $`(𝒰(P^{}),𝒱(P^{}))`$ is a torsion theory for $`\mathrm{𝖬𝗈𝖽}B`$.
(2) For any $`X𝒳(P^{})`$, by Lemmas 2.13, 4.1 (1) and 4.3 (3) we have
$`\mathrm{Hom}_B(\mathrm{H}^1((P^{})^{}),\mathrm{Hom}_A(\mathrm{H}^0(P^{}),X))`$ $`\mathrm{H}^1(F(G(X)))`$
$`\mathrm{H}^1(X)`$
$`=0.`$
Also, since by Lemma 3.2 (1) and Corollary 3.6 $`\mathrm{H}^1((P^{})^{})_A\mathrm{H}^0(P^{})=0`$,
$`\mathrm{H}^1((P^{})^{})_A\mathrm{H}^0(P^{})_BM=0`$ for all $`M𝒱(P^{})`$. The last assertion follows by Lemmas 2.13, 4.1 and 4.3.
(3) For any $`Y𝒴(P^{})`$, by Lemmas 2.13, 3.1, 4.1 (1) and 4.3 (2) we have
$`\mathrm{H}^0(P^{})_B\mathrm{H}^1((P^{})^{})_AY`$ $`\mathrm{H}^0(F(G(Y[1])))`$
$`\mathrm{H}^0(Y[1])`$
$`=0.`$
Also, since $`\mathrm{H}^1((P^{})^{})_A\mathrm{H}^0(P^{})=0`$, for any $`N𝒰(P^{})`$ we have
$`\mathrm{Hom}_A(\mathrm{H}^0(P^{}),\mathrm{Hom}_B(\mathrm{H}^1((P^{})^{}),N))`$ $`\mathrm{Hom}_B(\mathrm{H}^1((P^{})^{})_A\mathrm{H}^0(P^{}),N)`$
$`=0.`$
The last assertion follows by Lemmas 2.13, 4.1 and 4.3. ∎
###### Definition 4.5.
Let $`(𝒰,𝒱)`$ be a torsion theory for an abelian category $`𝒜`$. Then $`(𝒰,𝒱)`$ is called splitting if $`\mathrm{Ext}_𝒜^1(𝒱,𝒰)=0`$.
For a left $`A`$-module $`M`$, we denote by $`\mathrm{proj}\mathrm{dim}{}_{A}{}^{}M`$ (resp., $`\mathrm{inj}\mathrm{dim}{}_{A}{}^{}M`$) the projective (resp., the injective) dimension of $`M`$.
###### Proposition 4.6.
The torsion theory $`(𝒰(P^{}),𝒱(P^{}))`$ for $`\mathrm{𝖬𝗈𝖽}B`$ is splitting if and only if $`\mathrm{Ext}_A^2(𝒳(P^{}),𝒴(P^{}))=0`$. In particular, $`(𝒰(P^{}),𝒱(P^{}))`$ is splitting if either $`\mathrm{proj}\mathrm{dim}X1`$ for all $`X𝒳(P^{})`$ or $`\mathrm{inj}\mathrm{dim}Y1`$ for all $`Y𝒴(P^{})`$.
###### Proof.
For any $`X𝒳(P^{})`$ and $`Y𝒴(P^{})`$, we have
$`\mathrm{Ext}_B^1(\mathrm{Hom}_A(\mathrm{H}^0\left(P^{}\right),X),\mathrm{H}^1\left(\left(P^{}\right)^{}\right)_AY)`$ $`\mathrm{Hom}_{𝖣\left(\mathrm{𝖬𝗈𝖽}B\right)}(G\left(X\right),G\left(Y\left[1\right]\right)\left[1\right])`$
$`\mathrm{Hom}_{𝖣\left(\mathrm{𝖬𝗈𝖽}A\right)}(X,Y\left[2\right])`$
$`\mathrm{Ext}_A^2(X,Y).`$
## 5. Torsion Theories for Artin Algebras
In this section, we deal with the case where $`R`$ is a commutative artin ring, $`I`$ is an injective envelope of an $`R`$-module $`R/\mathrm{rad}(R)`$ and $`A`$ is a finitely generated $`R`$-module. We denote by $`\mathrm{𝗆𝗈𝖽}A`$ the full abelian subcategory of $`\mathrm{𝖬𝗈𝖽}A`$ consisting of finitely generated modules. Note that $`\mathrm{H}^n(P^{}),\mathrm{H}^n(\nu (P^{}))\mathrm{𝗆𝗈𝖽}A`$ for all $`n`$. We set
$$𝒳_c(P^{})=𝒳(P^{})\mathrm{𝗆𝗈𝖽}A\text{and}𝒴_c(P^{})=𝒴(P^{})\mathrm{𝗆𝗈𝖽}A.$$
###### Proposition 5.1.
The following are equivalent.
1. $`𝒳_c(P^{})𝒴_c(P^{})=\{0\}`$ and $`\mathrm{H}^0(P^{})𝒳_c(P^{})`$.
2. $`𝒳_c(P^{})𝒴_c(P^{})=\{0\}`$ and $`\tau (X)𝒳_c(P^{})`$, $`\pi (X)𝒴_c(P^{})`$ for all $`X\mathrm{𝗆𝗈𝖽}A`$.
3. $`(𝒳_c(P^{}),𝒴_c(P^{}))`$ is a torsion theory for $`\mathrm{𝗆𝗈𝖽}A`$.
4. $`𝒳_c(P^{})`$ consists of the modules generated by $`\mathrm{H}^0(P^{})`$ and $`𝒴_c(P^{})`$ consists of the modules cogenerated by $`\mathrm{H}^1(\nu (P^{}))`$.
###### Proof.
By the same arguments as in the proof of Theorem 3.5. ∎
###### Lemma 5.2.
The following are equivalent.
1. $`\{P^{}[i]:i\}`$ is a generating set for $`𝖣(\mathrm{𝗆𝗈𝖽}A)`$.
2. $`𝒳_c(P^{})𝒴_c(P^{})=\{0\}`$.
###### Proof.
By the same arguments as in the proof of Lemma 2.6. ∎
###### Lemma 5.3.
The following hold.
1. If $`DA𝒳_c(P^{})`$, then $`\mathrm{H}^1(P^{})=0`$, i.e. $`P^{}\mathrm{H}^0(P^{})`$ in $`𝖣(\mathrm{𝗆𝗈𝖽}A)`$.
2. $`\mathrm{H}^0(\nu (P^{}))𝒳_c(P^{})`$ if and only if $`\mathrm{H}^1(P^{})𝒴_c(P^{})`$.
###### Proof.
For any $`P\mathrm{𝗉𝗋𝗈𝗃}A`$, we have functorial isomorphisms
$$\nu (P)DA_AP\text{and}P\mathrm{Hom}_A(DA,\nu (P)).$$
Thus
$$\mathrm{H}^0(\nu (P^{}))DA_A\mathrm{H}^0(P^{})\text{and}\mathrm{H}^1(P^{})\mathrm{Hom}_A(DA,\mathrm{H}^1(\nu (P^{})))$$
and hence
$`\mathrm{Hom}_A(\mathrm{H}^0(\nu (P^{})),\mathrm{H}^1(\nu (P^{})))`$ $`\mathrm{Hom}_A(DA_A\mathrm{H}^0(P^{}),\mathrm{H}^1(\nu (P^{})))`$
$`\mathrm{Hom}_A(\mathrm{H}^0(P^{}),\mathrm{Hom}_A(DA,\mathrm{H}^1(\nu (P^{}))))`$
$`\mathrm{Hom}_A(\mathrm{H}^0(P^{}),\mathrm{H}^1(P^{})).`$
###### Lemma 5.4.
Assume $`𝒳_c(P^{})𝒴_c(P^{})=\{0\}`$ and $`\mathrm{H}^0(P^{})𝒳_c(P^{})`$. Then the following are equivalent.
1. $`\mathrm{H}^0(\nu (P^{}))𝒳_c(P^{})`$.
2. $`𝒳_c(P^{})`$ is stable under $`DA_A`$.
3. $`\mathrm{H}^1(P^{})𝒴_c(P^{})`$.
4. $`𝒴_c(P^{})`$ is stable under $`\mathrm{Hom}_A(DA,)`$.
###### Proof.
(1) $``$ (2). Let $`X𝒳_c(P^{})`$. Then by Proposition 5.1 $`X`$ is generated by $`\mathrm{H}^0(P^{})`$ and hence $`DA_AX`$ is generated by $`DA_A\mathrm{H}^0(P^{})\mathrm{H}^0(\nu (P^{}))𝒳_c(P^{})`$.
(2) $``$ (3). Since $`\mathrm{H}^0(\nu (P^{}))DA_A\mathrm{H}^0(P^{})𝒳_c(P^{})`$, by Lemma 5.3 (2) we have $`\mathrm{H}^1(P^{})𝒴_c(P^{})`$.
(3) $``$ (4) $``$ (1). By the dual arguments. ∎
###### Proposition 5.5.
The following are equivalent.
1. $`P^{}`$ is a tilting complex.
2. $`𝒳_c(P^{})𝒴_c(P^{})=\{0\}`$, $`\mathrm{H}^0(P^{})𝒳_c(P^{})`$ and $`\mathrm{H}^1(P^{})𝒴_c(P^{})`$.
3. $`(𝒳_c(P^{}),𝒴_c(P^{}))`$ is a torsion theory for $`\mathrm{𝗆𝗈𝖽}A`$ and $`\mathrm{H}^1(P^{})`$ $`𝒴_c(P^{})`$.
4. $`(𝒳_c(P^{}),𝒴_c(P^{}))`$ is a torsion theory for $`\mathrm{𝗆𝗈𝖽}A`$ and $`𝒳_c(P^{})`$ is stable under $`DA_A`$.
5. $`(𝒳_c(P^{}),𝒴_c(P^{}))`$ is a torsion theory for $`\mathrm{𝗆𝗈𝖽}A`$ and $`𝒴_c(P^{})`$ is stable under $`\mathrm{Hom}_A(DA,)`$.
###### Proof.
By Proposition 5.1, Lemmas 2.7, 5.2 and 5.4. ∎
###### Definition 5.6.
Let $`𝒜`$ be an abelian category and $`𝒞`$ a full subcategory of $`𝒜`$ closed under extensions. Then an object $`X𝒞`$ is called $`\mathrm{Ext}`$-projective (resp., $`\mathrm{Ext}`$-injective) if $`\mathrm{Ext}_𝒜^1(X,𝒞)=0`$ (resp., $`\mathrm{Ext}_𝒜^1(𝒞,X)=0`$).
###### Proposition 5.7.
Assume $`P^{}`$ is a tilting complex. Then the following hold.
1. $`\mathrm{H}^0(P^{})𝒳_c(P^{})`$ is $`\mathrm{Ext}`$-projective and generates $`𝒳_c(P^{})`$.
2. $`\mathrm{H}^1(\nu (P^{}))𝒴_c(P^{})`$ is $`\mathrm{Ext}`$-injective and cogenerates $`𝒴_c(P^{})`$.
###### Proof.
By Propositions 5.1, 5.5 and Lemmas 3.3, 3.4. ∎
###### Theorem 5.8.
Let $`(𝒳,𝒴)`$ be a torsion theory for $`\mathrm{𝗆𝗈𝖽}A`$ such that $`𝒳`$ contains an $`\mathrm{Ext}`$-projective module X which generates $`𝒳`$, $`𝒴`$ contains an $`\mathrm{Ext}`$-injective module Y which cogenerates $`𝒴`$, and $`𝒳`$ is stable under $`DA_A`$. Let $`M_X^{}`$ be a minimal projective presentation of $`X`$ and $`N_Y^{}`$ a minimal injective presentation of $`Y`$. Then
$$P^{}=M_X^{}\mathrm{Hom}_A^{}(DA,N_Y^{})[1]$$
is a tilting complex such that $`𝒳=𝒳_c(P^{})`$ and $`𝒴=𝒴_c(P^{})`$.
###### Proof.
According to Proposition 5.5, we have only to show that $`𝒳=𝒳_c(P^{})`$ and $`𝒴=𝒴_c(P^{})`$. It follows by \[Ho2\], Lemmas 2 and 3 that $`\mathrm{H}^0(P^{})𝒳`$ and $`\mathrm{H}^1(\nu (P^{}))𝒴`$. Since $`X`$ is a direct summand of $`\mathrm{H}^0(P^{})`$ and $`Y`$ is a direct summand of $`\mathrm{H}^1(\nu (P^{}))`$, it follows that $`\mathrm{H}^0(P^{})`$ generates $`𝒳`$ and $`\mathrm{H}^1(\nu (P^{}))`$ cogenerates $`𝒴`$. It now follows by Remark 2.9, Lemmas 3.2, 3.3 (2) that $`𝒳=𝒳_c(P^{})`$ and $`𝒴=𝒴_c(P^{})`$. ∎
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# 1 Abstract
## 1 Abstract
Interference effects in quantum transitions, giving rise to amplification without inversion, optical transparency and to enhancements in nonlinear optical frequency conversions are considered. Review of the relevant early theoretical and experimental results is given. The role of relaxation processes, spontaneous cascade of polarizations, local field effects, Doppler-broadening, as well as specific features of the interference in the spectral continuum are discussed.
Keywords: atomic coherence and interference, resonant nonlinear interactions, bound-free transitions, amplification without inversion, relaxation-induced processes, local field effects, inhomogeneous broadening, frequency-conversion, $`VUV`$ generation
## 2 Introduction
There has been considerable interest recently in the study of laser-induced quantum coherence and interference, which leads to fundamental effects in high resolution nonlinear spectroscopy, to amplification of radiation without the requirement of population inversion ($`AWI`$) and to resonantly enhanced refraction at vanishing (without) absorption ($`ERWA`$), to coherent population trapping and constructive contributions in resonantly enhanced nonlinear-optical frequency conversions and, at the same time, to distractive contributions in absorption of the fundamental and generated radiations <sup>1,2</sup>. Wide range of applications are expected <sup>3</sup>.
Resonant nonlinear optical interference effects have been subject of the extensive both theoretical and experimental studies since the discovery of masers and lasers (see for example <sup>2</sup> and ref. therein). In this paper we briefly review some early and recent results of Russian research groups on this topic.
## 3 Resonant nonlinear optical interference processes
### 3.1 Destructive and constructive interference in classical and quantum optical physics
Interference is one of the fundamental physical phenomena. Two oscillations at one and the same, or close, frequencies may interfere both in constructive and destructive ways. One can manipulate by the resulting oscillations with variation of the relative phase and the amplitudes of the interfering oscillators in order to enhance or, on the contrary, to eliminate the oscillations of any nature. Interference is widely used in optical physics, including quantum optics. The concept of interference is more general, then the notions of elementary quantum-optical processes, such as one-photon, multistep and multiphoton transitions. These notions were introduced and classify at their frequency-correlation properties in the framework of the perturbation theory. Indeed, in resonant interactions, these properties may be drastically changed with growth of the intensity of the coupled fields <sup>4;2b,c</sup>. The latter may give rise to such effects in nonlinear spectroscopy of Doppler broadened transition, as compensation of the residual inhomogeneous broadening in Raman-like and cascade configurations <sup>4;2b,c;5</sup>.
Quantum interference may occur when coherent superposition of real states is involved in a process<sup>6</sup>. Alternatively, interfering frequency-degenerate intraatomic oscillations may originate from different correlated quantum pathways, contributing in one and the same frequency. For example, in the weak-field approximation, these can be one- and two-photon contributions to an optical process, associated with the radiation at a given frequency. Such process may be thought as that started from the coherent superposition of closely spaced real energy-level and quasi-level (virtual state), created by the auxiliary strong field <sup>2b,c;4</sup>. Such a coherent superposition can be produced even more easily than in the case of real doublet state. In general, even in the cases, when many elementary processes contribute to an optical process and their classification is troublesome, one can explain and predict experimental results with the aid of the notion of interfering frequency-degenerated components of nonlinear polarization. The amplitudes of the components can be varied with the intensities and phases – with the frequency-detunings of the driving fields.
### 3.2 Effect of energy levels population and relaxation, density matrix approach
In general case of open energy-level configuration with all the levels being populated and various relaxation processes involved, density-matrix method is the most convenient for the analysis of a resonant nonlinear-optical response. Explicit formulae, describing spectral properties of a weak probe field in the presence of an auxiliary strong one, in cascade, $`V`$ and $`\mathrm{\Lambda }`$ configurations can be easily derived in the similar way <sup>2b,c</sup>. We shall show that on the example of the energy-level schematic, given on Fig.1.
Fields $`E_1`$ at frequency $`\omega _1\omega _{gl}`$ and $`E_3`$ at frequency $`\omega _3\omega _{mn}`$ are strong. Fields $`E_2`$ at frequency $`\omega _2`$ and $`E_4`$ at frequency $`\omega _4`$ are probe ones. We shall derive the conditions to achieve $`AWI`$ at the transition $`gn`$, as well as at transition $`ml`$, so that both $`V`$ and $`\mathrm{\Lambda }`$ configurations are embedded. Frequency of the probe field may be both higher and lower compared to the driving field.
Consider energy-level configuration, shown in the Fig.1. Density matrix equations in the interaction representation, relevant to the problem under consideration, can be written in the form:
$`\rho _{lg}=r_1exp(i\mathrm{\Omega }_1t)`$, $`\rho _{nm}=r_3exp(i\mathrm{\Omega }_3t)`$, $`\rho _{ng}=r_2exp(i\mathrm{\Omega }_1t)+\stackrel{~}{r}_2exp[i(\mathrm{\Omega }_1+\mathrm{\Omega }_3\mathrm{\Omega }_4)t]`$,
$`\rho _{lm}=r_4exp(i\mathrm{\Omega }_4t)+\stackrel{~}{r}_4exp[i(\mathrm{\Omega }_1\mathrm{\Omega }_2+\mathrm{\Omega }_3)t]`$, $`\rho _{ln}=r_{12}exp[i(\mathrm{\Omega }_1\mathrm{\Omega }_2)t]+r_{43}exp[i(\mathrm{\Omega }_4\mathrm{\Omega }_3)t]`$,
$`\rho _{ii}=r_i`$,
$`P_2r_2=iG_2\mathrm{\Delta }r_2iG_3r_{32}^{}+ir_{12}^{}G_1`$, $`d_2\stackrel{~}{r}_2=iG_3r_{41}^{}+ir_{43}^{}G_1`$,
$`P_4r_4=i\left[G_4\mathrm{\Delta }r_4G_1r_{41}+r_{43}G_3\right]`$, $`d_4\stackrel{~}{r}_4=iG_1r_{32}+ir_{12}G_3`$
$`P_{41}r_{41}=iG_1^{}r_4+ir_1^{}G_4`$, $`P_{43}r_{43}=iG_4r_3^{}+ir_4G_3^{}`$,
$`P_{32}r_{32}=iG_2^{}r_3+ir_2^{}G_3`$, $`P_{12}r_{12}=iG_1r_2^{}+ir_1G_2^{}`$,
$`\mathrm{\Gamma }_mr_m=2Re\{iG_3^{}r_3\}+q_m`$, $`\mathrm{\Gamma }_nr_n=2Re\{iG_3^{}r_3\}+\gamma _{gn}r_g+\gamma _{mn}r_m+q_n`$,
$`\mathrm{\Gamma }_gr_g=2Re\{iG_1^{}r_1\}+q_g`$, $`\mathrm{\Gamma }_lr_l=2Re\{iG_1^{}r_1\}+\gamma _{gl}r_g+\gamma _{ml}r_m+q_l`$,
$`\mathrm{\Delta }r_1=r_lr_g,\mathrm{\Delta }r_2=r_nr_g,\mathrm{\Delta }r_3=r_nr_m,\mathrm{\Delta }r_4=r_lr_m`$.
Where $`\mathrm{\Omega }_1=\omega _1\omega _{lg}`$, $`\mathrm{\Omega }_3=\omega _3\omega _{mn}`$, $`\mathrm{\Omega }_2=\omega _2\omega _{gn}`$, $`\mathrm{\Omega }_4=\omega _4\omega _{ml}`$,
$`G_1=𝐄_\mathrm{𝟏}𝐝_{lg}/2\mathrm{}`$, $`G_2=𝐄_\mathrm{𝟐}𝐝_{gn}/2\mathrm{}`$, $`G_3=𝐄_\mathrm{𝟑}𝐝_{nm}/2\mathrm{}`$, $`G_4=𝐄_\mathrm{𝟒}𝐝_{ml}/2\mathrm{}`$,
$`P_1=\mathrm{\Gamma }_{lg}+i\mathrm{\Omega }_1`$, $`P_2=\mathrm{\Gamma }_{ng}+i\mathrm{\Omega }_2`$, $`P_3=\mathrm{\Gamma }_{nm}+i\mathrm{\Omega }_3`$, $`P_4=\mathrm{\Gamma }_{lm}+i\mathrm{\Omega }_4`$,, $`P_{12}=\mathrm{\Gamma }_{ln}+i(\mathrm{\Omega }_1\mathrm{\Omega }_2)`$, $`P_{43}=\mathrm{\Gamma }_{ln}+i(\mathrm{\Omega }_4\mathrm{\Omega }_3)`$, $`P_{32}=\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_3\mathrm{\Omega }_2)`$, $`P_{41}=\mathrm{\Gamma }_{gm}+i(\mathrm{\Omega }_4\mathrm{\Omega }_1)`$,
$`d_2=\mathrm{\Gamma }_{ng}+i(\mathrm{\Omega }_1+\mathrm{\Omega }_3\mathrm{\Omega }_4)`$, $`d_4=\mathrm{\Gamma }_{lm}+i(\mathrm{\Omega }_1\mathrm{\Omega }_2+\mathrm{\Omega }_3)`$.
Here $`\mathrm{\Omega }_i`$ are frequency detuning from the resonances, $`G_i`$ — Rabi frequencies, $`\mathrm{\Delta }r_i`$ — power–depending population differences, $`\mathrm{\Gamma }_{ij}`$ — homogeneous half linewidths, $`\mathrm{\Gamma }_i^1`$ — lifetimes, $`\gamma _{ij}`$ — relaxation rates from $`i`$ to $`j`$ states, $`q_i`$ — population rate by a incoherent source. Density matrix amplitudes $`r_i`$ determine absorption/gain and refraction indexes, $`\stackrel{~}{r}_i`$ — determine four – wave mixing driving nonlinear polarizations.
The equations and their solution for the cascade atomic configurations can be derived by the simple change of the detunings signs <sup>2b</sup>.
### 3.3 Laser–induced atomic coherence and classification of resonant nonlinear effects
Solution of the coupled density – matrix equations may be represented in the form:
$$r_{1,3}=iG_{1,3}\mathrm{\Delta }r_1/P_1,r_{2,4}=iG_{2,4}R_{2,4}/P_{2,4},$$
$$R_2=\frac{\mathrm{\Delta }r_2(1+g_7+v_7)v_3(1+v_7g_8)\mathrm{\Delta }r_3g_3(1+g_7v_8)\mathrm{\Delta }r_1}{(1+g_2+v_2)+[g_7+g_2(g_7v_8)+v_7+v_2(v_7g_8)]},$$
$`(1)`$
$$R_4=\frac{\mathrm{\Delta }r_4(1+v_5+g_5)g_1(1+g_5v_6)\mathrm{\Delta }r_1v_1(1+v_5g_6)\mathrm{\Delta }r_3}{(1+g_4+v_4)+[v_5+v_4(v_5g_6)+g_5+g_4(g_5v_6)]},$$
$`(2)`$
$$\mathrm{\Delta }r_1=\frac{(1+\text{æ}_3)\mathrm{\Delta }n_1+b_1\text{æ}_3\mathrm{\Delta }n_3}{(1+\text{æ}_1)(1+\text{æ}_3)a_1\text{æ}_1b_1\text{æ}_3},\mathrm{\Delta }r_3=\frac{(1+\text{æ}_1)\mathrm{\Delta }n_3+a_1\text{æ}_1\mathrm{\Delta }n_1}{(1+\text{æ}_1)(1+\text{æ}_3)a_1\text{æ}_1b_1\text{æ}_3},$$
$$\mathrm{\Delta }r_2=\mathrm{\Delta }n_2b_2\text{æ}_3\mathrm{\Delta }r_3a_2\text{æ}_1\mathrm{\Delta }r_1,\mathrm{\Delta }r_4=\mathrm{\Delta }n_4a_3\text{æ}_1\mathrm{\Delta }r_1b_3\text{æ}_3\mathrm{\Delta }r_3;$$
$$r_m=n_m+(1b_2)\text{æ}_3\mathrm{\Delta }r_3,r_g=n_g+(1a_3)\text{æ}_1\mathrm{\Delta }r_1,r_n=n_nb_2\text{æ}_3\mathrm{\Delta }r_3+a_1\text{æ}_1\mathrm{\Delta }r_1,$$
$`(3)`$
$$r_l=n_lb_1\text{æ}_3\mathrm{\Delta }r_3+a_3\text{æ}_1\mathrm{\Delta }r_1,\mathrm{\Delta }r_i(E_1=0,E_3=0)=\mathrm{\Delta }n_i;$$
$$g_1=\frac{|G_1|^2}{P_{41}P_1^{}},g_2=\frac{|G_1|^2}{P_{12}^{}P_2},g_3=\frac{|G_1|^2}{P_{12}^{}P_1^{}},g_4=\frac{|G_1|^2}{P_{41}P_4},g_5=\frac{|G_1|^2}{P_{43}d_2^{}},g_6=\frac{|G_1|^2}{P_{41}d_2^{}},g_7=\frac{|G_1|^2}{P_{32}^{}d_4^{}},g_8=\frac{|G_1|^2}{P_{12}^{}d_4^{}},$$
$$v_1=\frac{|G_3|^2}{P_{43}P_3^{}},v_2=\frac{|G_3|^2}{P_{32}^{}P_2},v_3=\frac{|G_3|^2}{P_{32}^{}P_3^{}},v_4=\frac{|G_3|^2}{P_{43}P_4},v_5=\frac{|G_3|^2}{P_{41}d_2^{}},v_6=\frac{|G_3|^2}{P_{43}d_2^{}},v_7=\frac{|G_3|^2}{P_{12}^{}d_4^{}},v_8=\frac{|G_3|^2}{P_{32}^{}d_4^{}};$$
$$\text{æ}_1=\text{æ}_1^0\frac{\mathrm{\Gamma }_{lg}^2}{|P_1|^2},\text{æ}_1^0=\frac{2(\mathrm{\Gamma }_l+\mathrm{\Gamma }_g\gamma _{gl})}{\mathrm{\Gamma }_l\mathrm{\Gamma }_g\mathrm{\Gamma }_{lg}}|G_1|^2,\text{æ}_3=\text{æ}_3^0\frac{\mathrm{\Gamma }_{mn}^2}{|P_3|^2},\text{æ}_3^0=\frac{2(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn})}{\mathrm{\Gamma }_m\mathrm{\Gamma }_n\mathrm{\Gamma }_{mn}}|G_3|^2;$$
$$a_1=\frac{\gamma _{gn}a_2}{\mathrm{\Gamma }_n\gamma _{gn}}=\frac{\gamma _{gn}\mathrm{\Gamma }_la_3}{\mathrm{\Gamma }_n(\mathrm{\Gamma }_g\gamma _{gl})}=\frac{\gamma _{gn}\mathrm{\Gamma }_l}{\mathrm{\Gamma }_n(\mathrm{\Gamma }_l+\mathrm{\Gamma }_g\gamma _{gl})},$$
$$b_1=\frac{\gamma _{ml}\mathrm{\Gamma }_nb_2}{\mathrm{\Gamma }_l(\mathrm{\Gamma }_m\gamma _{mn})}=\frac{\gamma _{ml}b_3}{\mathrm{\Gamma }_l(\mathrm{\Gamma }_l\gamma _{ml})}=\frac{\gamma _{ml}\mathrm{\Gamma }_n}{\mathrm{\Gamma }_l(\mathrm{\Gamma }_m+\mathrm{\Gamma }_n\gamma _{mn})}.$$
By substituting frequency deviations $`\mathrm{\Omega }_i`$ for that Doppler-shifted $`\mathrm{\Omega }_ik_iv`$ ($`v`$ is atomic velocity) we can take into account the effect of atomic motion. Imaginary part of density-matrix amplitudes $`r_2`$ and $`r_4`$ represent absorption or gain at the corresponding probe-field frequencies. At $`G_3=0`$ equations (2) and (3) convert in solutions for $`\mathrm{\Lambda }`$ and $`V`$ schemes
$$r_2=i\frac{G_2}{P_2}\frac{\mathrm{\Delta }r_2g_3\mathrm{\Delta }r_1}{1+g_2},r_4=i\frac{G_4}{P_4}\frac{\mathrm{\Delta }r_4g_1\mathrm{\Delta }r_1}{1+g_4}.$$
$`(4)`$
Following<sup>2b,c</sup> we can classify resonant nonlinear effects as 1) power saturation of the populations (eq. (4)); 2) strong-field induced splitting of the probe-field resonances, or ac Stark effect (denominators in eqs. (4)); and 3) nonlinear interference effect ($`NIEF`$) (second and third terms in the nominators of eqs. (2)).
## 4 Difference in absorption and emission spectra due to the nonlinear interference effects, amplification without inversion, resonance–enhanced refraction without absorption
Power of emitted or absorbed radiation, for example at the frequency $`\omega _2`$, which is proportional to $`Re(iG_2^{}r_2)`$, can be considered as a difference between pure emission (the term, proportional to $`r_g`$) and pure absorption (the rest terms in eqs.$`(2)`$ ). The difference in frequency-dependence of these terms, induced by the auxiliary driving field, is the origin of $`AWI^{2b}`$. Refractive index at $`\omega _2`$ is determined by $`Im(iG_2^{}r_2)`$ and, in general, laser-induced minimum in absorption may coincide with the resonance-enhanced maximum in refraction <sup>1,3</sup>. Thus, laser-induced resonance splitting and $`NIEF`$ transform only spectral shape of absorption/gain and refractive indices, give rise to difference in the line shapes of spontaneous (or pure induced) emission and absorption, but do not affect the integral intensity of the spectral lines <sup>2b,c</sup>:
$$𝑑\mathrm{\Omega }_2Re(ir_2/G_2)=\mathrm{\Delta }r_2,𝑑\mathrm{\Omega }_4Re(ir_4/G_4)=\mathrm{\Delta }r_4.$$
$`(5)`$
Indeed, $`NIEF`$ give rise to electromagnetically induced transparency ($`EIT`$) and to $`AWI`$ at the transitions $`gn`$ (or $`ml`$), when contributions of second and third terms in the nominators of eqs.$`(2)`$ are equal or dominate over $`\mathrm{\Delta }r_2`$ (or over $`\mathrm{\Delta }r_4`$), correspondingly. From the above presented density-matrix equations one can see that the coherence at the transitions $`gm`$ and $`ln`$ ($`r_{32}`$ and $`r_{12}`$), induced in cooperation of the strong and the probe fields, is the source of the $`EIT`$ and $`AWI`$ effects.
A great number of elementary processes, introduced and defined for the bare states in the framework of the perturbation theory, may give contribution to the absorption/gain index $`\alpha (\mathrm{\Omega }_i)`$. Consider, for example, $`\alpha (\mathrm{\Omega }_4)`$ at the frequency $`\omega _4>\omega _1`$ (Fig.1), reduced by it’s maximum value $`\alpha ^0(0)`$ in the absence of the all strong fields, for the case when $`E_3=0`$. From the eqs.$`(2)`$ one finds:
$$\frac{\alpha (\mathrm{\Omega }_4)}{\alpha ^0(0)}=Re\{\frac{\mathrm{\Gamma }_4}{P_4}\frac{\mathrm{\Delta }r_4g_1\mathrm{\Delta }r_1}{\mathrm{\Delta }n_4(1+g_4)}\}$$
$`(6)`$
Consider two subcases:
a.Off resonance:$`|\mathrm{\Omega }_1||\mathrm{\Omega }_4|>>\mathrm{\Gamma }_1,\mathrm{\Gamma }_4;|g_4|<<1;|g_1|<<1;P_4i\mathrm{\Omega }_4;P_1i\mathrm{\Omega }_1i\mathrm{\Omega }_4`$.
Eq.$`(6)`$ takes the form:
$$\frac{\alpha (\mathrm{\Omega }_4)}{\alpha ^0(0)}\frac{\mathrm{\Gamma }_4^2\mathrm{\Delta }r_4}{\mathrm{\Omega }_4^2\mathrm{\Delta }n_4}Re\{\frac{\mathrm{\Gamma }_4(\mathrm{\Delta }r_4g_4+\mathrm{\Delta }r_1g_1)}{i\mathrm{\Omega }_4\mathrm{\Delta }n_4}\}\frac{\mathrm{\Gamma }_4^2\mathrm{\Delta }r_4}{\mathrm{\Omega }_4^2\mathrm{\Delta }n_4}\frac{\mathrm{\Gamma }_4\mathrm{\Gamma }_{14}}{\mathrm{\Gamma }_{14}^2+(\mathrm{\Omega }_4\mathrm{\Omega }_1)^2}\frac{|G_1|^2(\mathrm{\Delta }r_1\mathrm{\Delta }r_4)}{\mathrm{\Omega }_4^2\mathrm{\Delta }n_4}=$$
$$=\frac{\mathrm{\Gamma }_{lm}^2(r_lr_m)}{(n_ln_m)\mathrm{\Omega }_4^2}\frac{\mathrm{\Gamma }_{gm}\mathrm{\Gamma }_{lm}}{\mathrm{\Gamma }_{gm}^2+(\mathrm{\Omega }_4\mathrm{\Omega }_1)^2}\frac{|G_1|^2(r_mr_g)}{\mathrm{\Omega }_4^2(n_ln_m)}$$
$`(7)`$
The last terms in eqs.$`(7)`$ describe Raman-like coupling and originate both from the nominator and the denominator in eq.$`(6)`$. Population inversion between initial and final bare states ($`r_m=n_m>r_g`$) is required for amplification of the probe field.
b.Resonance: $`\mathrm{\Omega }_1=\mathrm{\Omega }_4=0`$.
Conditions for $`AWI`$ and $`EIT`$ are:
$$g_1\mathrm{\Delta }r_1\mathrm{\Delta }r_4,or\frac{|G_1|^2}{\mathrm{\Gamma }_{lg}\mathrm{\Gamma }_{gm}}(r_lr_g)r_lr_m$$
$`(8)`$
Eq.$`(8)`$ shows that due to $`NIEF`$, population inversion between initial and final bare states is not required in order to attain $`AWI`$ in this case. Small relaxation rate of the coherence, induced in cooperation of the driving and probe fields, compared to the other relaxation rates is the most important.
Analysis of the condition for $`EIT`$ and $`AWI`$ as well as of sign–changing line shape in the more details can be found in ref.<sup>2b</sup> both for open and closed ($`l`$ is ground state) atomic configurations. The analysis shows strong dependence of the line shape on the ratios of both population and coherence relaxation rates as well as on the ratios of initial unsaturated population differences on the coupled transitions.
### 4.1 Constructive and destructive interference due to the atomic velocity distribution
Furthermore, the analysis shows that the contributions of the coherence driving fields to the spectra may be both constructive and destructive, depending on the detunings of the probe as well as of the strong fields. This indicates that in gases with inhomogeneous broadening of the coupled transitions, dominating over homogeneous one, conditions for $`AWI`$ and $`EIT`$ may considerably differ from that for atoms at rest. Nevertheless, it was found out that under certain conditions sign-changing spectral profiles may be produced too <sup>2b,c;7a,b</sup>. At weak intensities of driving field narrow structure, superimposed on the Doppler background, appears. The shape of the structure is anisotropic and depends on the angle between the wave vectors of the interacting radiations. Optically-pumped unidirectional-emitting ring laser may operate by that. The line shape is also dependent on the intensity of the driving field and velocity-changing collisions. Special features may occur, when some of the coupled transitions are homogeneously, and some of them are inhomogeneously broadened. It was found out that destructive or constructive character of the effect of Maxwell’s velocity distribution depends on the fact whether a frequency of the probe field is less or greater than that of the strong one too. Analytical results describing general behavior of the velocity-averaged functions for some limiting cases,including Rabi frequencies larger then homogeneous linewidths, can be found in ref.<sup>2a,b,c;7a,b</sup>.
## 5 Coherence and nonlinear-optical conversion, Enhancements in nonlinear-optical conversion due to multiple resonance and induced transparency, Local-field effects
Nonlinear optical response of a medium experiences a giant enhancements in one- and multiphoton resonances. This reduces required fundamental powers down to $`cw`$ regime <sup>8</sup>, however imposes severe limitations on the number density of the medium due to absorption of fundamental and generated radiations. As it is discussed above, in the presence of a strong electromagnetic radiation resonances for a weak probe radiation experience splitting <sup>2,9</sup>, which exhibits itself in a different ways in real and imaginary parts of linear and nonlinear susceptibilities. Later makes possible to combine decrease in absorption with increase in squared module of nonlinear susceptibilities, responsible for optical generation, and at the same time with improvements in phase-matching and increasing density of the medium <sup>2f,10</sup>.
With the increase of the atom number density, local field, acting on an atom, may pretty much differ from the external field both in the amplitude and phase. This may drastically change shape of nonlinear spectroscopic structures, including electromagnetically induced transparency<sup>11,12</sup>.
Consider experimental schematics,proposed in ref.<sup>13</sup>, that combines the advantages of both multiple resonance enhancements and increase in atom number density of nonlinear medium due to the above mentioned Autler-Townes ($`ac`$ Stark splitting) as well as local-field effects.
Consider energy-level scheme, shown in Fig.$`2a.`$ Strong fields at frequencies $`\omega _3`$ and $`\omega _2`$ couple unpopulated levels $`3`$ and $`2`$ (Rabi frequency $`G_3)`$ and 2 and 1 (Rabi frequency $`G_2)`$, respectively. Field at $`\omega _1\omega _{10}`$ as well as generated $`\omega _s=\omega _1+\omega _2+\omega _3`$ are weak, do not change populations of the levels and are accounted for only in the lowest order of the perturbation theory. Absorption and refraction indexes for the probe fields at $`\omega _1`$ and $`\omega _s`$ are represented by the imaginary and real parts of
$$\chi _1(\omega _1;\omega _1)=(\chi _1^0/P_{01})f_1,\chi _s(\omega _s;\omega _s)=(\chi _s^0/P_{03})f_s$$
$`(9)`$
Nonlinear susceptibility is:
$$\chi ^{NL}(\omega _s;\omega _1+\omega _2+\omega _3)=(\chi _0^{NL}/P_{01}P_{02}D_{03})f,$$
$`(10)`$
where $`\chi _1^0,\chi _s^0`$ and $`\chi _0^{NL}`$ are resonant values of the susceptibilities at negligibly small $`G_2`$ and $`G_3`$. Factors $`f_1`$, $`f_2`$ and $`f`$ describe effects of the strong fields. Simple density - matrix calculations, similar to given in<sup>2b,e;10a.</sup> yield:
$$f_1=\{1+g_2/P_{01}P_{02}[1+(g_3/P_{02}D_{03})]\}^1,$$
$`(11)`$
$$f_s=\{1+g_3/P_{03}D_{02}[1+(g_2/D_{02}D_{01})]\}^1,$$
$`(12)`$
$$f=f_1[1+g_3/D_{03}P_{02}]^1=[1+(g_2/D_{02}D_{01})+(g_3/D_{03}P_{02})]^1$$
$`(13)`$
$$P_{01}=1+ix_1,P_{02}=1+ix_0,P_{03}=1+ix_s;D_{01}=1+iy_1,D_{02}=1+iy_0,D_{03}=1+iy_s;$$
$$x_1=(\omega _1\omega _{10})/\mathrm{\Gamma }_{10}=0,x_{02}=(\omega _1+\omega _2\omega _{21})/\mathrm{\Gamma }_{20}=0,x_s=(\omega _s\omega _{30})/\mathrm{\Gamma }_{30}=0;$$
$$y_1=(\omega _s\omega _3\omega _2\omega _{10})/\mathrm{\Gamma }_{10}=0,y_{02}=(\omega _s\omega _3\omega _{21})/\mathrm{\Gamma }_{20}=0,y_s=(\omega _1+\omega _2+\omega _3\omega _{30})/\mathrm{\Gamma }_{30}=0;$$
$$g_2=G_2^2/\mathrm{\Gamma }_{10}\mathrm{\Gamma }_2,g_3=G_3^2/\mathrm{\Gamma }_{30}\mathrm{\Gamma }_{20},$$
$`\mathrm{\Gamma }_{ij}`$ are homogeneous halfwidth of the corresponding transitions. In the case, when $`E_s`$ is not a probe field, but generated radiation, $`\omega _s=\omega _1+\omega _2+\omega _3`$ and $`D_{0i}=P_{0i}`$.
Factors $`f_1,f_s`$ and $`f`$ are different and describe splitting of the corresponding resonances. Frequency-dependence and difference from unity of the factors $`f_1,f_s`$ and $`f`$ is determined by the coherence, induced at the transition $`02`$ by the two combinations of strong and weak fields $`(E_1,E_2`$ and $`E_s,E_3`$). Generated power $`Pg_2g_3\chi ^{NL}^2`$, depends not only on imaginary but on real part of $`\chi ^{NL}`$ too, and because of that may not deplete in the spectral range of induced transparency and phase-matching. Each resonance increases $`\chi ^{NL}^2`$ by the factor of $`x_i^2`$ , which may be on the order of $`10^6`$. Laser-induce spectral structures in real parts of $`\chi _1`$ and $`\chi _s(`$dispersion caused by the coherence at the $`02`$ transition), provide additional means to phase-match frequency - conversion by the small detunings of the fundamental radiations from the resonances. Triple resonance may yield total enhancement in generated power on the order of $`10^{18}`$. Due to the induced transparency, number density of the atoms $`N`$ and consequently $`PN^2`$ may be increased by several orders of the magnitude in addition.
At high number density of the atoms, local fields may significantly differ from the external electromagnetic fields both in amplitudes and in phases. As it was shown in <sup>12,13</sup>, that may drastically change spectral properties of the induced transparency as well as of the generating nonlinear polarization. Similar to <sup>11,12</sup>, making use Lorentz-Lorenz approximation, local field effects can be accounted for by the substituting one- photon resonances on that red-shifted (by substituting $`x_1`$ and $`x_s`$ for $`x_1+C_1`$ and $`x_s+C_s,C_1=Nd_{10}^2/3ϵ_0\mathrm{\Gamma }_{10};C_s=Nd_{30}^2/3ϵ_0\mathrm{\Gamma }_{30}`$, $`ϵ_0`$\- is permittivity of free space) . Due to the fact that this does not influence transition frequencies between the excited states and that of the multiphoton transitions, the introduced shifts may drastically change effects of strong electromagnetic fields at $`\omega _2`$ and $`\omega _3`$ on both dressed linear and nonlinear responses.
Equations, given above, can be easily generalized on the cases of the higher order processes. For example, when $`1`$$`0`$ and/or $`3`$$`2`$, $`2`$$`1`$ are multiphoton transitions, generalization can be done simply by substituting one-photon Rabi frequencies and detunings for the corresponding multiphoton ones. Manipulations by the nonlinear susceptibility, absorption and refractive indexes for the generating radiation with the auxilliary strong fields, coupled to the adjacent transitions (both bound and continuum states, Figs. $`2b,c.`$), were considered in ref.<sup>2f,10</sup>.
## 6 Nonlinear interference effects at bound-free transitions, Laser-induced autoionizing-like resonances (laser induced continuum structure)
Nonlinear interference phenomena, similar to those at bound-bound transitions, including $`AWI`$ and $`EIT`$, can occur at the transitions to ionization continuum. Appropriate theory was developed in ref. <sup>2f,10a,14</sup>. Similar case, relevant to the zone bands in crystals, was considered in ref.<sup>15</sup>. Laser induced autoionizing like resonances – laser induced continuum structure ($`LICS`$) was observed in the experiments ref.<sup>16</sup>, and since the end of $`80^{}s`$ studies of the resonant interference processes in the context of $`LICS`$, $`AWI`$ and $`EIT`$, first at bound–free and then at bound–bound transitions, have involved a number of research groups<sup>17,18</sup>.
Potential feasibilities to manipulate both by $`LICS`$ and by the splitting of the discrete resonances in order to enhance short - wavelengths frequency - mixing output and to decrease resonant absorption of the both fundamental and generated radiations can be shown with the example of Fig.$`2b.`$, generalized for the case, when $`\omega _1`$ is close to $`\omega _{10}`$, and radiations at $`\omega _2`$, $`\omega _3`$ and $`\omega `$ are strong. The example combines opportunities to manipulate by two $`LICS`$ and by depletion of absorption at the discrete transitions. Contribution of strong off - resonant $`k`$ levels are taken into account too. By that, the detunings $`\omega _1`$-$`\omega _{gm}`$, $`\omega _1+\omega _2\omega _{gn}`$ and $`\omega \omega _3\omega _{nl}`$ are assumed being much less than all the rest. Density - matrix calculations give the expressions for nonlinear susceptibility $`\chi ^{(3)}(\omega _\mu =\omega _1+\omega _2+\omega _3)`$, which determines generated power at the frequency $`\omega _\mu `$, as well as for absorption indexes $`\alpha (\omega _1)`$ and $`\alpha (\omega _\mu )`$ for probe radiations at corresponding frequences as follows <sup>19</sup>:
$$\chi ^{(3)}(\omega _\mu =\omega _1+\omega _2+\omega _3)/\chi _{0\mu }^{(3)}=K/(D_{gm}X),$$
$`(14)`$
$$\alpha (\omega _1)/\alpha _{01}=Re\{[1g_{mn}/(D_{gm}X)]/D_{gm}\},$$
$`(15)`$
$$\alpha (\omega _\mu )/\alpha _{0\mu }=1k_3\beta _l+k_3\beta _l(y_l+q_{gl})^2/(1+y_{l}^{}{}_{}{}^{2})$$
$$Re\{k_4g_{nn}A^2(1iq_{gn})^2/Y\}$$
$`(16)`$
where $`\chi _{0\mu }^{(3)}`$, $`\alpha _{01}`$ and $`\alpha _{0\mu }`$ \- are corresponding resonant values at the intencities of all the fields beeing negligibly weak. The rest parameters are as follows:
$$K=1k_1\beta _l[(1iq_{nl})(1iq_{lg})]/[(1iq_{ng})(1+ix_l)],$$
$`(17)`$
$$A=1k_1\beta _l[(1iq_{ln})(1iq_{gl})]/[(1iq_{gn})(1+iy_l)],$$
$`(18)`$
$$X=(1+g_{nn})[1+ix_n+g_{mn}/D_{gm}(1+g_{nn})k_2\beta _l\beta _n(1iq_{nl})^2/(1+ix_l)],$$
$`(19)`$
$$Y=(1+g_{nn})[1+iy_n+g_{mn}/p_{gm}(1+g_{nn}k_2\beta _l\beta _n(1iq_{nl})^2/(1+iy_l)],$$
$`(20)`$
$$D_{gm}=1+i(\omega _1\omega _{gm})/\mathrm{\Gamma }_{gm},p_{gm}=1+i(\omega _\mu \omega _3\omega _2\omega _{gm})/\mathrm{\Gamma }_{gm},$$
$`(21)`$
$$x_l=(\omega _1+\omega _2+\omega _3\omega \omega _{gl}\delta _{ll})/(\mathrm{\Gamma }_{gl}+\gamma _{ll}),x_n=(\omega _1+\omega _2\omega _{gn}\delta _{nn})/(\mathrm{\Gamma }_{gn}+\gamma _{nn}),$$
$`(22)`$
$$y_l=(\omega _\mu \omega \omega _{gl}\delta _{ll})/(\mathrm{\Gamma }_{gl}+\gamma _{ll}),y_n=(\omega _\mu \omega _3\omega _{gn}\delta _{nn})/(\mathrm{\Gamma }_{gn}+\gamma _{nn}),$$
$`(23)`$
$$k_1=(\gamma _{gl}\gamma _{ln})/(\gamma _{gn}\gamma _{nn}),k_2=(\gamma _{nl}\gamma _{ln})/(\gamma _{ll}\gamma _{nn}),k_3=(\gamma _{gl}\gamma _{lg})/(\gamma _{gg}\gamma _{ll}),k_4=(\gamma _{gn}\gamma _{ng})/(\gamma _{gg}\gamma _{nn}),$$
$`(24)`$
$$g_{mn}=G_{mn}^2/\mathrm{\Gamma }_{gm}\mathrm{\Gamma }_{gn},\beta _l=g_{ll}/(1+g_{ll}),\beta _n=g_{nn}/(1+g_{nn}),$$
$`(25)`$
$$g_{ii}=\gamma _{ii}/\mathrm{\Gamma }_{gi},q_{ij}=\delta _{ij}/\gamma _{ij},\gamma _{ij}=\pi \mathrm{}G_{iϵ}G_{ϵj}|_{ϵ=\mathrm{}\omega _\mu }+Re\{\underset{k}{}G_{ik}G_{kj}/p_{gk}\},$$
$`(26)`$
$$\delta _{ij}=\mathrm{}P𝑑ϵG_{iϵ}G_{ϵj}/(\mathrm{}\omega _\mu ϵ)+Im\{\underset{k}{}G_{ik}G_{kj}/p_{gk}\}$$
$`(27)`$
Factors $`0k_i1`$, depending on whether continuum states are not degenerate or degenerate (unity).
Comparing eqs.$`(14)`$ and $`(16)`$ with corresponding equations from ref.<sup>10a,2f</sup>, one can see additional interference $`LICS`$ structures in generating nonlinear polarization, absorption and refraction indexes, produced in cooperation by the $`E_3`$ and $`E`$ fields (terms, proportional to $`\beta _n`$ and $`g_n`$), which provide with the supplementary means in absorption spectroscopy and for enhancements of generated short-wavelength radiation.
## 7 Relaxation-induced coherence processes
As it was outlined above, relaxation may influence interference processes both in negative and positive ways. Consider examples, when role of relaxation is positive.
### 7.1 $`AWI`$ due to interference in spontaneous cascade of polarizations
The features in absorption and emission spectra, discussed above, are caused by interference of contributions of probe field and combination of probe and auxiliary strong field in atomic polarization. As it was outlined above, there may be other sources of interfering intraatomic oscillations. One of the means to obtain $`AWI`$ without making use of auxiliary strong fields has been suggested recently in ref.<sup>20</sup>. The origin is interference through the correlations in spontaneous decay.
Consider four-level atomic configuration shown in Fig.$`3a.`$ All four transitions are allowed. Suppose, that the transition frequency $`\omega _{mn}`$ is close to $`\omega _{m_1n_1}`$, and $`\omega _{m_1m}`$ is close to $`\omega _{n_1n}`$, that is difference $`\mathrm{\Delta }`$
$$\mathrm{\Delta }=\omega _{m_1n_1}\omega _{mn}=\omega _{m_1m}\omega _{n_1n}$$
$`(28)`$
is small. In this case interference between considered four radiating channels is possible. It is caused by the coherence transfer due to interaction with the vacuum oscillations, besides the populations decay and spontaneous emissions of photons. For the absorption index in the frequency range around $`\omega _{mn}`$ calculations give:
$$\alpha (\mathrm{\Omega })=\frac{\lambda ^2}{4\pi }\{N_{nm}A_{mn}\frac{\mathrm{\Gamma }}{\mathrm{\Gamma }^2+\mathrm{\Omega }^2}+N_{n_1m_1}A_{m_1n_1}[\frac{\mathrm{\Gamma }_1}{\mathrm{\Gamma }_1^2+(\mathrm{\Omega }\mathrm{\Delta })^2}+\frac{KC}{\mathrm{\Gamma }\mathrm{\Gamma }_1}f(\mathrm{\Omega })]\},$$
$`(29)`$
$$C=\sqrt{A_{m_1m}A_{n_1n}A_{mn}/A_{m_1}n_1},K=(1)^{J_m+J_{n_1}}\sqrt{2J_m+1}\sqrt{2J_{n_1+1}}\{\begin{array}{ccc}J_m& J_n& 1\\ J_{n_1}& J_{m_1}& 1\end{array}\},$$
$`(30)`$
$$f(\mathrm{\Omega })=Re\frac{\mathrm{\Gamma }\mathrm{\Gamma }_1}{(\mathrm{\Gamma }i\mathrm{\Omega })[\mathrm{\Gamma }_1i(\mathrm{\Omega }\mathrm{\Delta })]}=\frac{\mathrm{\Gamma }\mathrm{\Gamma }_1[\mathrm{\Gamma }\mathrm{\Gamma }_1\mathrm{\Omega }(\mathrm{\Omega }\mathrm{\Delta })]}{(\mathrm{\Gamma }^2+\mathrm{\Omega }^2)[\mathrm{\Gamma }_1^2+(\mathrm{\Omega }\mathrm{\Delta })^2]},$$
$`(31)`$
$$N_{nm}=(2J_m+1)(\rho _n\rho _m),N_{n_1m_1}=(2J_{m_1}+1)(\rho _{n_1}\rho _{m_1}).$$
$`(32)`$
Here $`\mathrm{\Omega }=\omega \omega _{mn}`$, $`A_{ij}`$ – Einstein coefficients, $`\mathrm{\Gamma }`$, $`\mathrm{\Gamma }_1`$ – are line halfwidths for the interfering transitions, $`J_i`$ – energy level momenta, $`N_{ij}`$ – population differences.
The interference term is described by the function $`f(\mathrm{\Omega })`$, $`f(\mathrm{\Omega })𝑑\mathrm{\Omega }=0`$ . Coefficient $`K`$ is determined by the moments of four levels under consideration and may vary in the interval $`1K1`$. The case $`K0`$ corresponds to constructive interference (enhancements in the oscillations), the case $`K0`$ — to destructive interference. The analysis of the lineshape eq.$`29`$ shows it sign-changing behavior. For example, at $`\mathrm{\Omega }\mathrm{\Delta }`$
$$\alpha (\mathrm{\Omega })=\frac{\lambda ^2}{4\pi \mathrm{\Omega }^2}\{N_{nm}A_{mn}\mathrm{\Gamma }+N_{n_1m_1}A_{m_1n_1}(\mathrm{\Gamma }_1KC)\},$$
$`(33)`$
According to eq.33, absorption index may occur negative ($`AWI`$), if the requirements
$$K>O,(KC/\mathrm{\Gamma }_11)N_{n_1m_1}A_{m_1n_1}\mathrm{\Gamma }_1>N_{nm}A_{mn}\mathrm{\Gamma }$$
$`(34)`$
are met. When $`K0`$, $`\mathrm{\Delta }=0`$, the condition
$$(|K|C/\mathrm{\Gamma }1)N_{n_1m_1}A_{m_1n_1}\mathrm{\Gamma }>N_{nm}A_{mn}\mathrm{\Gamma }_1$$
$`(35)`$
means appearance of $`AWI`$ in the line center $`(\mathrm{\Omega }=0)`$. Similar phenomena may occur in the spectral range of the doublet $`\omega _{m_1m}`$, $`\omega _{n_1n}`$. Thus, in the considered atomic configuration $`AWI`$ may be provided by the correlations in the spontaneous decay without any external action.
### 7.2 Collision-induced four-wave mixing
Consider example, when collisions and spontaneous relaxation, as well as external magnetic field, break destructive interference<sup>8a</sup> . This remove elimination of for-wave mixing process and provides with the test, selectively sensitive to the specific modes of relaxation. The experiment was carried out with $`HeNe`$ laser, $`\lambda =1.52\mu m`$, which is resonant to $`2s_2`$$`2p_4`$ transition of $`Ne`$. The upper level consist of three Zeeman’s sublevel ($`J_1=1`$), the lower one is singlet ($`J_0=0`$). Fundamental beam consisted of two linear and orthogonal polarized components $`E_1`$ and $`E_2`$, frequency-shift $`\mathrm{\Delta }=\omega _2\omega _1`$ being much less than natural transition linewidth. Intensity of the radiation at $`\omega _1`$ was much greater then that at $`\omega _2`$. Collision and magnetic field sensitive four-wave mixing output $`E_s`$ at $`\omega _s=2\omega _1\omega _2=\omega _22\mathrm{\Delta }`$ and with the same polarization as $`E_2`$ was detected. Growth of the $`FWM`$ signal with the increase of collision rate and strength of magnetic field was observed, that can be explained as follows.
Each field and emitting nonlinear polarization $`P^{NL}(\omega _s)`$ may be represented as combination of two circular polarized components $`P_+^{NL}(\omega _s)`$ and $`P_{}^{NL}(\omega _s)`$. Formulae for each of these components of nonlinear polarization consist of two terms. One of them describes $`FWM`$ of the radiations with one and the same polarizations in two-level subsystem, another one – $`FWM`$ of the waves with opposite polarizations in three-level Zeeman’s subsystem (Fig.$`3b`$). In the schematic under consideration, it turned out, that the two contributions interfere in the distractive way and completely eliminate each other, provided by the relaxation rates of population and quadruple moment (alignment) in the upper level are equal. It is obvious that trapping of the spontaneous radiation from the upper level, anisotropic collisions, as well as external magnetic field break the counterbalance and, therefore, induce $`FWM`$ output. Such dependence was observed in the experiments. External magnetic field turns the second channel into fully resonant double-$`V`$ schematics.
## 8 Review of early theory and experiments on $`NIEF`$, $`AWI`$ and related phenomena
Coherence phenomena in three-level systems were studied since discovery of masers. Feasibility to attain $`AWI`$ in these systems was discussed in some of publications of that period both for microwave<sup>21</sup> and optical transitions<sup>22</sup>. $`AWI`$ in optical two-level systems was predicted in ref.<sup>23</sup> and first was observed in radio-frequency transitions<sup>24,2d</sup>. In optical range $`AWI`$ and corresponding features in refractive index were observed in ref.<sup>2e,25</sup>. Studies of coherence and interference phenomena in quantum transitions is growing research area, since they are embedded in many optical processes of basic and practical importance.
## 9 Concluding remarks
As it was outline, interference is basic and very general phenomenon of optical physics, which may play a crucial role in many experimental schematics of resonant nonlinear optics. Some of such schematics are shown in the Fig.$`4`$.
Fig.$`4a.`$ shows upconversion of weak infrared radiation at the frequency $`\omega _2`$. Fields $`E_1`$ and $`E_3`$ are strong. Destructive interference of oscillations at the frequency $`\omega _s\omega _3=\omega _{ng}=\omega _1+\omega _2`$ was shown to be one of the main process, limiting the conversion efficiency<sup>26</sup>. Fig.$`4b.`$ – interference of multiphoton transition and one-photon, induced by the generating radiation eliminates population of the upper level. Fig.$`4c.`$ – off-resonant 7th-order seventh-harmonic generation interfere with resonant 9th-order seventh-harmonic generation, that was used for detection of the processes<sup>26</sup>. Figs.$`4d,e.`$ – interference of contributions of the doublet sublevels in two-photon and off-resonant one-photon transitions. Figs.$`4f,g.`$ – interference of doublet sublevels in $`FWM`$.
Pressure-induced resonance was first proposed and experimentally proved in<sup>8a</sup> and later in<sup>27</sup>. The entire analogy between the schemes 3b and 4f is seen from the formula for the driving coherence (scheme 4f)
$$\rho _{n^{}n}^{(2)}V_{n^{}g}\rho _{gn}^{(1)}+\rho _{n^{}g}^{(1)}V_{gn}[\frac{1}{\mathrm{\Omega }_2+i\mathrm{\Gamma }_{n^{}g}}\frac{1}{\mathrm{\Omega }_1i\mathrm{\Gamma }_{ng}}]\frac{1}{\mathrm{\Omega }+i\mathrm{\Gamma }_{n^{}n}}=$$
$$=\frac{1}{(\mathrm{\Omega }_2+i\mathrm{\Gamma }_{n^{}g})(\mathrm{\Omega }_1i\mathrm{\Gamma }_{ng})}[1i\frac{\mathrm{\Gamma }_{nn^{}}\mathrm{\Gamma }_{n^{}g}\mathrm{\Gamma }_{ng}}{\mathrm{\Omega }+i\mathrm{\Gamma }_{nn^{}}}].$$
$`(36)`$
Here $`\mathrm{\Omega }_1=\omega _1\omega _{ng}`$, $`\mathrm{\Omega }_2=\omega _2\omega _{n^{}g}`$, $`\mathrm{\Omega }=\omega _2\omega _1\omega _{n^{}n}`$. At spontaneous relaxation, $`\mathrm{\Gamma }_{ij}=(\mathrm{\Gamma }_i+\mathrm{\Gamma }_j)/2`$, and resonance $`\mathrm{\Omega }=0`$ disappears. Collisions induce this resonance.
## 10 Acknowledgements
This work was supported in part by the International Science Foundation (Grants No 4000, No 4300), by the Russian Foundation for Fundamental Research (Grants N 93-02-03460, N 4300) and by the Krasnoyarsk Regional science Foundation, (Grant 4F 0123).
## 11 References.
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T.Ya. Popova, A.K. Popov, S.G. Rautian, R.I. Sokolovskii, Zh. Eksp. Teor. Fiz., Vol. 57, 850, 1969 (Engl.: JETP, Vol. 30, 466, 1970, quant-ph/0005094;
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$`d.`$ A.M. Bonch-Bruevich, V.A. Khodovoi, N.A. Chigir, Zh. Eksp. Teor. Fiz., Vol.67, 2069, 1974;
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5. $`a.`$ C. Cohen-Tannoudji, F. Hoffbeck, S. Reynaud, Opt. Commun., Vol. 27, 71, 1978;
$`b.`$ A.K. Popov, V.M. Shalaev, Opt. Commun., Vol. 35, 189, 1980.
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$`b.`$ E.B. Aleksandrov, G.I. Khvostenko, M.P. Chaika, Interference of atomic states (in Russ.), Nauka, Moscow, 1991.
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M.S. Feld , In: Fundamental and Applied Laser Physics : Proceedings of the Esfahan Symposium,
August 29 — September 5, 1971/ Ed. by M.S. Feld, A. Javan, N. Kurnit, N.Y., Willey, pp. 369 – 420, 1973;
$`b.`$ Th. Hänsch, P. Toschek, Z. Physik, Bd.236, 213, 1970.
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$`b.`$ M.P. Bondareva, et al, Opt. Spektrosk. Vol. 38, 219, 1975.
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and ref. therein.
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# Is MS1054-03 an exceptional cluster? A new investigation of ROSAT/HRI X-ray data
## 1. Introduction
MS1054-03 is the most distant cluster of galaxies found in the X-ray selected sample of the Einstein Extended Medium Sensitivity Survey (EMSS; Gioia et al., 1990). Due to the extreme sensitivity of the high–mass end of the cluster mass function to the density parameter, $`\mathrm{\Omega }_0`$, the existence of massive, virialized clusters at high z is of great cosmological significance (e.g. Oukbir & Blanchard 1992). The very existence of even a few massive clusters at redshifts approaching unity and beyond strongly argue against cosmological models with $`\mathrm{\Omega }_0=1`$.
There is a general consensus that MS1054-03 is such a massive cluster, based on its high X-ray temperature, $`kT=12.3_{2.2}^{+3.1}`$ keV, measured with ASCA (Donahue et al. 1998 hereafter D98), its high velocity dispersion (Tran et al. 1999) and intense weak lensing signal (Luppino & Kaiser 1997). On the other hand its apparent complex morphology (D98) might indicate that MS1054-03 has not yet reached an equilibrium state, which could cast some doubt on the dynamical mass estimates and on the relevance of the cluster for cosmological tests. Tran et al. (1999) emphasized the good agreement between the various mass estimates, but the large error bars must be noted.
In this paper we re-investigate the ROSAT/HRI observations of MS1054, including a new exposure taken in 1997. We first try to better understand the cluster morphology. A wavelet analysis of the HRI image is performed to unravel significant substructures and identify the main cluster component. A $`\beta `$–model is then fitted to the data. From the fit results, we estimate for the first time the gas mass of this cluster. A comparison of the gas mass fraction of MS1054-032, with the gas mass fraction of nearby clusters, provides a consistency check on the total mass estimate, assuming that the gas mass fraction in clusters does not evolve with redshift.
Throughout the paper we assume $`H_0=50`$ km/sec/Mpc, $`q_0=0.5`$, $`\mathrm{\Lambda }=0`$ and all quoted error bars are 1 $`\sigma `$ unless otherwise stated. At the redshift of the cluster, $`z=0.83`$, one arcmin corresponds to $`497\mathrm{h}_{50}^1`$ kpc.
## 2. Observations
MS1054-03 was observed by the ROSAT HRI (Trümper 1992) for 190,754 sec in total: an exposure of $`122`$ ksec in 1996, which was analyzed and presented in D98 and a pointing of 68 ksec taken in 1997. In this analysis we combined the exposures of 1996 and 1997 and we selected only HRI channels 1 to 7 (David et al. 1997) in order to maximize the signal-to-noise ratio. The background level is thus lowered by a factor of about 2. The HRI image smoothed with a Gauss filter is shown on Fig. 1 and looks similar to the image presented by D98.
## 3. Morphology
In order to remove noise and to identify the significant structures we performed a wavelet-analysis of the HRI raw image with the Multi-scale Vision Model (MVM) package (Rué & Bijaoui 1997). We explicitly took into account the Poisson statistics and the significance level was set to $`3\sigma `$. The algorithm is optimized for the detection of diffuse components by normalizing the wavelet coefficients by their energy (Rué & Bijaoui 1997, Arnaud et al. 2000).
The reconstructed image consists of point sources and one large scale diffuse component at the cluster position. To extract possible substructures, we reapplied the multi-scale analysis to the reconstructed diffuse source only. Two components are detected that are shown in Fig. 2, superimposed on the cluster optical image. The main component ($`90\%`$ of the flux) is centered at RA = $`10^h56^m58.47^s`$, Dec.= $`03\mathrm{deg}37\mathrm{}31.2\mathrm{}`$, about $`20\mathrm{}`$ ($`166\mathrm{h}_{50}^1`$ kpc) from the dominant galaxy (D98). We identify this structure with the main cluster component. It is very elliptical with an elongation in the east-west direction. The orientation follows the filamentary structure found in the optical and in the weak lensing map (Luppino & Kaiser 1997). The second component is compact and lies in the west, with an offset of 0.6 arcmin with respect to the main cluster center. It coincides with the brightest image peak, identified by D98, but its reconstructed flux is only $`10\%`$ of the main cluster flux.
This western source might be a subcluster gravitationally connected to the main cluster, or a projection effect due to a group of galaxies in the field-of-view (FOV) or even a point source. The available data are insufficient to settle this issue. There is a weak indication that the source is extended, from our comparison of the source surface brightness profile with the point spread function (PSF) of the ROSAT/HRI (David et al. 1997). However, due to the lack of bright sources in the FOV, we cannot correct for possible alignment uncertainties, and the extent might be due to errors in the aspect reconstruction or residual main cluster emission. Very close to the maximum of this source we find several cluster members so this structure might be a subgroup falling right onto the cluster. Finally it must be noticed that the reconstructed surface mass density by Luppino & Kaiser (1997) presents an extension in the west. Although it is unresolved, it might indicate that the western X–ray component does trace a gravitational potential.
## 4. Isothermal beta-model fits
We performed $`\beta `$-model fits to the data using spherically symmetric (1–D) and elliptical (2–D) $`\beta `$–model fits. We excluded from the fits circular regions corresponding to serendipitous point sources. In the 1–D case the X-ray surface brightness profile is given by:
$$S=S_0\left(1+r^2/r_c^2\right)^{3\beta +0.5}+B$$
(1)
where $`r_c`$ is the core radius, $`\beta `$ is a slope parameter and $`B`$ is the background level. The functional form for the elliptical $`\beta `$–model can be found in Neumann & Böhringer (1997).
### 4.1. Spherical model
We binned the data in concentric annuli with the center at the position determined from the above wavelet reconstruction. The $`\beta `$–model was convolved with the PSF of the ROSAT/HRI. The free parameters in the fit are $`S_0`$, $`r_c`$, $`\beta `$ and $`B`$. We first consider the main cluster component in a common way by excluding non-symmetric features in the data, i.e. a circular region around the substructure in the west. Fig. 3 shows the observed surface brightness profile and the best fit model. The reduced $`\chi ^2`$ value is $`\chi ^2=0.97`$, and the shape parameters are well constrained (see Tab.1). On the other hand if the substructure is not removed, a good fit is still obtained ($`\chi ^2=0.99`$) but the best fit shape parameters are unreasonably large ($`\beta =2.,r_c=840\mathrm{kpc}`$), and poorly constrained ($`\beta >1.1`$). An artificial increase of the best fit shape parameters is a well known effect of not excising a secondary sub-cluster in the $`\beta `$–model fit (Jones & Forman 1999). A good fit is still obtained due to the large errors on the observed profile, whereas no upper limit on $`\beta `$ can be set because the best fit core radius becomes similar to the maximum radius of cluster emission detected.
### 4.2. Elliptical model
We also fitted an elliptical isothermal $`\beta `$–model to the HRI image. Again we excise the substructure in the west as well as other weak unresovled sources in the FOV. Due to the limited statistics of the observation, we did not attempt to fit the whole image with a $`\beta `$–model and an additional component for the substructure, since the elliptical $`\beta `$–model has already 8 fit parameters.
The iso-contours of the best fit model are superimposed on the cluster image in Fig. 1. The fitting procedure is described in full detail in Neumann (1999). We binned the data into an image with a pixel size of $`5\mathrm{}\times 5\mathrm{}`$. The fit included all pixels less than 4.6 arcmin away from the central pointing position of $`10^h57^m00.00^s`$, and Dec.= $`03\mathrm{deg}37\mathrm{}12.0\mathrm{}`$ (J2000). The central position of the cluster was left free to vary.
As we are dealing with low number statistics in each image pixel, which shows non-Gaussian behavior, we smoothed the data with a Gauss filter ($`\sigma =10\mathrm{}`$) before fitting. The modeling accounts for the Gauss filtering and for the PSF (the fitted model is convolved with the Gauss filtered PSF). To calculate the errors of the $`\beta `$–model parameters we performed a Monte-Carlo analysis in which we added Poisson noise to the data and fitted the $`\beta `$–model subsequently. We performed 100 Monte-Carlo realizations.
The validity of this approach was discussed in Neumann & Böhringer (1997) and Neumann (1999). In order to see whether it is still free of systematics in this regime of extremely low signal-to-noise data (the central cluster intensity is only a factor of two higher than the background level) we simulated 100 realizations of the image of the best fit $`\beta `$–model cluster including background. The Poisson statistics were defined according to the length of the actual exposure time. We subsequently fitted a $`\beta `$–model to these artificial smoothed images, as for the real image. The results are satisfactory, as the mean output values are practically identical to the input values, with differences much smaller than the actual determined error bars for MS1054-03. The width of their distribution is slightly smaller than the errors determined from the real data.
### 4.3. Comparison of the 1–D and 2–D models
The values of the spherical and elliptical $`\beta `$–model parameters for the main cluster component together with their 1 $`\sigma `$ errors are given in Tab.1 and the two $`\beta `$–models profiles are plotted on Fig. 3. The best fit center of the elliptical model is only $`5\mathrm{}`$ away from the central position from the wavelet analysis, which we also choose as the center for the 1–D fit. As already mentioned, these central positions are also in excellent agreement with the position of the brightest cluster member. The 1–D and 2–D best fit parameters are consistent. The background level is slightly higher in the 1–D fit which explains the lower central value and the higher value for $`\beta `$. Nevertheless the $`68\%`$ confidence intervals for the fitted shape parameters show a large large region of overlap.
## 5. Luminosity
The background estimated from the 1–D model fit was subtracted from the radial profile. An additional $`5\%`$ systematic uncertainty was assumed on its level. The cluster emission is detected up to $`2\mathrm{}`$ or $`1\mathrm{h}_{50}^1\mathrm{Mpc}`$ at the $`68\%`$ confidence level. The observed count rate within this aperture is $`(9\pm 2)\times 10^3`$ cts/s, $`25\%`$ higher than the value given by D98 (subtracting additional sources, including the western structure for consistency). We recall that the count rate of the western substructure is only $`10\%`$ of the cluster flux. The count rates derived from the best fit 2–D model is consistent but higher $`11\times 10^3`$cts/s, a consequence of the lower best fit background, in this case.
The observed count rate was translated into luminosity, assuming a cluster temperature of 12.3 keV, as measured by D98 from ASCA data. The derived X–ray luminosity is $`L_\mathrm{X}=1.2\times 10^{45}\mathrm{h}_{50}^2`$ ergs/s (0.1–2.4 keV rest frame). The bolometric luminosity within $`2\mathrm{}`$ in radius, $`L_{\mathrm{bol}}=4.3\pm 0.9\times 10^{45}\mathrm{h}_{50}^2`$ergs/s, is consistent with the ASCA value of $`4.4\times 10^{45}\mathrm{h}_{50}^2`$ergs/s (D98). This might be fortuitous though since the later includes all components in the field. The bolometric luminosity is lower than the value expected from the $`L_\mathrm{X}`$$`T`$ relation of Arnaud & Evrard (1999), assuming this relation does not evolve with redshift. For a cluster with $`T=12.3_{2.2}^{+3.1}`$ keV this correlation predicts a luminosity of $`9.1_{4.0}^{+8.2}\times 10^{45}\mathrm{h}_{50}^2`$ergs/s. Up to now it is not clear whether the $`L_\mathrm{X}`$$`T`$ relation evolves with redshift or not (see for example Schindler 1999). Also, the total X-ray luminosity of MS1054 might be actually higher since the cluster might extend beyond the radius of detection. However, due to the low S/N ratio of this observation such extrapolation is very sensitive to the assumed background level and $`\beta `$ parameter.
## 6. Mass content
### 6.1. Mass estimates
The cluster gas mass profile $`M_{\mathrm{gas}}(r)`$ can be derived from the best fit $`\beta `$–model , given the observed temperature and $`N_H`$ values. As the emissivity in the ROSAT/HRI energy band is nearly insensitive to the temperature, the uncertainty on the gas mass is overwhelmingly dominated by the errors on the $`\beta `$–model parameters. The association of the western source with the cluster is unclear. However, in practice, this additional uncertainty has no significant impact on the gas mass estimate. If the substructure is included in the $`\beta `$–model analysis, the derived gas mass differs by less than $`5\%`$ from the value obtained when excising it; the derived mass is increased within the detection radius (as a consequence of the $`10\%`$ higher flux) and artificially decreased when one extrapolates the data to higher radii (due to the higher derived $`\beta `$ value). In the following we only consider the main cluster component and its corresponding $`\beta `$ model.
The total mass can be estimated with the isothermal $`\beta `$–model approach (BM), which is thought to be roughly valid even if the cluster is not fully in hydrostatic equilibrium (e.g Schindler 1996). The alternative is to employ the virial theorem (VT) at given density contrast, over the mean mass density of the Universe at the cluster redshift, normalized from numerical simulations (Evrard et al. 1996). The resulting $`M`$$`T`$ relation (see Appendix) depends both on redshift and on the cosmological parameters $`\mathrm{\Omega }_0`$ and $`\mathrm{\Lambda }`$ (Voit & Donahue 1998). We used the analytical expression from Bryan & Norman (1998).
We first fix the temperature to the best fit ASCA value of k$`T=12.3\mathrm{keV}`$ and consider an $`\mathrm{\Omega }_0=1`$ Universe. The uncertainty introduced by the errors on the temperature is discussed later. The statistical errors on $`M_{\mathrm{gas}}(r)`$ and on $`M_{\mathrm{BM}}(r)`$ due to the uncertainties in the $`\beta `$–model , are estimated following the general method described in Elbaz et al. (1995). The derived mass profiles are plotted in Fig. 4. Tab.2 summarizes the mass estimates at $`1\mathrm{h}_{50}^1\mathrm{Mpc}`$ (the maximum radius of detection) and $`1.65\mathrm{h}_{50}^1\mathrm{Mpc}`$ (the virial radius $`R_\mathrm{V}`$ for z=0.83 and k$`T=12.3\mathrm{keV}`$ using the Evrard et al. (1996) simulations).
The 2–D model gas mass at $`1\mathrm{h}_{50}^1\mathrm{Mpc}`$ is $`15\%`$ higher than the value derived from the 1–D model. This discrepancy, larger than the formal statistical uncertainties, is essentially due to the systematic difference in the background estimates, which dominates the error. As the derived $`\beta `$ value is consistently smaller in the 2–D fit, this discrepancy is amplified for extrapolated gas masses. It reaches $`25\%`$ at the virial radius. Similarly the total BM mass estimate which scales as $`\beta `$ is smaller for the 2–D model than for the 1–D model. Both are consistent with the VT estimate, which we will adopt in the following discussion. The gas mass is taken as the average of the 1–D and 2–D estimates. Their difference is an estimate of the systematic uncertainties that we add in quadrature to statistical ones. The gas mass fraction at the virial radius is thus $`f_{gas}=14\pm 3\%`$, for k$`T=12.3\mathrm{keV}`$ and $`\mathrm{\Omega }_0=1`$, the error coming from the uncertainty on the gas mass from the imaging data.
We performed the same analysis for the extreme values of the temperature, as allowed by the ASCA data at the $`90\%`$ confidence level (D98). For the lower limit on the temperature, k$`T=10.1\mathrm{keV}`$, the virial radius is $`1.58\mathrm{Mpc}`$ and the gas mass fraction within that radius is $`f_{gas}=16.5\pm 3\%`$. For k$`T=15.4\mathrm{keV}`$, we get $`R_\mathrm{V}=1.95\mathrm{Mpc}`$ and $`f_{gas}=11\pm 3\%`$.
In summary the gas mass fraction at the virial radius is $`f_{gas}=14[3,+2.5]\pm 3\%`$ for $`\mathrm{\Omega }_0=1`$. We have separated the uncertainties due to the errors on the temperature (in bracket), which affect essentially the total mass estimate, and the uncertainties from the imaging data, which only affect the gas mass estimate.
### 6.2. Comparison with nearby clusters
As mentioned above, a precise determination of the total mass of distant luminous clusters like MS1054-03 is crucial for the determination of $`\mathrm{\Omega }_0`$ based on cluster abundances at high redshifts. Unfortunately there are still important uncertainties on the total mass estimate: statistical errors due to the errors on the temperature measurement but also possible systematic errors if the temperature is a biased estimator of the cluster mass. In particular the gas temperature might differ from the true virial temperature if the cluster were not in hydrostatic equilibrium.
By considering the additional information on the gas mass, we can however perform a consistency check on the total mass estimate. From a study of the intrinsic dispersion in the $`L_\mathrm{X}`$$`T`$ relation, Arnaud & Evrard (1999) showed that the fractional variations of $`f_{gas}`$ at fixed cluster mass is very small. Since it is unlikely that the gas mass fraction evolves with redshift, the derived gas mass fraction of MS1054-03 should be consistent with the typical value of hot nearby clusters.
As a reference we first consider the best fit ASCA temperature (k$`T=12.3\mathrm{keV}`$) and a $`\mathrm{\Omega }_0=1`$ universe. The corresponding gas mass fraction of MS1054-03 is significantly smaller (at the 95$`\%`$ confidence level) than the value found by Arnaud & Evrard (1999) for hot nearby clusters, $`f_{gas}=20.1\pm 1.6\%`$. Obviously MS1054-03 might be a true outlier. Such outliers are rare but some, such as A1060, a cluster of exceptional low gas content in the local Universe (Arnaud & Evrard 1999), do exist. If this is not the case for MS1054, the low derived $`f_{gas}`$ value suggests that the actual virial temperature is lower than $`12.3\mathrm{keV}`$ and/or that the adopted cosmology is wrong. We examine each possibility in turn.
The derived gas mass fraction depends on the assumed values of the cosmological parameters $`\mathrm{\Omega }_0`$ and $`\mathrm{\Lambda }`$. In the BM approach, $`f_{gas}d_\mathrm{A}^{3/2}`$, where $`d_\mathrm{A}`$ is the angular distance (Pen 1997). In the VT approach adopted here the dependence of $`f_{gas}`$ is slightly different. $`f_{gas}`$ is both dependent on the geometry factor (the variation of $`d_\mathrm{A}`$) and on the normalization of the $`M_\mathrm{V}(\mathrm{or}R_\mathrm{V})`$$`T`$ relation. A detailed derivation of this dependence and necessary equations to compute the variation of $`f_{gas}`$ with $`\mathrm{\Omega }_0`$ are given in the Appendix. The variation of the gas mass fraction of MS1054-03 with $`\mathrm{\Omega }_0`$ is plotted on Fig. 5 for an open Universe ( $`\mathrm{\Omega }<1`$, $`\mathrm{\Lambda }=0`$) and a flat Universe ($`\mathrm{\Omega }_0+\mathrm{\Lambda }=1`$). For the best 1–D(2–D) $`\beta `$–model , the estimated gas mass fraction of MS1054-03 would be a factor of 1.2(1.3) larger if $`\mathrm{\Omega }_0=0.3`$, $`\mathrm{\Lambda }=0.0`$ and 1.4(1.5) times larger if $`\mathrm{\Omega }_0=0.3`$,$`\mathrm{\Lambda }=0.7`$. The gas mass fraction of MS1054-03 would in the later case be perfectly consistent with the local value, which would itself remain essentially unchanged ($`<5\%`$ increase).
The second possibility is that the virial temperature, and thus the virial mass, are smaller than given by the best fit ASCA value. The lower limit on the ASCA temperature (k$`T=10.1\mathrm{keV}`$) yields a gas mass fraction marginally consistent with the local value. If the virial temperature were actually even somewhat lower, k$`T=8\mathrm{keV}`$, the virial mass would decrease to $`0.84\times 10^{15}\mathrm{M}_{\mathrm{}}`$, the virial radius to 1.3 Mpc and the gas mass within that radius, $`1.7\times 10^{15}\mathrm{M}_{\mathrm{}}`$, would reach $`20\%`$ of the virial mass and be a perfect match with the local value. A virial temperature of $`8\mathrm{keV}`$ is formally excluded by ASCA measurements at the $`90\%`$ confidence level but would be in better agreement with the measured velocity dispersion of $`\sigma =1170\pm 150`$ kms/s (Tran et al. 1999). In that case MS1054-03 would fit perfectly in the $`\sigma T`$ relation established for $`z=0.190.55`$ clusters and the virial relation (Fig. 3 of Tran et al. 1999). The temperature would also be more consistent with the measured bolometric luminosity. A possible explanation for the measured temperature being significantly higher than the virial temperature (apart from contamination by a hard source like an absorbed AGN) would be the presence of shock waves in the gas if the cluster were indeed undergoing a recent merger.
## 7. Conclusion
Our wavelet analysis finds evidence for two components in the X–ray image of MS1054-03: a main diffuse component, with emission peaking at $`20\mathrm{}`$ from the brightest cluster member, and a compact substructure in the west, which coincides with the X-ray maximum in the ROSAT image. We emphasize that this brightest peak is not the centroid of the cluster. Indeed, when the western structure is excluded, the cluster emission is fitted well by a classical $`\beta `$–model , centered within $`20\mathrm{}`$ of the central cD galaxy. This indicates that we are then identifying the main cluster component.
Unlike previous attempts (D98, Ebeling et al. 1999), we obtained a good fit to the data with a $`\beta `$–model (low $`\chi ^2`$ value) and derived well constrained parameters. This is a natural consequence of the higher S/N ratio of our data (optimum channel selection and longer exposure time) and our identification and removal of the western structure from our fits.
There are indications from the X–ray data that the cluster is not fully relaxed. The core radius ($`400`$ kpc) and ellipticity are relatively high and the western substructure could be an in-falling group. This is not unusual, such mergers do exist in the nearby Universe. Actually MS1054-03 appears very similar to A521 at z=0.27 (Arnaud et al. 2000), where the brightest peak is also associated with the subcluster.
The gas mass fraction derived for the best fit ASCA temperature of k$`T=12.3\mathrm{keV}`$ is only consistent with the local value if $`\mathrm{\Omega }_0<1`$, a flat $`\mathrm{\Lambda }`$–dominated Universe being favored. The local value can be matched as well for $`\mathrm{\Omega }_0=1`$, provided that the actual virial temperature is close to the lower ASCA limit ($`10\mathrm{keV}`$), with an even lower value of $`8\mathrm{keV}`$ giving the best match. To decide between the two options requires a systematic analysis of a large sample of distant clusters. MS1054-03 would appear as an outlier in the second case (low virial temperature) and not in the first case (flat $`\mathrm{\Lambda }`$ dominated Universe). If the cluster’s actual mass is indeed lower than previously estimated, this might have consequences for the measurements of $`\mathrm{\Omega }_0`$ based on the abundance of massive clusters at high redshift.
Finally we want to stress that the current data allow us to constrain the physics of MS1054 only with large error bars. Therefore better data, in particular spectroscopic data from Chandra and XMM, are strongly needed.
## Acknowledgments
We want to thank Isabella Gioia and Megan Donahue for providing the optical image (see Fig. 2). We are grateful for discussions with Megan Donahue, Marshall Joy, Sandeep Patel and Jack Hughes.
## APPENDIX
In this Appendix we consider the gas mass fraction derived from X-ray data in the VT approach (see Sect.6.1) and examine how it depends on the assumed cosmological parameters, $`\mathrm{\Omega }_0`$ and $`\mathrm{\Lambda }`$.
The total mass $`M_\mathrm{V}`$ is estimated from the measured temperature and the mass-temperature relation derived from the virial theorem at given density contrast:
$$M_\mathrm{V}\left(\frac{\mathrm{\Delta }_\mathrm{c}(\mathrm{\Omega }(z),\mathrm{\Lambda })\mathrm{\Omega }_0}{\mathrm{\Omega }(z)}\right)^{1/2}(1+z)^{3/2}\left(\mathrm{k}T\right)^{3/2}$$
(1)
It depends on $`\mathrm{\Omega }_0`$ and $`\mathrm{\Lambda }`$, via the density contrast, $`\mathrm{\Delta }_\mathrm{c}(\mathrm{\Omega }(z),\mathrm{\Lambda })`$, and the density parameter of the universe at redshift z, $`\mathrm{\Omega }(z)`$. The analytical expression of $`\mathrm{\Delta }_c`$ and $`\mathrm{\Lambda }(z)`$ can be found in Bryan & Norman (1998). The corresponding virial radius scales as:
$$R_\mathrm{V}\left(\frac{\mathrm{\Delta }_\mathrm{c}(\mathrm{\Omega }(z),\mathrm{\Lambda })\mathrm{\Omega }_0}{\mathrm{\Omega }(z)}\right)^{1/2}(1+z)^{3/2}\left(\mathrm{k}T\right)^{1/2}$$
(2)
We assume that the gas density $`n_\mathrm{g}(r)`$, follows a $`\beta `$–model: $`n_\mathrm{g}(r)=n_{\mathrm{g},0}\left[1+(r/r_\mathrm{c})^2\right]^{3\beta /2}`$. X-ray imaging data provides the central surface brightness $`S_0`$, the slope parameter $`\beta `$ and the angular core radius $`\theta _c=r_\mathrm{c}/d_\mathrm{A}`$, where $`d_A`$ is the angular distance.
The emission measure along the line of sight through the cluster center can be derived from $`S_0`$, independently of any cosmological parameters :
$$EM_0=\frac{4\pi (1+z)^4S_0}{\mathrm{\Lambda }(T,z)}$$
(3)
where $`\mathrm{\Lambda }(T,z)`$ is the emissivity in the detector band, taking into account the interstellar absorption and the instrument spectral response.
Assuming that the X-ray atmosphere extends up to the virial radius, the central emission measure along the line of sight is linked to the gas density by $`EM_0_0^{R_\mathrm{V}}n_\mathrm{g}^2(r)𝑑r`$, whereas the gas mass within $`R_V`$ is $`M_{\mathrm{gas}}_0^{R_\mathrm{V}}n_\mathrm{g}(r)r^2𝑑r`$. By combining the two expressions, we can derive an expression for the gas mass that varies as:
$$M_{\mathrm{gas}}\alpha R_\mathrm{V}^{5/2}\sqrt{EM_0}Q(\beta ,x_\mathrm{c})$$
(4)
where we have introduced the form factor $`Q(\beta ,x_\mathrm{c})=n_\mathrm{g}/\sqrt{n_\mathrm{g}^2_{los}}`$. Here $`n_\mathrm{g}`$ is the average gas density within $`R_\mathrm{V}`$ and $`n_\mathrm{g}^2_{los}`$ is the average along the line–of–sight passing through the cluster center. For the $`\beta `$–model, this form factor depends on $`\beta `$ and the scaled core radius $`x_c=r_c/R_V`$:
$`Q(\beta ,x_\mathrm{c})`$ $`=`$ $`{\displaystyle \frac{3_0^1\left[1+(x/x_c)^2\right]^{3\beta /2}x^2𝑑x}{\sqrt{_0^1\left[1+(x/x_c)^2\right]^{3\beta }𝑑x}}}={\displaystyle \frac{\text{BETACF}(\frac{3}{2},\frac{3\beta }{2},\frac{1}{1+x_c^2})}{\sqrt{\text{BETACF}(\frac{1}{2},3\beta \frac{1}{2},\frac{1}{1+x_c^2})}}}`$ (5)
BETACF is the continued fraction entering the expression of the incomplete Beta function $`B_x(a,b)`$ (see Press et al. 1986) with:
$$B_x(a,b)=\frac{x^a(1x)^b}{a}\text{BETACF}(x,a,b)$$
(6)
From Eq.1, Eq.2, Eq.3 and Eq.4, the gas mass fraction, estimated from given X–ray data, depends on the assumed cosmological parameters $`\mathrm{\Omega }_0`$ and $`\mathrm{\Lambda }`$ as:
$`f_{\mathrm{gas}}`$ $`\alpha `$ $`R_\mathrm{V}^{3/2}Q(\beta ,\theta _\mathrm{c}d_\mathrm{A}/R_\mathrm{V})`$ (7)
where $`R_\mathrm{V}\left(\mathrm{\Delta }_\mathrm{c}(\mathrm{\Omega }(z),\mathrm{\Lambda })\mathrm{\Omega }_0/\mathrm{\Omega }(z)\right)^{1/2}`$ and $`d_\mathrm{A}`$ is the angular distance. Note that the variation with the cosmological parameters is not the same for all clusters. It depends on the cluster temperature and shape (via the $`\theta _\mathrm{c}d_\mathrm{A}/R_\mathrm{V}`$ factor).
The variation of $`f_{\mathrm{gas}}`$ with $`\mathrm{\Omega }_0`$, for an open Universe ($`\mathrm{\Lambda }=0`$) and a flat Universe ($`\mathrm{\Omega }_0+\mathrm{\Lambda }=1`$) is plotted on Fig. 5 for the X–ray best fit parameters of MS1054-03 ($`\mathrm{k}T=12.3\mathrm{keV},\beta =0.96,\theta _\mathrm{c}=0.89\mathrm{}`$).
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# Extremely charged static perfect fluid distributions with dilaton in curved spacetimes
## 1 Introduction
Recently, there has been much interest in the study of Majumdar-Papapetrou metrics , discribing the static equilibrium state of extremely charged black holes. For the static Einstein-Maxwell equation with charged dust as the external source of the fields, one can reduce the electrovacuum field equations to the Poisson equation in the flat space . In such a system, one can show that the charge and mass densities are equal.
In the low energy limit of string theory, the dilatonic forces as well as gravitational and electric forces are acting among charged matters as long-range forces. In this paper, we study the charged perfect fluid distributions which also couple to a dilaton field in static $`(N+1)`$ dimensional spacetimes. We find that field equations reduce to a non-linear type of Poisson equation and that Maxwell equation and an equation for a dilaton show the relation among the charge, mass and dilatonic charge densities. We also examine some simple exact solutions.
The organization of this paper is as follows. In the next section, we will show the action and the assumptions on the charged perfect fluid distributions in the static $`(N+1)`$ dimensional spacetimes, which we consider in the present paper. We reduce the field equations to the non-linear version of Poisson equation in section 3. We find some simple solutions and discuss them in section 4. Finally, section 5 is devoted to conclusion and discussion.
## 2 The model
The action for the fields which mediate long-range forces is
$$S=d^{N+1}x\frac{\sqrt{g}}{16\pi }\left[R\frac{4}{N1}(\varphi )^2e^{\frac{4a}{N1}\varphi }F^2\right],$$
(1)
where $`N`$ $`(N3)`$ denotes the dimension of space, $`R`$ is the scalar curvature and $`\varphi `$ is the dilaton field. $`F^2=F^{\mu \nu }F_{\mu \nu }`$ and $`F_{\mu \nu }`$ denotes the Maxwell field strength. The dilaton coupling to the Maxwell term is governed by a constant $`a`$. The Newton constant is normalized to unity.
Incorporating coupling to matter, we obtain our basic equations:
$$G_{\mu \nu }\frac{4}{N1}\left[_\mu \varphi _\nu \varphi \frac{1}{2}g_{\mu \nu }(\varphi )^2\right]e^{\frac{4a}{N1}\varphi }\left[2F_{\mu \nu }^2\frac{1}{2}g_{\mu \nu }F^2\right]=8\pi T_{\mu \nu },$$
(2)
$$\frac{8}{N1}^2\varphi +\frac{4a}{N1}e^{\frac{4a}{N1}\varphi }F^2=4\pi \frac{8a}{N1}\rho _{dil},$$
(3)
$$_\mu \left[e^{\frac{4a}{N1}\varphi }F^{\mu \nu }\right]=4\pi j^\nu ,$$
(4)
where $`G_{\mu \nu }`$ is the Einstein tensor. The energy momentum tensor $`T_{\mu \nu }`$ for a perfect fluid is given by
$$T_{\mu \nu }=(\rho +p)u_\mu u_\nu +pg_{\mu \nu },$$
(5)
where $`\rho `$ is the energy density and $`u^\mu `$ is the four velocity. The electric current vector $`j^\mu `$ is defined as
$$j^\mu =\rho _eu^\mu ,$$
(6)
where $`\rho _e`$ is the charge density. We have introduced the dilatonic charge density $`\rho _{dil}`$ in the right hand side of the equation for a dilaton field.
## 3 Deriving the non-linear version of Poisson equation
We assume that the fluid is static and the metric of the static spacetime takes the form:
$$ds^2=U^2dt^2+U^{\frac{2}{N2}}\stackrel{~}{g}_{ij}dx^idx^j,$$
(7)
where $`i,j=1,\mathrm{},N`$, and both the background metric $`\stackrel{~}{g}_{ij}`$ and $`U`$ depend only on the space-like coordinates $`x^i`$.
The Ricci and Einstein tensor components derived from the metric (7) are given by
$$R_{00}=U^{2\frac{2}{N2}}\stackrel{~}{}_l\left(\frac{\stackrel{~}{}^lU}{U}\right),$$
(8)
$$G_{ij}=\frac{N1}{N2}\frac{\stackrel{~}{}_iU}{U}\frac{\stackrel{~}{}_jU}{U}+\frac{1}{2}\frac{N1}{N2}\stackrel{~}{g}^{kl}\frac{\stackrel{~}{}_kU}{U}\frac{\stackrel{~}{}_lU}{U}\stackrel{~}{g}_{ij}+\stackrel{~}{G}_{ij},$$
(9)
where $`\stackrel{~}{}_i`$ denotes the $`N`$ dimensional covariant derivative in terms of $`\stackrel{~}{g}_{ij}`$. $`\stackrel{~}{G}_{ij}`$ is constructed from $`\stackrel{~}{g}_{ij}`$.
Here we should assume that there is only the electric field, namely $`F_{0i}0`$ and the others are set to be zero. Then we get
$`F_{00}^2`$ $`=`$ $`U^{\frac{2}{N2}}\stackrel{~}{g}^{kl}F_{0k}F_{0l},`$
$`F_{ij}^2`$ $`=`$ $`U^2F_{0i}F_{0j},`$
$`F^2`$ $`=`$ $`2U^{2\frac{2}{N2}}\stackrel{~}{g}^{kl}F_{0k}F_{0l}.`$
In addition, we also put an assumption on the dilatonic field:
$$e^{\frac{4a}{N2}\varphi }=U^{2\alpha },$$
(10)
where $`\alpha `$ is a constant.
Then the $`(00)`$ component in the left hand side of Eq. (2) becomes
$$R_{00}e^{\frac{4a}{N1}\varphi }\left[2F_{00}^2\frac{1}{N1}g_{00}F^2\right]=U^{2\frac{2}{N2}}\stackrel{~}{}^2\mathrm{ln}U2U^{2\alpha \frac{2}{N2}}\frac{N2}{N1}\stackrel{~}{g}^{kl}F_{0k}F_{0l},$$
(11)
while the $`(ij)`$ component in the left hand side of Eq. (2) is
$`G_{ij}`$ $``$ $`{\displaystyle \frac{4}{N1}}\left[_i\varphi _j\varphi {\displaystyle \frac{1}{2}}g_{ij}(\varphi )^2\right]e^{\frac{4a}{N1}\varphi }\left[2F_{ij}^2{\displaystyle \frac{1}{2}}g_{ij}F^2\right]`$ (12)
$`=`$ $`{\displaystyle \frac{N1}{N2}}{\displaystyle \frac{\stackrel{~}{}_iU}{U}}{\displaystyle \frac{\stackrel{~}{}_jU}{U}}(N1){\displaystyle \frac{\alpha ^2}{a^2}}{\displaystyle \frac{\stackrel{~}{}_iU}{U}}{\displaystyle \frac{\stackrel{~}{}_jU}{U}}+2U^{2\alpha +2}F_{0i}F_{0j}`$
$`+{\displaystyle \frac{1}{2}}{\displaystyle \frac{N1}{N2}}\stackrel{~}{g}^{kl}{\displaystyle \frac{\stackrel{~}{}_kU}{U}}{\displaystyle \frac{\stackrel{~}{}_lU}{U}}\stackrel{~}{g}_{ij}+{\displaystyle \frac{1}{2}}(N1){\displaystyle \frac{\alpha ^2}{a^2}}\stackrel{~}{g}^{kl}{\displaystyle \frac{\stackrel{~}{}_kU}{U}}{\displaystyle \frac{\stackrel{~}{}_lU}{U}}\stackrel{~}{g}_{ij}U^{2\alpha +2}\stackrel{~}{g}^{kl}F_{0k}F_{0l}\stackrel{~}{g}_{ij}`$
$`+\stackrel{~}{G}_{ij}.`$
Here we should suppose that $`U^{\alpha +1}=V`$, then Eq. (11) is changed into
$`R_{00}e^{\frac{4a}{N1}\varphi }\left[2F_{00}^2{\displaystyle \frac{1}{N1}}g_{00}F^2\right]`$ $`=`$ $`U^{2\frac{2}{N2}}[{\displaystyle \frac{1}{\alpha +1}}{\displaystyle \frac{\stackrel{~}{}^2V}{V}}{\displaystyle \frac{1}{\alpha +1}}\stackrel{~}{g}^{kl}{\displaystyle \frac{\stackrel{~}{}_kV}{V}}{\displaystyle \frac{\stackrel{~}{}_lV}{V}}`$ (13)
$`+2{\displaystyle \frac{N2}{N1}}V^2\stackrel{~}{g}^{kl}F_{0k}F_{0l}].`$
In order that the second term is canceled by the third one in the right hand side of Eq. (13), we adopt
$$F_{0k}=\pm \sqrt{\frac{N1}{2(N2)}}\sqrt{\frac{1}{\alpha +1}}\frac{\stackrel{~}{}_kV}{V^2}.$$
(14)
On the other hand, assuming Eq. (14), we can reduce Eq. (12) to
$`G_{ij}`$ $``$ $`{\displaystyle \frac{4}{N1}}\left[_i\varphi _j\varphi {\displaystyle \frac{1}{2}}g_{ij}(\varphi )^2\right]e^{\frac{4a}{N1}\varphi }\left[2F_{ij}^2{\displaystyle \frac{1}{2}}g_{ij}F^2\right]`$ (15)
$`=`$ $`\left({\displaystyle \frac{1}{N2}}{\displaystyle \frac{1}{\alpha +1}}{\displaystyle \frac{\alpha ^2}{a^2}}{\displaystyle \frac{1}{\alpha +1}}+{\displaystyle \frac{1}{N2}}\right){\displaystyle \frac{N1}{\alpha +1}}{\displaystyle \frac{\stackrel{~}{}_iU}{U}}{\displaystyle \frac{\stackrel{~}{}_jU}{U}}`$
$`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{1}{N2}}{\displaystyle \frac{1}{\alpha +1}}{\displaystyle \frac{\alpha ^2}{a^2}}{\displaystyle \frac{1}{\alpha +1}}+{\displaystyle \frac{1}{N2}}\right){\displaystyle \frac{N1}{\alpha +1}}\stackrel{~}{g}^{kl}{\displaystyle \frac{\stackrel{~}{}_kU}{U}}{\displaystyle \frac{\stackrel{~}{}_lU}{U}}\stackrel{~}{g}_{ij}`$
$`+\stackrel{~}{G}_{ij}.`$
In order to eliminate the first and second terms in the right hand side of Eq. (15), we take
$$\alpha =\frac{a^2}{N2}.$$
(16)
Consequently, we reduce the left hand side of Eq. (2) to the following equations:
$$R_{00}e^{\frac{4a}{N1}\varphi }\left[2F_{00}^2\frac{1}{N1}g_{00}F^2\right]=U^{2\frac{2}{N2}}\frac{N2}{N2+a^2}\frac{1}{V}\stackrel{~}{}^2V,$$
(17)
$$G_{ij}\frac{4}{N1}\left[_i\varphi _j\varphi \frac{1}{2}g_{ij}(\varphi )^2\right]e^{\frac{4a}{N1}\varphi }\left[2F_{ij}^2\frac{1}{2}g_{ij}F^2\right]=\stackrel{~}{G}_{ij}.$$
(18)
We should remember that
$$e^{\frac{4a}{N1}\varphi }=U^{\frac{2a^2}{N2}}=V^{\frac{2a^2}{N2+a^2}},$$
(19)
$$F_{0k}=\pm \sqrt{\frac{N1}{2(N2+a^2)}}\frac{\stackrel{~}{}_kV}{V^2}.$$
(20)
Finally, using Eqs. (5), (6) and (17-20), we reduce the field equations (2), (3) and (4) simply to the following equations:
$$\stackrel{~}{}^2V+8\pi \frac{N2+a^2}{N1}V^{\frac{N+a^2}{N2+a^2}}(\rho +\frac{N}{N2}p)=0,$$
(21)
$$\stackrel{~}{R}_{ij}=\frac{16\pi p}{N2}V^{\frac{2}{N2+a^2}}\stackrel{~}{g}_{ij},$$
(22)
$$\rho _{dil}=\rho +\frac{N}{N2}p,$$
(23)
$$\rho _e=\pm e^{\frac{2a}{N1}\varphi }\sqrt{\frac{2(N2+a^2)}{N1}}\left(\rho +\frac{N}{N2}p\right).$$
(24)
Therefore, these equations represent the Einstein, Maxwell and dilaton equations.
Here we think about Eq. (24) for the dust case $`(p=0)`$. The action for particles, of which coordinates are denoted by $`x^\mu `$, can be written as:
$$I=\underset{a}{}𝑑s_a\left[m_ae^{\frac{2a}{N1}\varphi }+e_aA_\nu \frac{dx_a^\nu }{ds_a}\right],$$
(25)
where $`m_a`$ and $`e_a`$ stand for the mass and electric charges of the particles. Suppose that the distribution of these particles represents the matter densities. One can find that the dilatonic charge density is proportional to the charge density. Thus, we can recognize that the relationship between the charge density and the mass density is $`\rho _e\pm e^{\frac{2a}{N1}\varphi }\rho `$, because each electric charge $`e_a`$ is a constant. In the next section, we discuss the some explicit solutions of Eq. (21).
## 4 Exact solutions
For the dust case $`(p=0)`$, we find some simple exact solutions of Eq. (21), which do not have the singularities. When spherical symmetry is assumed, the non-linear version of Poisson equation takes the following form:
$$\frac{d^2V}{dr^2}+\frac{N1}{r}\frac{dV}{dr}+8\pi \rho \frac{N2+a^2}{N1}V^{\frac{N+a^2}{N2+a^2}}=0.$$
(26)
If we put the following condition on the energy density:
$$\rho =\frac{A}{8\pi }\frac{N1}{N2+a^2}V^{\frac{N+a^2}{N2+a^2}},$$
(27)
we can find that the solution is
$$V(r)=B\frac{Ar^2}{6N}.$$
(28)
Here $`A`$ and $`B`$ are constants.
We show that the energy density $`\rho `$ for a certain value of total mass plotted against $`r`$ for $`a^2=0`$, $`a^2=\frac{N1}{2N}`$, $`a^2=1`$ and $`a^2=N`$ in Fig. 1(a) in the case of $`N=3`$. In Fig. 1(b), $`1U^2`$ is plotted against $`r`$ for the same coupling constants. Here the energy density is matched to the one for the vacuum solution at $`r=2`$. Fig. 2 is drawn with the same conditions of Fig. 1, except for $`N=5`$ and Fig. 3 is also, except for $`N=9`$.
If we put another condition:
$$\rho =\frac{C^2}{8\pi }\frac{N1}{N2+a^2}V^{\frac{2}{N2+a^2}},$$
(29)
the solution is
$$V(r)=D\frac{J_{(N2)/2}(Cr)}{r^{(N2)/2}}.$$
(30)
Here $`C`$ and $`D`$ are constants, and $`J_\nu (z)`$ is the Bessel function. If we choose $`N=3`$ and $`a^2=0`$, then we can obtain the same results of Gürses .
We show that the energy density $`\rho `$ for a certain value of total mass plotted against $`r`$ for $`a^2=0`$, $`a^2=\frac{N1}{2N}`$, $`a^2=1`$ and $`a^2=N`$ in Fig. 4(a) in the case of $`N=3`$. In Fig. 4(b), $`1U^2`$ is plotted against $`r`$ for the same coupling constants. Here the energy density is matched to the one for the vacuum solution at $`r=2`$. Fig. 5 is drawn with the same conditions of Fig. 4, except for $`N=5`$ and Fig. 6 is also, except for $`N=9`$.
Varela considered the case that Eq. (26) can be reduced to the sine-Gordon equation . Using the new radial coodinate $`\tau =\frac{1}{r^{N2}}`$ to rearrange Eq. (26), we obtain:
$$\frac{d^2V}{d\tau ^2}+8\pi \rho \frac{N2+a^2}{(N2)^2(N1)}\tau ^{\frac{2(N1)}{N2}}V^{\frac{N+a^2}{N2+a^2}}=0.$$
(31)
If we assume
$$\rho =\frac{E^2}{8\pi }\frac{(N2)^2(N1)}{N2+a^2}\tau ^{\frac{2(N1)}{N2}}(\mathrm{sin}V)V^{\frac{N+a^2}{N2+a^2}},$$
(32)
then, Eq. (31) reduces to the sine-Gordon equation
$$\frac{d^2V}{d\tau ^2}+E^2\mathrm{sin}V=0,$$
(33)
which has the solutions
$$V(\tau )=2\mathrm{arcsin}[\mathrm{tanh}(E\tau +F)]+2n\pi ,$$
(34)
where $`n`$ is an arbitrary integer, $`F`$ is an integration constant, and $`E`$ is assumed to be positive. We consider only the case $`n=0`$. If we choose the integration constant $`F`$ for
$$F=\frac{1}{2}\mathrm{ln}\left[\frac{1+\mathrm{sin}(1/2)}{1\mathrm{sin}(1/2)}\right],$$
(35)
then the spacetime corresponding to Eq. (34) and Eq. (35) becomes asmptotically flat .
We show that the energy density $`\rho `$ for a certain value of total mass plotted against $`r`$ for $`a^2=0`$, $`a^2=\frac{N1}{2N}`$, $`a^2=1`$ and $`a^2=N`$ in Fig. 7(a) in the case of $`N=3`$. In Fig. 7(b), $`1U^2`$ is plotted against $`r`$ for the same coupling constants.
Here, in these figures, we find that the energy density decreases as the coupling constant $`a^2`$ increases. We also find that the difference between the energy densities gets narrow for the various values of $`a^2`$ and the contrast (i.e., the difference between the energy density at $`r=0`$ and the one at $`r=2`$) decreases as the dimension of space $`N`$ increases.
## 5 Conclusion and discussion
In this paper, we have investigated charged static perfect fluid distributions with the dilaton field in the frame-work of general relativity. As shown in section 3, the Einstein equations have reduced to the non-linear version of Poisson equation, and the Maxwell equation and the equation for the dilaton have implied the relation among the charge, mass and dilatonic charged densities. For the dust case, one can find that the relationship between the charge density and the mass density is $`\rho _e\pm e^{\frac{2a}{N1}\varphi }\rho `$, because, for point particles, the dilaton does not couple to the electric charge but to the mass.
In section 4, we have found simple exact solutions of Eq. (21) corresponding to certain energy densities. We have found that the energy density decreases as the coupling constant $`a^2`$ increases. We have found that the energy density decreases as the coupling constant $`a^2`$ increases. We have found that the difference between the energy densities gets narrow for the various values of $`a^2`$ and the contrast (i.e., the difference between the energy density at $`r=0`$ and the one at $`r=2`$) decreases as the dimension of space $`N`$ increases.
We have not yet dealt with Eq. (21) on the condition for $`p0`$. Recently, Ida found some exact charged solutions in this situation . We will study the non-zero pressure case with a dilaton field in $`(N+1)`$ dimensions. Also we have not yet considered the case of $`N=2`$, which we have only thought of the equilibrium between the dilatonic attractions and the electric repulsions. We must continue to make every effort to study these situations.
## Acknowledgement
The authors would like to thank N. Kan, M. Ooho and T. Watabe for useful advice and support.
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# I. Introduction
## I. Introduction
Observations of active galactic nuclei and other massive objects have increased the interest in equilibrium plasma models within a Schwarzschild or Kerr geometry. Even if the metric is given and all physical quantities are assumed to be independent of a toroidal angle $`\phi `$ and of time $`t`$, there is still some effort needed to reduce the whole set of equations for the electromagnetic field and the fluid quantities. This has been done within ideal magnetohydrodynamics (MHD) by several authors with various degrees of completeness and sophistication. A characteristic feature of the MHD description is the conservation of magnetic flux in a system co-moving with the center-of-mass four-velocity $`u^\mu `$ (of ions, essentially); this equation (“Ohm’s law”, Eq. (2) below) has been “integrated” for general stationary and axisymmetric systems by Bekenstein and Oron who give also a basic discussion and some historical background of general-relativistic MHD. As a result, one can represent the electromagnetic field tensor $`F_{\mu \nu }`$ completely by the particle flux $`nu^\mu `$, the two Killing vectors associated with the translational symmetry in $`t`$ and $`\phi `$, and two “flux functions” which are constant along the poloidal stream lines (Eq. (17) below). Further reductions of the whole set of MHD equations and discussions of the astrophysical background have been given by Camenzind , Mobarry and Lovelace , Nitta, Takahashi and Tomimatsu , and Beskin and Par’ev . Thus one arrives, as in the non-relativistic case, at a single potential equation for the magnetic flux function $`\mathrm{\Psi }`$ (the covariant toroidal component $`A_\phi `$ of the vector potential $`𝑨`$), together with some constraints.
One problem with these MHD models is the large number of arbitrary flux functions – no dissipation mechanism has been included to reduce this number –, and a fluid picture may be questionable if the collision frequencies are too low. This, however, is not our point here, and we admit an ideal fluid picture on reasons of simplicity. The question is, however, whether the usual MHD theory is consistent for a really rotating black hole . In this case the metric is necessarily non-diagonal; the relevant element $`g_{t\phi }`$ is the invariant scalar product of the two Killing vectors mentioned above , so it can by no means be transformed to zero. On the other hand, it couples the components $`A_t`$ and $`A_\phi `$ in Ampère’s law which will, in general, contradict the strong coupling of $`A_t`$ and $`A_\phi `$ in the MHD models (here they have to be functions of each other). So we are led to the question of how to reduce the set of multifluid plasma equations and Maxwell’s equations for stationary axisymmetric systems without the discrepancy mentioned above.
In Section II we give a short account of the work of Bekenstein and Oron, leading to the strong coupling of $`A_t`$ and $`A_\phi `$. The same solution, however, can be used for any generalized Helmholtz equation in ideal fluid descriptions. Two examples of these Helmholtz equations are derived in Sections III and IV, respectively: One refers to the ordinary non-relativistic MHD equation (momentum balance), the other to the relativistic multifluid plasma, assuming a constant temperature of all species. The latter example is used for the reduction problem in Section V. Here it is shown that the multifluid equations are reduced to one single potential equation for each species, and the poloidal components of Ampère’s equation are integrated. A numerical solution of this full set of four potential equations would be facilitated by the existence of a variational principle; the solution could then be approximated with finite elements by minimizing the corresponding functional. So it may be interesting that such a functional exists, as is shown in Section VI. Finally we discuss the results in Section VII.
## II. Flux conservation in axisymmetric plasmas
Let us assume a plasma configuration where all physical quantities, including the metric tensor components $`g_{\mu \nu }`$, are independent of time $`t`$ and of a toroidal angle $`\phi `$. We choose a coordinate system $`(x^\mu )`$ with $`x^0=ct(c=1)`$, $`x^1=\phi `$, and $`x^2,x^3`$ some poloidal coordinates. Assuming in addition that all physical quantities are invariant to the simultaneous inversion of $`t`$ and $`\phi `$ – which is reasonable for any rotating equilibrium – the most general line element $`ds`$ can be represented as follows :
$`(ds)^2=g_{rs}dx^rdx^s+g_{ab}dx^adx^b,`$ (1)
where the indices $`r,s`$ run from $`0`$ to $`1`$, and $`a,b`$ from $`2`$ to $`3`$. We want to determine an electromagnetic field tensor $`F_{\mu \nu }`$ which is consistent with Eq. (1), and which, in addition, obeys the condition of magnetic flux conservation in a medium moving with the Eulerian four-velocity $`u^\mu `$. The latter condition is familiar from ideal magnetohydrodynamics (MHD) and means that a certain four-vector $`E_\mu `$ (“co-moving electric field”) should vanish:
$`E_\mu F_{\mu \nu }u^\nu =0.`$ (2)
The field tensor $`F_{\mu \nu }`$, of course, should solve the homogeneous Maxwell equations; therefore it can be written as the curl of a four-potential $`A_\mu `$ in the usual manner:
$`F_{\mu \nu }=A_{\nu ,\mu }A_{\mu ,\nu },`$ (3)
where $`()_{,\mu }`$ is the partial derivative of the quantity in brackets with respect to $`x^\mu `$. Finally we use the continuity equation
$`(\sqrt{g}nu^\mu )_{,\mu }=0,`$ (4)
with $`g`$ the determinant of $`g_{\mu \nu }`$ and $`n`$ the particle number density in the local inertial rest frame.
To solve Eq. (2) we remember that any skew-symmetric tensor $`F_{\mu \nu }`$ can be represented by two four-vectors, $`E_\mu `$ and $`B_\mu `$ (“co-moving magnetic field”) ; denoting all components in the local inertial rest frame by primes (with $`u^0=c=1;u^i=0`$ for $`i=1,2,3`$) we define
$`E_0^{}=0`$ ; $`E_i^{}=F_{i0}^{},i=1,2,3`$
$`B_0^{}=0`$ ; $`B_i^{}=F_{jk}^{},`$
where in the last equation $`(i,j,k)`$ is a cyclic permutation of $`(1,2,3)`$. In this coordinate system we use the Minkowski metric
$`(g_{\mu \nu }^{})=diag(1,1,1,1),`$
and $`E^i=E_i^{},B^i=B_i^{}`$ are the local electric and magnetic fields, respectively. It is then easily seen that the covariant representation of $`E_\mu `$ in the laboratory frame with any velocity field $`u^\mu (x)`$ is given by the left part of Eq. (2), while $`B_\mu `$ is given by
$`B_\mu ={\displaystyle \frac{1}{2}}\epsilon _{\mu \nu \varrho \sigma }u^\nu F^{\varrho \sigma },`$ (5)
where the totally antisymmetric Levi-Cività tensor $`\epsilon _{\mu \nu \varrho \sigma }`$ includes a factor $`\sqrt{g}`$ in order to be a tensor. Both vectors $`E_\mu ,B_\mu `$ are obviously orthogonal to $`u^\mu `$; this corresponds to $`2\times 3`$ independent components. Therefore we can represent the six independent elements of $`F_{\mu \nu }`$ by these vector components, namely:
$`F_{\mu \nu }=E_\mu u_\nu E_\nu u_\mu \epsilon _{\mu \nu \varrho \sigma }u^\varrho B^\sigma .`$ (6)
The covariant equation (6) can easily be proved by writing it in the local inertial rest frame. Eq. (6) is generally valid; in the case of flux conservation according to Eq. (2) it simplifies to
$`F_{\mu \nu }=\epsilon _{\mu \nu \varrho \sigma }u^\varrho B^\sigma .`$ (7)
Obviously the field tensor $`F_{\mu \nu }`$ is then orthogonal not only to $`u^\mu `$, as required by Eq. (2), but also to $`B^\mu `$:
$`B^\mu F_{\mu \nu }=0.`$ (8)
The remaining task is now to construct $`B^\mu `$ for axisymmetric equilibria.
For this purpose we consider the two Killing vectors $`(\xi ^\mu )`$ associated with these symmetries, namely:
$`(\xi ^\mu )=`$ $`(k^\mu )`$ $`(1,0,0,0),`$
$`(\xi ^\mu )=`$ $`(m^\mu )`$ $`(0,1,0,0).`$
In both cases they lead to a vanishing partial derivative of any physical quantity $`A`$ in the direction of $`\xi ^\mu `$:
$`\xi ^\mu A_{,\mu }=0.`$ (9)
Let us assume that all components $`A_\nu `$ of the vector potential share this property (though a gauge transformation could destroy it). Then Eq. (3) leads to
$`F_{r\nu }\xi ^\mu F_{\mu \nu }=\xi ^\mu A_{\mu ,\nu }=A_{r,\nu },`$ (10)
where $`r=0`$ for $`\xi ^\mu =k^\mu `$ and $`r=1`$ for $`\xi ^\mu =m^\mu `$. The right-hand side of this expression is obviously non-zero only for $`\nu =a=2`$ or $`3`$. Differentiating Eq. (10) with respect to $`x^b`$, $`b=2`$ or $`3`$, but $`ba`$, we find
$`F_{ra,b}=A_{r,ab}=F_{rb,a}.`$ (11)
To obtain the same expressions from Eq. (7) we write the factor $`\sqrt{g}`$ of $`\epsilon _{\mu \nu \varrho \sigma }`$ explicitely, with the remaining constant permutation symbol $`\stackrel{~}{\epsilon }_{\mu \nu \varrho \sigma }`$:
$`\epsilon _{\mu \nu \varrho \sigma }=\sqrt{g}\stackrel{~}{\epsilon }_{\mu \nu \varrho \sigma }.`$
Then our integrability condition from Eqs. (7) and (11) reads as follows:
$`0`$ $`=`$ $`F_{ra,b}F_{rb,a}`$
$`=`$ $`\stackrel{~}{\epsilon }_{rsab}[\left(\sqrt{g}u^sB^b\right)_{,b}\left(\sqrt{g}u^bB^s\right)_{,b}`$
$`+\left(\sqrt{g}u^sB^a\right)_{,a}\left(\sqrt{g}u^aB^s\right)_{,a}].`$
The right-hand side of this equation is understood with fixed and mutually different values for $`r,s,a`$ and $`b`$. We can re-write it by restoring the summation convention with respect to the index $`a`$:
$`\left(\sqrt{g}u^sB^a\right)_{,a}\left(\sqrt{g}u^aB^s\right)_{,a}=0;s=0,1.`$ (12)
The general solution $`B^\mu `$ of Eq. (12) with $`B^\mu u_\mu =0`$ can be written as follows:
$`B^\mu =\alpha u^\mu +b^\mu ;\alpha b^\mu u_\mu ,`$
where $`b^\mu `$ solves the same Eq. (12) as $`B^\mu `$. Since it is linear in $`b^\mu `$ we can solve it separately for the poloidal part $`b^a`$ and for $`b^s`$. Ignoring now $`b^a`$ we can solve Eq. (12) for $`b^s`$ if the poloidal flow $`u^a`$ is not identically zero. It is useful to represent $`b^s`$ by a linear combination of the Killing vectors, namely:
$`b^\mu `$ $`=`$ $`nCK^\mu ,`$ (13)
$`K^\mu `$ $``$ $`k^\mu +\beta m^\mu ,`$ (14)
where the coefficients $`C`$ and $`\beta `$ have to be determined suitably. The factor $`n`$ has been included in order to take advantage of the mass conservation law, Eq. (4), where the replacement $`\mu a`$ is allowed due to the symmetries. Then we have a solution of Eq. (12) if $`C`$ and $`\beta `$ are constant along the poloidal stream lines (“flux functions”):
$`u^aC_{,a}`$ $`=`$ $`0=u^a\beta _{,a},`$ (15)
$`B^\mu `$ $`=`$ $`nC\left[K^\mu (K^\lambda u_\lambda )u^\mu \right],`$ (16)
$`F_{\mu \nu }`$ $`=`$ $`nC\epsilon _{\mu \nu \varrho \sigma }u^\varrho K^\sigma .`$ (17)
The last equation has been obtained by inserting Eq. (16) into Eq. (7). While $`C`$ is an arbitrary flux function, $`\beta `$ is fixed by the condition that $`F_{\mu \nu }`$ is perpendicular to both $`u^\mu ,B^\mu `$ or $`u^\mu ,K^\mu `$. Using Eqs. (2), (8), (10), (14) and (16) we find
$`0=K^\mu F_{\mu \nu }=(A_{0,\nu }+\beta A_{1,\nu }).`$ (18)
This equation can only be fulfilled if $`A_0`$ and $`A_1`$ are functions of each other and constant along the poloidal stream lines. This is indeed the case as can now easily be shown from Eq. (2) for $`\mu =r`$ and Eq. (10):
$`0=F_{r\nu }u^\nu =A_{r,a}u^a.`$ (19)
This completes our particular solution for $`B^\mu `$ and $`F_{\mu \nu }`$ if $`b^a=0`$. It coincides with the result of Ref. (their $`A`$ is our $`\beta `$). Why is a poloidal part of $`b^\mu `$ not possible? Our symmetry requires $`F_{rs}=0`$ according to Eq. (10); from Eq. (7) we have
$`F_{rs}=\epsilon _{rs\varrho \sigma }u^\varrho B^\sigma =\epsilon _{rsab}u^aB^b.`$
Here we see that for a non-zero poloidal flow $`u^a`$ the poloidal part of $`B^\mu `$ must be proportional to $`u^\mu `$, otherwise $`F_{rs}`$ would be non-zero. Therefore the solution as given in Eqs. (16) and (17) is unique, up to the specification of two flux functions, $`C`$ and $`\beta `$.
Sometimes it is useful to re-write the result of this section in particular poloidal coordinates as defined by the poloidal stream lines, $`\mathrm{\Psi }`$ = const., and an angle-like coordinate $`\theta `$ varying along the poloidal stream lines:
$`x^2=\mathrm{\Psi };x^3=\theta .`$
From Eq. (19) we know that $`A_r=A_r(\mathrm{\Psi })`$, so we may identify one of both components with $`\mathrm{\Psi }`$ itself, while the other component defines $`\beta `$ according to Eq. (18), e.g.:
$`A_1\mathrm{\Psi };\beta =dA_0(\mathrm{\Psi })/d\mathrm{\Psi }.`$ (20)
In this coordinate system we have from Eq. (10):
$`F_{13}=0;F_{12}=1.`$
Inserting here Eq. (17) we find
$`u^2=0;C=1/(\sqrt{g}nu^3).`$ (21)
The remaining elements of $`F_{\mu \nu }`$ in this coordinate system are then:
$`F_{0a}=\beta F_{1a};F_{23}=(\beta u^0u^1)/u^3.`$
Then our vector $`B^\mu `$ and field tensor elements $`F_{\mu \nu }`$ are, for a given geometry and flow field, completely determined by one single flux function $`\beta `$.
## III. Vorticities in MHD plasmas
Flux conservation laws are generally expected if dissipative processes are negligible. Let us first discuss the non-relativistic case. The oldest example is the Kelvin/Helmholtz theorem for a neutral fluid with pressure $`p=p(\varrho )`$, where $`\varrho `$ is the mass density. Euler’s equation may then be written using the vorticity $`𝝎`$ and the Bernoulli function $`U`$ as follows:
$`{\displaystyle \frac{𝒗}{t}}+𝝎\times 𝒗+U=0,`$ (22)
and we obtain immediately Helmholtz’s equation for the vorticity by taking the curl:
$`{\displaystyle \frac{𝝎}{t}}+\times (𝝎\times 𝒗)=0.`$ (23)
This equation is the prototype of any vector field $`𝝎𝛀`$ which is “frozen in”, moving with the fluid velocity $`𝒗`$, and which is the curl of another field, say $`𝑽`$:
$`{\displaystyle \frac{𝛀}{t}}+\times (𝛀\times 𝒗)=0;𝛀\times 𝑽.`$ (24)
To get an equation for $`𝑽`$ we integrate Eq. (24), introducing a scalar potential $`\mathrm{\Phi }`$:
$`{\displaystyle \frac{𝑽}{t}}+𝛀\times 𝒗+\mathrm{\Phi }=0.`$ (25)
For $`𝑽=𝑨`$ and $`𝛀=𝑩`$ we obtain
$`𝑬+𝒗\times 𝑩=0,`$ (26)
where the electric field $`𝑬`$ is the usual expression in terms of the potentials $`𝑨,\mathrm{\Phi }`$ (with $`c=1`$). Eq. (26) is just the non-relativistic limit of Eq. (2). For other ideal fluid models one may also find a Kelvin/Helmholtz theorem, though the corresponding vectors $`𝑽,𝛀`$ are more complicated. The momentum balance of ideal MHD theory refers to a “center-of-mass” fluid with total pressure $`p`$, total mass density $`\varrho `$, center-of-mass velocity $`𝒗`$, and the Lorentz force due to the electric current density $`𝒋`$:
$`\varrho {\displaystyle \frac{d𝒗}{dt}}+p=𝒋\times 𝑩.`$
The picture of an idealized fluid requires not only zero resistivity according to Eq. (26), but, more generally, zero entropy production. So we expect for ideal MHD theory a second Kelvin/Helmholtz theorem and associated vectors $`𝑽`$, $`𝛀`$; they seem to be unknown, but they can be constructed using appropriate Lagrangian and Lin variables , . Here we use three Lin variables: the entropy per mass $`s`$ and the two Euler potentials $`q_\lambda (\lambda =1,2)`$ of the magnetic field:
$`𝑩=q_1\times q_2.`$
The constraints of entropy and magnetic flux conservation are then expressed by these three material invariants:
$`{\displaystyle \frac{ds}{dt}}={\displaystyle \frac{dq_1}{dt}}={\displaystyle \frac{dq_2}{dt}}=0,`$
and the Ansatz for $`𝑽`$ reads as follows:
$`𝑽=𝒗rs{\displaystyle \underset{\lambda =1}{\overset{2}{}}}v_\lambda q_\lambda .`$ (27)
The coefficients $`r,v_\lambda `$ are now determined in order to match Eq. (25) with an arbitrary potential $`\stackrel{~}{\mathrm{\Phi }}`$. From Eq. (27) and its curl, the equation of $`𝛀`$, we find after some rearrangements the following purely kinematical relation:
$`{\displaystyle \frac{𝑽}{t}}+𝛀\times 𝒗`$ $`=`$ $`{\displaystyle \frac{𝒗}{t}}+𝝎\times 𝒗{\displaystyle \frac{dr}{dt}}s+{\displaystyle \frac{ds}{dt}}r`$ (30)
$`{\displaystyle \underset{\lambda }{}}{\displaystyle \frac{dv_\lambda }{dt}}q_\lambda +{\displaystyle \underset{\lambda }{}}{\displaystyle \frac{dq_\lambda }{dt}}v_\lambda `$
$`\left(r{\displaystyle \frac{s}{t}}+{\displaystyle \underset{\lambda }{}}v_\lambda {\displaystyle \frac{q_\lambda }{t}}\right).`$
Here we insert our particular fluid model for $`𝒗/t`$. The ideal MHD equation can be written as follows:
$`{\displaystyle \frac{𝒗}{t}}+𝝎\times 𝒗=\left({\displaystyle \frac{𝒗^2}{2}}+h\right)+Ts+{\displaystyle \frac{𝒋}{\varrho }}\times 𝑩,`$
where $`h`$ is the enthalpy per mass, $`T`$ the temperature, and $`𝒋`$ the current density as determined from Ampère’s law. The right-hand side of Eq. (30) is then just $`\stackrel{~}{\mathrm{\Phi }}`$, with
$`\stackrel{~}{\mathrm{\Phi }}={\displaystyle \frac{𝒗^2}{2}}+h+r{\displaystyle \frac{s}{t}}+{\displaystyle \underset{\lambda }{}}v_\lambda {\displaystyle \frac{q_\lambda }{t}},`$
provided the coefficients $`r,v_\lambda `$ obey the following equations of motion:
$`{\displaystyle \frac{dr}{dt}}=T;{\displaystyle \frac{dv_1}{dt}}={\displaystyle \frac{1}{\varrho }}𝒋q_2;{\displaystyle \frac{dv_2}{dt}}={\displaystyle \frac{1}{\varrho }}𝒋q_1.`$
So we find a second flux conservation law in ideal MHD; it refers to the following generalized vorticity:
$`𝛀\times 𝑽=𝝎r\times s{\displaystyle \underset{\lambda }{}}v_\lambda \times q_\lambda .`$ (31)
The generalized Helmholtz equation (24) allows, of course, the “trivial” solution $`𝛀=0`$, corresponding to a potential flow for $`𝑽`$, $`𝑽=S`$; then Eq. (27) leads to the so-called Clebsch representation of $`𝒗`$:
$`𝒗=S+rs+{\displaystyle \underset{\lambda =1}{\overset{2}{}}}v_\lambda q_\lambda .`$
(For $`𝛀0`$ we would need a further pair of Clebsch variables $`v_\lambda ,q_\lambda `$ with $`dv_\lambda /dt=dq_\lambda /dt=0`$). In contrast to the first Kelvin/Helmhotz theorem associated with Eq. (26) we have no simple advantage from this second vorticity law; in particular, the equations of motion for $`r`$ and $`v_\lambda `$ have to be solved, in addition to the conservation laws for $`s`$ and $`q_\lambda `$ and (not shown here) the Bernoulli equation for $`S`$. The situation, however, becomes more transparent if we leave the MHD description and treat the plasma as a fluid of different species $`j(j=e:`$ electrons, $`j=i:`$ ions of any kind) interacting only via the electromagnetic field, with the following momentum balance for each species:
$`\varrho _j\left({\displaystyle \frac{}{t}}+𝒗_j\right)𝒗_j+p_j=e_jn_j(𝑬+𝒗_j\times 𝑩).`$
There is an obvious formal bridge to the MHD description for a two-component plasma: Neglecting the electron mass ($`m_e\varrho _e0`$) we obtain $`𝒗_i`$ as the center-of-mass velocity $`𝒗`$; assuming, in addition, quasi-neutrality with simply charged ions ($`n_e=n_i=n`$) we can replace $`𝒗_e`$ as follows:
$`𝒗_e=𝒗{\displaystyle \frac{𝒋}{|e|n}}.`$
Adding both momentum equations gives then the correct ideal MHD balance with $`p=p_e+p_i`$, while the momentum equation of electrons leads to a generalized Ohm’s law replacing Eq. (26):
$`𝑬+𝒗\times 𝑩={\displaystyle \frac{1}{|e|n}}(𝒋\times 𝑩p_e).`$
This “modified” MHD theory is the bridge mentioned above, but a simpler form of the two Kelvin/Helmholtz theorems is obtained if we come back to the multfluid description, now within general relativity.
## IV. Vorticities in isothermal plasmas
The flux conservation theorems associated with perfect fluid models depend crucially on the equation of state for the pressure, e.g., $`p=p(\varrho ,s)`$. The non-relativistic case in standard textbooks on neutral fluids usually assumes constant entropy $`s`$ throughout the fluid volume, leading to the usual Helmholtz equation, Eq. (23). The corresponding equation in general relativity has been derived by Taub . If the entropy of a non-relativistic neutral fluid varies in space, the situation is more subtle; it has been discussed already by Eckart and later in Ref. . The result is that the vector $`𝑽`$ whose vorticity flux is conserved differs from $`𝒗`$ by $`(rs)`$, the second term on the right-hand side of Eq. (27). For a non-relativistic multifluid plasma Eq. (27) is replaced by
$`𝑽_j=𝒗_jr_js_j+{\displaystyle \frac{e_j}{m_j}}𝑨,`$
leading to a “three-circulation theorem” corresponding to the three constituents of $`𝑽(𝑨`$ is the vector potential). In general relativity we have, instead of Eq. (25) (from which the general Helmholtz equation (24) follows), the equation corresponding to Eqs. (2) and (3) for each species (suppressing now the species index $`j`$):
$`u^\nu \mathrm{\Omega }_{\mu \nu }`$ $`=`$ $`0,`$
$`\mathrm{\Omega }_{\mu \nu }`$ $``$ $`V_{\nu ,\mu }V_{\mu ,\nu }.`$
The vector $`V_\nu `$ for a multifluid plasma with $`s=`$ const. has been given by Lichnérowicz and Carter (“single constituent perfect fluid”, his example (c) in §4), and for varying $`s`$ by Ref. , namely:
$`V_\mu ={\displaystyle \frac{\sigma }{\varrho }}u_\mu +rs_{,\mu }+{\displaystyle \frac{e}{m}}A_\mu ,`$
where $`\sigma `$ is the relativistic enthalpy per volume, and $`u_\mu `$ the covariant Eulerian four-velocity (with $`c=1)`$. One of Carter’s results refers also to a neutral fluid with two constituents and varying entropy, but the resulting canonical momentum per volume ($`=\varrho V_\mu `$) reads, in our notation, as follows (from Eqs. (4.25), (4.26), (4.32) and (4.33) of Ref. ):
$`\varrho V_\mu n\pi _\mu =\sigma u_\mu ,`$
so the term $`rs_{,\mu }`$ is missing.
Here we consider a different physical situation which may be of astrophysical interest: We assume that a radiation field acts like a heat reservoir for electrons and protons, maintaining a constant temperature of them. Then the term $`rs_{,\mu }`$ is again absent, and we find the relativistic Helmholtz equation for the ordinary canonical vorticity of both species. In this case, $`\sigma `$ will turn out to be the free relativistic enthalpy per volume. We start with the material energy-momentum-stress tensor $`T_\mu ^\nu `$ of an ideal electron or ion fluid without denoting the species index explicitely. Since both fluids are only coupled via the electromagnetic field by the Coulomb/Lorentz force, we can write the energy-momentum balance for both species as follows:
$`T_\mu {}_{}{}^{\nu }{}_{;\nu }{}^{}=enF_{\mu \nu }u^\nu .`$ (32)
The semi-colon indicates the covariant derivative. The coupling of both species by the right-hand side of Eq. (32) implies now that the magnetic flux is not conserved, neither for the electron nor for the ion fluid. The tensor $`T_\mu ^\nu `$ for an ideal fluid is well-known (see, e.g., Ref. ):
$`T_\mu ^\nu =\sigma u_\mu u^\nu p\delta _\mu ^\nu .`$
The last term in this equation is the scalar pressure $`p`$ in the local inertial rest frame times the Kronecker symbol, and the scalar $`\sigma `$ depends on the equation of state. Using also the mass conservation, Eq. (4), we obtain
$`T_\mu {}_{}{}^{\nu }{}_{;\nu }{}^{}=\varrho u^\nu \left({\displaystyle \frac{\sigma }{\varrho }}u_\mu \right)_{;\nu }p_{,\mu }.`$
This four-vector must be orthogonal to $`u^\mu `$, as is also the right-hand side of Eq. (32). With the normalization of $`u^\mu (u_\mu u^\mu =1`$) we find then
$`0`$ $`=`$ $`u^\mu T_\mu {}_{}{}^{\nu }{}_{;\nu }{}^{}=\varrho u^\nu \left({\displaystyle \frac{\sigma }{\varrho }}\right)_{,\nu }u^\nu p_{,\nu },`$
$`d\left({\displaystyle \frac{\sigma }{\varrho }}\right)={\displaystyle \frac{1}{\varrho }}dp,`$
where the differential $`d`$ means variation along the path of a fluid element. In the co-moving inertial rest frame we use the Gibbs-Duhem relation for the free enthalpy $`\mu `$ per mass:
$`d\mu =sdT+{\displaystyle \frac{1}{\varrho }}dp.`$ (33)
Ignoring temperature variations, and including the relativistic rest energy $`\varrho c^2`$ (with $`c1`$ for the moment) the resulting expression for $`\sigma `$ is then as follows:
$`\sigma =\varrho \left(1+{\displaystyle \frac{\mu }{c^2}}\right).`$ (34)
To derive a flux conservation law we insert the vector potential for $`F_{\mu \nu }`$ in Eq. (32) and put all terms to the left-hand side:
$`u^\nu \left[\left({\displaystyle \frac{\sigma }{\varrho }}u_\mu \right)_{;\nu }+{\displaystyle \frac{e}{m}}A_{\mu ;\nu }{\displaystyle \frac{e}{m}}A_{\nu ;\mu }\right]{\displaystyle \frac{1}{\varrho }}p_{,\mu }=0.`$
To eliminate the pressure term we calculate the partial derivative from Eqs. (33) and (34) with $`dT=0`$ (and this time $`c=1`$):
$`{\displaystyle \frac{1}{\varrho }}p_{,\mu }=\left({\displaystyle \frac{\sigma }{\varrho }}\right)_{,\mu }=u^\nu \left({\displaystyle \frac{\sigma }{\varrho }}u_\nu \right)_{;\mu },`$
where in the last step we used again the normalization of $`u^\nu `$. This is now just the term leading to flux conservation for the canonical vorticity of each species; we define
$`V_\mu `$ $``$ $`{\displaystyle \frac{\sigma }{\varrho }}u_\mu +{\displaystyle \frac{e}{m}}A_\mu ,`$ (35)
$`\mathrm{\Omega }_{\mu \nu }`$ $``$ $`V_{\nu ;\mu }V_{\mu ;\nu }=V_{\nu ,\mu }V_{\mu ,\nu },`$ (36)
and we find
$`\mathrm{\Omega }_{\mu \nu }u^\nu =0.`$ (37)
The solution of Eq. (37) for axisymmetric equilibria is now simply obtained from the previous section: We replace there $`F_{\mu \nu }`$ by $`\mathrm{\Omega }_{\mu \nu }`$ and $`A_\mu `$ by $`V_\mu `$. The mean velocities of electrons and ions, however, are usually different; the fluxes of their general vorticities $`\mathrm{\Omega }_{\mu \nu }`$ are therefore conserved in different frames. It is interesting to rewrite Eq. (37) for the spatial components of the canonical velocity, $`V^i`$, in the special-relativistic case (no gravity), namely:
$`{\displaystyle \frac{𝑽}{t}}+𝛀\times 𝒗+U=0,`$
where $`U(cV_0)`$ is the relativistic Bernoulli function, and $`\mathrm{\Omega }^i(\mathrm{\Omega }_{jk})`$ are the spatial vector components associated with $`\mathrm{\Omega }_{\mu \nu }`$. This equation is now again of the same form as Eq. (25), and we recover Eq. (24) as the prototype of any special-relativistic generalized Helmholtz equation.
## V. Multifluid plasma equations for axisymmetric equilibria
Let us use a general poloidal coordinate system to solve the mass conservation law, Eq. (4), for each species separately by introducing an appropriate stream function $`\chi `$, namely $`(\stackrel{~}{\epsilon }^{1ab}`$ is again the permutation symbol, here for spatial indices):
$`u^a={\displaystyle \frac{1}{\sqrt{g}n}}\stackrel{~}{\epsilon }^{1ab}\chi _{,b}.`$ (38)
The poloidal stream lines are then given by $`\chi =`$ const., and Helmholtz’ equation (37) in the symmetry plane, $`\mu =r=0,1`$, reads as follows:
$`0=V_{r,\nu }u^\nu =V_{r,a}u^a.`$
Similarly as in Eq. (19) we conclude that the components $`V_r`$ can be any flux functions with respect to $`\chi `$:
$`V_r=V_r(\chi ).`$
These two flux functions and the corresponding components of $`A_\nu `$ fix two components of $`u_\nu `$ for each species up to a factor $`\varrho /\sigma `$, namely:
$`u_r={\displaystyle \frac{\varrho }{\sigma }}\left(V_r(\chi ){\displaystyle \frac{e}{m}}A_r\right).`$ (39)
Assuming that $`\chi `$ and $`A_\nu `$ are given elsewhere we may read the normalization condition for $`u_\nu `$ as an equation for $`n`$:
$`1=g^{rs}u_ru_s+g_{ab}u^au^b.`$ (40)
So all fluid quantities besides $`\chi `$ are determined by Eqs. (38) - (40), and we find all elements of $`\mathrm{\Omega }_{\mu \nu }`$ except $`\mathrm{\Omega }_{ab}`$:
$`\mathrm{\Omega }_{rs}=0;\mathrm{\Omega }_{ra}=V_{r,a}=V_r^{}\chi _{,a},`$ (41)
where the prime means differentiation with respect to $`\chi `$. The general solution for $`\mathrm{\Omega }_{\mu \nu }`$, however, can be obtained from Eqs. (17) and (18) with appropriate changes of notation, namely:
$`\mathrm{\Omega }_{\mu \nu }`$ $`=`$ $`nC(\chi )\epsilon _{\mu \nu \varrho \sigma }u^\varrho K^\sigma ,`$ (42)
$`\beta `$ $`=`$ $`V_0^{}(\chi )/V_1^{}(\chi ).`$ (43)
Comparing this with the results above (Eqs. (38) and (41) ) we find the flux function $`C(\chi )`$:
$`C(\chi )=V_1^{}(\chi ).`$ (44)
The equation for the stream function $`\chi `$ is then obtained from Eq. (42) with $`\mu =2`$ and $`\nu =3`$, where the left-hand side is determined according to the definitions (36) and (35); the result is then the following:
$`\mathrm{\Omega }_{23}`$ $`=`$ $`\sqrt{g}nV_r^{}(\chi )u^r,`$ (45)
$`\mathrm{\Omega }_{23}`$ $``$ $`\left({\displaystyle \frac{\sigma }{\varrho }}u_3\right)_{,2}\left({\displaystyle \frac{\sigma }{\varrho }}u_2\right)_{,3}+{\displaystyle \frac{e}{m}}F_{23}.`$ (46)
The complete set of fluid equations (38) - (40) and (45) - (46) for the unknown variables $`u_r`$, $`n`$ and $`\chi `$ is written for any poloidal coordinate system, thus allowing any number of particle species. To solve finally Ampère’s equation in the poloidal plane we are then free to use particular coordinates and a particular gauge of $`A_\nu `$. To be consistent with the usual notation we denote the projections of the $`j^\mu `$-lines onto the poloidal plane the lines $`\mathrm{\Psi }=`$ const., where the flux function $`\mathrm{\Psi }`$ may be identified with a “radial” coordinate $`x^2`$ as in Sec. II, and the stream function $`I`$ of $`j^a`$ is denoted as a flux function, $`I=I(\mathrm{\Psi })`$. The continuity equation for $`j^\mu `$ in the poloidal plane is then solved as follows:
$`j^2=0;4\pi \sqrt{g}j^3=I^{}(\mathrm{\Psi }),`$ (47)
where the prime of $`I`$ means differentiation with respect to $`\mathrm{\Psi }`$. The flux function $`I(\mathrm{\Psi })`$ is, of course, not independent from the stream functions $`\chi `$ of the different species. From
$`j^\mu ={\displaystyle \underset{j}{}}enu^\mu `$
and Eq. (47) we find the following relation:
$`{\displaystyle \underset{j}{}}e\chi ={\displaystyle \frac{I(\mathrm{\Psi })}{4\pi }}+const.,`$ (48)
where the sum over $`j`$ is the sum with respect to the different particle species. A convenient gauge of $`A_\nu `$ is, as in the non-relativistic case, the condition that $`𝑨`$ is tangential to the surfaces $`\mathrm{\Psi }`$=const. of the current lines:
$`𝑨\mathrm{\Psi }=0.`$
This equation can usually be fulfilled by an appropriate gauge function since $`\mathrm{\Psi }`$ is a radial coordinate; for $`\mathrm{\Psi }=x^2`$ it reads
$`A_2=0\text{or}F_{23}=A_{3,2}.`$ (49)
Let us now start with Ampère’s equation for $`A_\nu `$:
$`\left[\sqrt{g}g^{\mu \varrho }g^{\nu \sigma }\left(A_{\sigma ,\varrho }A_{\varrho ,\sigma }\right)\right]_{,\nu }=4\pi \sqrt{g}j^\mu .`$ (50)
In the symmetry plane, $`\mu =r=0,1`$, this equation is decoupled from the poloidal components of $`A_\nu `$; inserting the fluid quantities for $`j^r`$ we have then the following set of two equations for the components $`A_s,s=0,1`$:
$`\left[\sqrt{g}g^{rs}g^{ab}A_{s,b}\right]_{,a}=4\pi \sqrt{g}g^{rs}{\displaystyle \underset{j}{}}en{\displaystyle \frac{\varrho }{\sigma }}\left(V_s(\chi ){\displaystyle \frac{e}{m}}A_s\right).`$ (51)
In the poloidal plane, Eq. (50) can be re-written with indices $`a,b,c,d`$ which run from 2 to 3 only:
$`\left[\sqrt{g}g^{ab}g^{cd}F_{bd}\right]_{,c}=4\pi \sqrt{g}j^a`$
The left-hand side is easily evaluated due to the antisymmetry of $`F_{bd}`$; we introduce the determinants of $`g_{ab}`$ and $`g_{rs}`$ explicitely:
$`g_{pol}det(g_{ab});g_{sym}=det(g_{rs}),`$
then we obtain the following form of Ampère’s equation in the poloidal plane for $`a=2`$:
$`\left[\sqrt{{\displaystyle \frac{g_{sym}}{g_{pol}}}}F_{23}\right]_{,3}=4\pi \sqrt{g}j^2,`$ (52)
and for $`a=3`$:
$`\left[\sqrt{{\displaystyle \frac{g_{sym}}{g_{pol}}}}F_{23}\right]_{,2}=4\pi \sqrt{g}j^3.`$ (53)
Inserting here Eq. (47) we identify the expression in brackets as a flux function, namely:
$`\sqrt{{\displaystyle \frac{g_{sym}}{g_{pol}}}}F_{23}=I(\mathrm{\Psi }).`$ (54)
In the non-relativistic case this is simply the covariant toroidal component of the magnetic field. Finally, the component $`A_3`$ is determined from Eqs. (54) and (49); it is not a flux function, and it is not needed in the remaining set of equations.
## VI. A variational principle
The numerical solution of our potential equations for $`\chi `$, Eqs. (45) - (46), and $`A_r`$, Eq. (51), is simplified by the existence of a functional of $`\chi `$ and $`A_r`$ which is stationary in equilibrium. We need, however, still another quantity whose variation leads to the normalization condition, Eq. (40). So we look for a functional $`W`$ of three quantities, $`W=W(\sigma ,\chi ,A_r)`$ say, whose independent variations lead to the equilibrium conditions. The particle density $`n`$ (which is not varied), and $`\sigma `$ of each species are then obtained afterwards from a combination of the normalization condition an an equation of state according to Eq. (34).
Let us start with the normalization condition, where $`u^a`$ and $`u_r`$ are given by Eqs. (38) and (39), respectively. The latter equation gives
$`\delta u_r={\displaystyle \frac{1}{\sigma }}u_r\delta \sigma +{\displaystyle \frac{\varrho }{\sigma }}V_r^{}(\chi )\delta \chi {\displaystyle \frac{en}{\sigma }}\delta A_r,`$ (55)
while $`\delta u^a`$ depends on $`\delta \chi `$ only:
$`\delta u^a={\displaystyle \frac{1}{\sqrt{g}n}}\stackrel{~}{\epsilon }^{1ab}(\delta \chi )_{,b}.`$ (56)
It is then easily realized that $`W`$ could be of the following form:
$`W(\sigma ,\chi ,A_r)={\displaystyle d^2x\sqrt{g}\underset{j}{}\frac{\sigma }{2}(g_{ab}u^au^bg^{rs}u_ru_s1)}+\mathrm{},`$
where the terms indicated by dots are independent of $`\sigma `$, and the integration is done in a fixed region of the poloidal plane $`(x^2,x^3)`$. It is a remarkable effect that in the expression above the usual invariant $`u_\nu u^\nu `$ is replaced by $`u_au^au_ru^r`$ which is only invariant under transformations in the poloidal plane. A similar replacement will be needed in the invariant which produces Maxwell’s equations in vacuum. Using our solution for $`F_{23}`$, Eq. (54), and $`F_{rs}=0,F_{ar}=A_{r,a}`$, we find
$`{\displaystyle \frac{1}{16\pi }}F_{\mu \nu }F^{\mu \nu }={\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{I^2(\mathrm{\Psi })}{(g_{sym})}}{\displaystyle \frac{1}{8\pi }}g^{rs}g^{ab}A_{r,a}A_{s,b}.`$
Changing now the sign of the last expression above, we are led to the following functional:
$`W(\sigma ,\chi ,A_r)={\displaystyle d^2x\sqrt{g}}`$ $`[`$ $`{\displaystyle \underset{j}{}}{\displaystyle \frac{\sigma }{2}}(g_{ab}u^au^bg^{rs}u_ru_s1)`$ (57)
$``$ $`{\displaystyle \frac{1}{8\pi }}{\displaystyle \frac{I^2(\mathrm{\Psi })}{(g_{sym})}}+{\displaystyle \frac{1}{8\pi }}g^{rs}g^{ab}A_{r,a}A_{s,b}]`$ (58)
Variations with respect to $`\chi `$ and $`A_r`$ are now done by eliminating derivatives of $`\delta \chi `$ and $`\delta A_r`$ by partial integrations, assuming that $`\delta \chi `$ and $`\delta A_r`$ vanish at the boundary. Furthermore, we have to vary $`I(\mathrm{\Psi })`$ according to Eq. (48), namely:
$`\delta I(\mathrm{\Psi })=(4\pi e)\delta \chi .`$ (59)
The total variation of $`W(\sigma ,\chi ,A_r)`$ is then obtained with the following result:
$`\delta W={\displaystyle d^2x\sqrt{g}}`$ $`\{{\displaystyle \underset{j}{}}{\displaystyle \frac{1}{2}}(g_{ab}u^au^b+g^{rs}u_ru_s1)\delta \sigma `$ (62)
$`+{\displaystyle \underset{j}{}}{\displaystyle \frac{1}{(g_{sym})}}\left(\stackrel{~}{\mathrm{\Delta }}\chi +eI(\mathrm{\Psi })+g_{sym}\varrho u^rV_r^{}(\chi )\right)\delta \chi `$
$`+[{\displaystyle \underset{j}{}}enu^r{\displaystyle \frac{1}{4\pi \sqrt{g}}}(\sqrt{g}g^{rs}g^{ab}A_{s,b})_{,a}]\delta A_r\},`$
$`\stackrel{~}{\mathrm{\Delta }}\chi `$ $``$ $`\sqrt{{\displaystyle \frac{g_{sym}}{g_{pol}}}}{\displaystyle \underset{a,b}{}}(\stackrel{~}{g}_{ab}\chi _{,b}\stackrel{~}{g}_{bb}\chi _{,a})_{,a},`$ (63)
$`\stackrel{~}{g}_{ab}`$ $``$ $`{\displaystyle \frac{\sigma }{\sqrt{g}n^2}}g_{ab}.`$ (64)
The vanishing factor of $`\delta A_r`$ is easily identified with Ampère’s equation in the symmetry plane, Eq. (51). To identify the vanishing factor of $`\delta \chi `$ we note that from Eqs. (46), (54), and (38) we have
$`m\mathrm{\Omega }_{23}={\displaystyle \underset{a,b}{}}(\stackrel{~}{g}_{ab}\chi _{,b}\stackrel{~}{g}_{bb}\chi _{,a})_{,a}+e\sqrt{{\displaystyle \frac{g_{pol}}{g_{sym}}}}I(\mathrm{\Psi }).`$
Inserting this in Eq. (45) with a slight rearrangement, we find
$`\stackrel{~}{\mathrm{\Delta }}\chi +eI(\mathrm{\Psi })+g_{sym}\varrho u^rV_r^{}(\chi )=0.`$ (65)
This is just the condition that $`W`$ is stationary with respect to variation of $`\chi `$, assuming that $`g_{sym}`$ is finite at this point.
## VII. Discussion
A plasma equilibrium near a rotating black hole has been considered in the ideal fluid picture. The usual MHD equations imply two flux conservation laws. One is the well-known conservation law of magnetic flux (Ohm’s law with vanishing resistivity); it leads to two constants of motion for stationary axisymmetric systems: The covariant time-like and toroidal components, $`A_t`$ and $`A_\phi `$, of the vector potential are constant on the poloidal stream lines of the plasma bulk velocity. This well-known fact is in contradiction with the coupling of $`A_t`$ and $`A_\phi `$ for a metric with $`g_{t\phi }0`$, as is appropriate for a rotating black hole. Simple Grad-Shafranov type MHD equilibria (see, e. g., Ref. are then ruled out in this case. The second flux conservation law of MHD has been derived in Section III for the non-relativistic case (Eqs. (24) and (27) - (31)); it is, however, not as simple as the magnetic flux conservation law. A more reasonable description for $`g_{t\phi }0`$ is given by the multifluid equations: They can be cast into the form of Helmholtz equations for each fluid component ; they are particularly simple for isothermal plasmas, as is shown here (Eqs. (35) - (37)). Since the form of these equations is exactly the same as the magnetic flux conservation law, we obtain two constants of motion for each species of an axisymmetric equilibrium, the time-like and toroidal components of the canonical velocity $`V_\nu `$. The component $`V_t`$ is the relativistic Bernoulli function, and $`V_\phi `$ is the canonical angular momentum per mass of a fluid particle of a particular species; this has to be expected for axisymmetric equilibria if the fluid components interact only through the electromagnetic field. The poloidal components of Ampère’s equation can be integrated, too, by adjusting the coordinates to the lines of the electric current, $`\mathrm{\Psi }=`$ const.; the relevant flux function $`I(\mathrm{\Psi })`$, Eq. (54), corresponding to the toroidal magnetic field, is now simply related to the stream functions $`\chi `$ of the different plasma components according to Eq. (48). The whole set of potential equations can now be summarized in a functional $`W(\sigma ,\chi ,A_r)`$, where $`\sigma `$ is the free enthalpy density of a particular species, $`\chi `$ the stream function of its poloidal velocity, and $`A_r`$ stands for $`A_t`$ and $`A_\phi `$. The total variation of $`W`$ produces then the nomalization condition for the Eulerian four-velocity of each species, the potential equation for $`\chi `$ (Eq. (65)), and Ampère’s equations for $`A_t`$ and $`A_\phi `$.
It is interesting to consider possible solutions of these equations in a given geometry with $`g_{t\phi }0`$ like Kerr’s metric. Equilibria with poloidal velocity fields are of interest from an observational point of view, because they are able to exchange mass, angular momentum etc. between inner and outer parts. Plasma models with pure poloidal velocity fields, however, are not possible; the reason is that the toroidal velocity $`u_\phi `$ is proportional to $`(V_\phi (e/m)A_\phi )`$ (Eq. (39)), where $`V_\phi `$ is constant on the stream lines of the particular species which is considered, but $`A_\phi `$ generally not. Equilibria with pure toroidal velocities are possible for an electron-positron plasma. In this case Eq. (65) is solved trivially with $`\chi 0`$, $`V_r`$ const., $`V_r^{}0`$, and the velocity components $`u_t`$, $`u_\phi `$ become equal but opposite in sign, up to a constant $`V_r`$ ; they are determined from Ampère’s equation (51) which is highly nonlinear due to the normalization condition. This solution, however, seems to be artificial, and an acceptable solution will exhibit inevitably also poloidal velocity components.
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# Predictions for Decays of Radially Excited Baryons
## I Introduction
In this note, we consider the strong decays of the lowest-lying radially excited baryons. These are the states comprising the first excited 56-plet, in SU(6) language. We make predictions for the decay widths of these states that are independent of any specific model for the binding potential or spatial wave functions.
A reason for taking up studies of excited baryons now is that we anticipate new results from the CLAS Collaboration at the Thomas Jefferson National Accelerator Facility. Not all the states in the radially excited 56-plet are yet discovered, and many of the measured decay widths of the observed states have large uncertainties. A set of predictions for the decay widths of unobserved states will be of great use to experimenters. Of necessity, we also provide predictions for the masses of the unobserved states. To this end, we present a derivation of the Gürsey-Radicati mass formula from large-$`N_c`$ QCD that we believe has not appeared explicitly in the literature.
Another and more specific reason for our analysis is to shed light on the Roper resonance, the $`N(1440)`$. The Roper is the first excitation of the nucleon with the same $`J^P=(1/2)^+`$ quantum numbers. The most direct explanation of the Roper is as a three-quark radially excited state. However, certain of its apparent peculiarities have led to plausible suggestions that it may be a hybrid state (a state whose lowest significant Fock component is three quarks plus gluonic excitations) or that it may be a cross section enhancement that does not correspond to any resonance . An argument against interpreting the Roper as a three-quark state has been that its calculated mass is too high in quark models with one-gluon-exchange, i.e. spin-color, interactions . Recently, it has been found that the observed mass is perfectly consistent with the three-quark picture if mass splittings are due to spin-flavor interactions instead . Moreover, there is evidence of a Roper signal from a lattice calculation using a three-quark source operator, although uncertainties in the mass determination as yet preclude saying if it is lighter than its negative parity counterpart . In the present work, we assume the Roper resonance is a three-quark state embedded in the excited 56-plet; confirmation of the decay widths predicted here will therefore support the three-quark interpretation of this state.
The value of a model-independent analysis should be clear: We are free of any assumptions regarding the binding potential or quark spatial wave functions. We do make assumptions about the dominance of single-quark operators and about configuration mixing that we now consider in turn.
First, we assume that the decays of interest proceed via a single quark interaction vertex. We may write the decay operator in the same SU(6) tensor product space in which we define the baryon spin-flavor wave functions:
$$_{\mathrm{eff}}G_{}^{ia}k^i\pi ^a.$$
(1)
Here $`G_{}^{ia}=S_{}^iT_{}^a`$ is an SU(6) generator, where $`S_{}^i`$ and $`T_{}^a`$ are the generators of SU(2) and SU(3), respectively, $`k^i`$ is the meson three-momentum, and $`\pi ^a`$ represents a meson field operator. The asterisk in Eq. (1) indicates that that the operator acts on the spin-flavor wave function of the excited quark. Here we adopt a convenient picture of the baryon state that follows from large-$`N_c`$ QCD, namely, that the baryon consists of a ‘core’ of $`N_c1`$ ground state quarks generating a collective potential in which a single quark is excited. Note that the complete symmetry of the 56-plet spatial wave functions implies that matrix elements of the operator in Eq. (1) are proportional to matrix elements of the same spin-flavor operator summed over all the quarks in the baryon state,
$$G^{ia}\underset{\mathrm{quarks}\alpha }{}S_\alpha ^iT_\alpha ^a.$$
(2)
Thus, the operator in Eq. (1) works equally well in parameterizing the decays of radially excited baryons for $`N_c=3`$, where the state may be described more realistically as a collective excitation of all three quarks. Matrix elements of this operator may be written
$$\mathrm{\Psi }(B_f,\pi ^a)|_{\mathrm{eff}}|\mathrm{\Psi }(B_i)=f(k)k^jB_f|G^{ja}|B_i$$
(3)
where $`k=|\stackrel{}{k}|`$ and $`f(k)`$ is a function that parameterizes the momentum dependence of the amplitude. This includes momentum dependence that originates from the underlying quark-level vertex, as well as from the overlap of the baryon spatial wave functions. The baryon states shown represent only the spin-flavor part of the wave function, and are nonrelativistically normalized. Use of the single-body decay approximation is strongly motivated by its success in describing decays of the orbitally excited, 70-plet baryons, both in the algebraic SU(6) approach , as well as more recent large-$`N_c`$ effective field theory analyses . In the latter case, both one- and two-body operators have been shown to contribute to 70-plet decay matrix elements at leading order in the $`1/N_c`$ expansion. Nonetheless, comparison to data for both strong decays and photoproduction amplitudes reveals that the two-body operators are phenomenologically irrelevant. It is not unreasonable to assume that similar dynamics may play a role in the strong decays of other excited SU(6) multipletsNote that the suppression of two-body operators in the present case implies that F-wave decay amplitudes should be highly suppressed..
Our second assumption is that mixing between states in the excited 56-plet and other SU(6) multiplets, i.e. configuration mixing, can be neglected. In the case of the 70-plet, successful large-$`N_c`$ effective theory descriptions of the strong decays , photoproduction amplitudes , and the nonstrange mass spectrum all neglect mixing with other SU(6) multiplets. In the present case, one might worry that the nearest observed states with appropriate quantum numbers for mixing, members of the positive-parity 70-plet, are only between 200 and 300 MeV heavier than the states of interest, so that a significant effect cannot be precluded. Fortunately, evidence from explicit quark models suggests that configuration mixing involving the excited 56 is also small , and even smaller than one would expect from $`1/N_c`$ arguments . We simply adopt this as a working assumption, without wedding ourselves to a particular quark model or choice of baryon spatial wave function.
## II Analysis of Observed States
We first determine the functional form and normalization of $`f(k)`$ by considering the decay modes of observed members of the excited 56-plet, the $`N(1440)`$, $`\mathrm{\Delta }(1600)`$, $`\mathrm{\Lambda }(1600)`$ and the $`\mathrm{\Sigma }(1660)`$. We will refer to this multiplet as the 56 henceforth. (The RPP refers to the 56 as the $`(56,0_2^+)`$.) Our input values are shown in Table I, and have been extracted from data in the Review of Particle Physics (RPP) . The one standard deviation errors on the baryon masses were taken to be half of the corresponding mass ranges given in the RPP, and partial widths were computed using the quoted full widths, branching fractions, and associated errors. The values for $`f(k)`$ corresponding to each observed decay mode are shown in Fig. 1. Since we wish to model the function $`f(k)`$ for momenta from $`100`$ to $`500`$ MeV, we require only that we find a simple functional form that works well within this range. The most successful result is $`f(k)=(2.8\pm 0.2)/k`$, which has a $`\chi ^2`$ per degree of freedom of $`1.1`$. Note that the functional forms $`a+kb`$, $`a/k^{0.5}`$, $`a/k^{1.5}`$ lead to a $`\chi ^2`$ per degree of freedom of $`1.9`$, $`1.6`$, and $`2.8`$, respectively. It is somewhat remarkable that such a simple functional form can account for the data in the momentum range of interest, though admittedly, the experimental uncertainties are large. It is not inconceivable, for example, that improvement in the data could lead to a preference for the less singular-looking linear fit, but given the present errors, the difference between a linear and the $`1/k`$ fit makes little difference in our decay predictions. Continuing in our spirit of model independence, we do not try to ascertain the origin of $`f(k)`$, but explore how far we can go in predicting decay modes for unobserved states.
## III Mass Predictions
We use the Gürsey-Radicati (GR) formula to predict masses for the unobserved members of the 56-plet:
$$M=A+BN_s+C[I(I+1)N_s^2/4]+DJ(J+1).$$
(4)
Here J (I) is the baryon spin (isospin) and $`N_s`$ is the number of valence strange quarks. The GR formula predicts spin- and flavor-dependent mass splitting for completely symmetric SU(6) baryon multiplets in terms of the four parameters $`A`$, $`B`$, $`C`$ and $`D`$; in our case, these parameters are determined from the masses given in Table I. While the GR formula was originally a conjecture , and later ‘derived’ by assuming SU(6) symmetry and a specific set of mass operators transforming in symmetry-breaking representations of small dimensionality , we begin by showing how Eq. (4) may be obtained more rigorously from large-$`N_c`$ QCD. This observation follows from the discussion in Ref. , but is not presented there explicitly. In a large-$`N_c`$ operator analysis for the 56-plet, we may write the general mass formula
$`M`$ $`=`$ $`a_1𝟙+𝕒_\mathrm{𝟚}{\displaystyle \frac{𝕊^\mathrm{𝟚}}{_𝕔}}+ϵ𝕒_\mathrm{𝟛}𝕋^\mathrm{𝟠}+ϵ𝕒_\mathrm{𝟜}{\displaystyle \frac{𝕊^𝕚𝔾^{𝕚\mathrm{𝟠}}}{_𝕔}}++ϵ𝕒_\mathrm{𝟝}{\displaystyle \frac{𝕊^\mathrm{𝟚}𝕋^\mathrm{𝟠}}{_𝕔^\mathrm{𝟚}}}+ϵ^\mathrm{𝟚}𝕒_\mathrm{𝟞}{\displaystyle \frac{𝕋^\mathrm{𝟠}𝕋^\mathrm{𝟠}}{_𝕔}}`$ (5)
$`+`$ $`ϵ^2a_7{\displaystyle \frac{T^8S^iG^{i8}}{N_c^2}}+ϵ^3a_8{\displaystyle \frac{T^8T^8T^8}{N_c^2}},`$ (6)
where $`S`$, $`T`$, and $`G`$ are the spin, flavor, and spin-flavor generators of SU(6) (with a sum over quarks left implicit), and $`ϵ1/31/N_c`$ parameterizes the size of SU(3) breaking. Notice that the eight order-one coefficients, $`a_1\mathrm{}a_8`$, completely span the space of observables, namely the eight baryon mass eigenvalues. Assuming that the baryon states have spin, isospin, and strangeness of order one in the large $`N_c`$ limit, we may discard all but the first four terms if we choose to work only up to subleading order. Acting on a large-$`N_c`$ baryon state
$$S^2=S(S+1),$$
(7)
$$T^8=(N_c3N_s)/\sqrt{12},$$
(8)
and for totally-symmetric spin-flavor wave functions
$$S^iG^{i8}=\frac{1}{4\sqrt{3}}\left[3I(I+1)S(S+1)3N_s(N_s+2)/4\right].$$
(9)
These identities imply that the only effects of the discarded terms up through order $`1/N_c`$ are redefinitions of the first four coefficients, $`a_1\mathrm{}a_4a_1^{}\mathrm{}a_4^{}`$. Thus we may write the mass eigenvalues
$`M=a_1^{}N_c`$ $`+`$ $`{\displaystyle \frac{a_2^{}}{N_c}}S(S+1)+ϵa_3^{}(N_c3N_s)`$ (10)
$`+`$ $`ϵ{\displaystyle \frac{a_4^{}}{2\sqrt{12}N_c}}\left[3I(I+1)S(S+1)3N_s(N_s+2)/4\right]+𝒪(1/N_c^2).`$ (11)
Eq. (4) then follows from a trivial redefinition of the coefficients. In modern language, the SU(6) breaking representations that lead to the GR mass formula are precisely those that arise at leading and subleading order in the $`1/N_c`$ expansion.
Using the observed masses in Table I, we may solve for the coefficients in Eq. (4). In units of MeV, we find
| $`A=1406.3\pm 31.3`$, | $`B=195.0\pm 43.0`$, |
| --- | --- |
| $`C=15.0\pm 38.1`$, | $`D=43.3\pm 46.0`$ , |
from which we predict the following states:
| $`\mathrm{\Xi }^{}(1825\pm 98),`$ | $`\mathrm{\Sigma }_{}^{}{}_{}{}^{}(1790\pm 192),`$ | $`\mathrm{\Xi }_{}^{}{}_{}{}^{}(1955\pm 196),`$ | $`\mathrm{\Omega }^{}(2120\pm 234)`$ . |
| --- | --- | --- | --- |
## IV Decay Width Predictions
From Eq. (3), the decay width is given by
$$\mathrm{\Gamma }=\frac{M_f}{6\pi M_i}k^2f(k)^2|𝒢|^2$$
(12)
where $`M_i`$ and $`M_f`$ are the masses of the initial and final baryons and
$$|𝒢|^2|B_f|G_{ja}|B_i|^2,$$
(13)
where the sum is over final state spins and isospins. The quantity $`|𝒢|^2`$ may be obtained either by use of symbolic math code or by using
$$B_f|G_{ja}|B_i=B_f(I_f,\alpha _f,S_f,m_f)|G_{ja}|B_i(I_i,\alpha _i,S_i,m_i)=𝒢\left(\begin{array}{ccc}I_f& I_a& I_i\\ \alpha _f& \alpha _a& \alpha _i\end{array}\right)\left(\begin{array}{ccc}S_f& 1& S_i\\ m_f& j& m_i\end{array}\right)$$
(14)
and evaluating the left-hand side for one particular final state. Here, $`\alpha `$ and $`m`$ stand for isospin and spin projections, and the sum in equation (13) is over $`\alpha _f`$, $`\alpha _a`$, $`m_f`$, and $`j`$.
We have already discussed fitting the function $`f(k)`$ to the measured decays, and how to use the Gürsey-Radicati mass formula to obtain the masses of the unobserved 56 states. Using the results of the fit with $`f(k)1/k`$, we predict the remainder of the strong decays of the 56. The results are shown in Table II, along with the masses we use for the initial baryons.
Table II does not quote the uncertainties for the widths. Numerical uncertainties come from two sources. One is the uncertainty in the function $`f(k)`$, and the other is the uncertainty in the masses. The uncertainty in the fitted $`f(k)`$ leads to about $`\pm 15`$% uncertainties in the widths. Further width uncertainty induced by mass uncertainty can be large, as discussed below. Of course, once the state is found, the widths in Table II can be recalculated easily and accurately with the correct mass.
## V Discussion
The uncertainty in predicting the width of the 56 decays is dominated not by uncertainty in the matrix elements but by uncertainty in the masses of the states. Four of the 56 masses are measured and four are not. We must predict the unmeasured masses, and do so using the Gürsey-Radicati mass formula. The accuracy of masses predicted from the Gürsey-Radicati formula is ultimately limited by the the approximations inherent in its derivation, and is easily estimated in a large-$`N_c`$ scheme. Currently, however, uncertainties in the predicted masses are dominated by uncertainties in the measured masses, and they are not small.
Table II shows the predicted partial width of the various decay modes, using the central values of the predicted masses. To help see how the mass uncertainties affect us, consider the equal spacing rule that follows from the GR formula. The decuplet spacing is given by
$$\mathrm{\Delta }m_{10^{}}=\frac{3}{2}m(\mathrm{\Lambda }^{})\frac{1}{2}m(\mathrm{\Sigma }^{})m(N^{})=(165\pm 110)\mathrm{MeV}.$$
(15)
(For comparison, the equivalent prediction for the ground state decuplet is $`\mathrm{\Delta }m_{10}=139`$ MeV, to be compared to 153, 148, and 139 MeV for the $`\mathrm{\Delta }`$-$`\mathrm{\Sigma }^{}`$,$`\mathrm{\Sigma }^{}`$-$`\mathrm{\Xi }^{}`$, and $`\mathrm{\Xi }^{}`$-$`\mathrm{\Omega }`$ mass splittings, respectively.) Then, as one example, the predicted central value for the $`\mathrm{\Omega }^{}`$ mass is 2.12 GeV, and this leads to a sizeable decay width for $`\mathrm{\Omega }^{}\mathrm{\Xi }^{}\overline{K}`$. However, the uncertainties allow the $`\mathrm{\Omega }^{}`$ to lie below the threshold for this decay.
Moreover, a simple weighted averaging of the published results that the RPP uses as the basis of their mass estimates for the $`\mathrm{\Lambda }^{}`$ and $`\mathrm{\Sigma }^{}`$ gives somewhat different central masses and much tighter uncertainties than they conservatively quote. One finds
$`m(\mathrm{\Lambda }^{})=(1600\pm 7)\mathrm{MeV}`$ (16)
$`m(\mathrm{\Sigma }^{})=(1672\pm 5)\mathrm{MeV}.`$ (17)
(The error limit on the $`\mathrm{\Lambda }^{}`$ mass includes a scale factor $`S=1.4`$, in accordance with procedures described in the RPP narrative.) Using these masses together with the previous $`N^{}`$ and $`\mathrm{\Delta }^{}`$ masses yields
$$\mathrm{\Delta }m_{10^{}}=(114\pm 20)\mathrm{MeV}$$
(18)
This reduced spacing leads to a number of decays listed in Table II being kinematically forbidden.
Unfortunately, there is not much one can do theoretically to reduce the uncertainties originating from the 56 masses. Table II simply presents the most reliable predictions we can make with the current data. As unobserved states are discovered, and their masses measured, one can easily revise the widths in Table II. To reiterate, the relevant formula for our $`kf(k)=const.=a`$ fit is
$$\mathrm{\Gamma }=\frac{kM_f}{6\pi M_i}a^2|𝒢|^2,$$
(19)
where $`M_i`$, $`M_f`$, and $`k`$ are defined already, the matrix element sum $`|𝒢|^2`$ is given in Table II, and $`a`$ is 2.8 $`\pm `$ 0.20, based on the current data.
The function $`f(k)`$ is not taken from any model, but fit to the limited data with the fulfilled hope that there would be a good fit using a simple form. The form one would expect from an atomic physics calculation is rather different from the one we found. For the hydrogen atom, in a nonrelativistic calculation, $`f(k)`$ would be proportional to a matrix element of $`\mathrm{exp}(i\stackrel{}{k}\stackrel{}{r})`$ between 2S and 1S wave functions. Further in a hydrogen atom, the wavelength of the outgoing radiation is long compared to the Bohr radius, and the small $`k`$ result that $`f(k)k^2`$ is valid. This leads to a very small width and the famous metastability of the hydrogen 2S state. For particle decay, the wavelength of the outgoing meson is comparable to and often smaller than the size scale of the baryon states. One should have no expectation that a low $`k`$ approximation will work. This also applies to barrier factors in the non-metastable case.
A model or an eventual ab initio calculation for the baryon states will produce some definite $`f(k)`$ that may or may not look analytically like our form, but to fit the data must not be numerically dissimilar to our result.
To summarize, we have without reference to models, but using the limited available data, studied the strong decays of the lowest-lying radially excited baryons. We assumed, justified by previous studies, that one-body operators dominated and that configuration mixing could be neglected. Simple forms for the one momentum-dependent function give a good fit to existing data, and allow prediction of 22 additional decay modes for this multiplet.
Acknowledgments
We thank R. Horgan for useful communications, and R. F. Lebed for a careful reading of the manuscript. CDC and CEC thank the National Science Foundation for support under Grant No. PHY-9900657. In addition, CDC thanks the National Science Foundation for support under grant No. PHY-9800741 and the Jeffress Memorial Trust for support under Grant No. J-532.
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# Generalized Faddeev Equations for 𝑁-Particle Scattering
## I INTRODUCTION
There has been an increased interest in many-body scattering dynamics due to the rapid growth of computer power. The fundamental method of formal scattering theory is the Lippmann-Schwinger equation. This equation is equivalent to the Schrödinger equation for a two-body problem with an appropriate scattering boundary condition. However, it is well recognized that there are problems when solving the Lippmann-Schwinger equation for more than two particles. In particular for three particles, but the same problem occurs when there are more particles, a perturbation expansion in powers of the potential, or better in terms of pair particle T-matrices, allows a third particle to be a spectator and a dangerous delta function to arise in the momentum representation of the Lippmann-Schwinger equation . Faddeev rearranged the equation so as to eliminate the dangerous delta function and in the process he was able to show that the kernel of the rewritten integral equation was compact, essentially by showing that it became connected after one iteration. Lovelace and later, Alt, Grassberger and Sandhas (AGS) gave two different ways of rewriting the Faddeev equations in terms of scattering amplitudes for various collision processes. The purpose of this work is to propose a way of generalizing the Faddeev equation to an arbitrary number of particles.
There has been a large effort in extending Faddeev’s work to an $`N`$-particle system, in particular to the four-particle system, both for theoretical interest and for practical use in nuclear and chemical reactions . Most work puts the emphasis on deriving equations with a connectable kernel, namely, a kernel that is connected after a finite number of iterations. Glöckle showed that a set of generalized Lippmann-Schwinger equations has a unique solution for the three-particle system. He also formulated similar equations for a four-particle scattering system. Special formulations were given by Weinberg and Rosenberg, the former emphasized the quasiparticle picture of an interacting $`N`$-particle system, the latter provided a set of equations that do not explicitly depend on the scattering potentials, a feature that may be useful for relativistic scattering. A channel array approach was taken by Baer et al., to achieve a connected kernel for a three-body system which was later generalized to an $`N`$-particle scattering system by Tobocman. Yakubovsky presented a formal generalization of the Faddeev equations to $`N`$-particle scattering by appropriately organizing pair-particle interactions. Yakubovsky’s equations , like those of Faddeev, have considerable advantage in dealing with the bound state problem since they are equivalent to the Schrödinger equation, thus there being less chance of spurious solutions. Unfortunately, the complexity and the lack of a simple method of iterating the Yakubovsky equations hinder their practical usage in physical and chemical problems. In a later treatment, Narodetsky and Yakubovsky proposed an alternative two-cluster approach and obtained a new set of equations for the $`N`$-particle scattering. The equivalence of these two-cluster equations to the Schrödinger equation has not be proved. Sloan made an important contribution to solving the four-particle problem by allowing only two-body channels in the coupled equations for scattering amplitudes. Bencze obtained generalized Lovelace equations and generalized AGS equations by resuming the two-cluster equations obtained by Narodetsky and Yakubovsky and showed that these equations are equivalent to those of Sloan when $`N=4`$. Redish generalized Sloan’s equations to an arbitrary number of particles and demonstrated that his results are identical to those of Bencze. The Sloan-Bencze-Redish equations are connected after one iteration provided that the kernels of the subpartitions are connected.
These previously presented methods appear to emphasize the pair particle transition operator and expand everything in terms of it. This approach of building up everything from pair particle properties involves combining the pair transition operators in more and more elaborate ways, the details of which may often be very tedious. In fact, there are so many terms that they are very cumbersome to implement. Even the comparison of the numerical calculations using different methods appears to raise questions. What is proposed here might be thought of as an approach of pulling the problem apart into smaller pieces. Thus the $`N`$-particle system (if $`N`$ is even) is divided into all possible combinations of $`N/2`$-particle subsystems with the coupling of the subsystems carried out in a manner parallel to the method of Faddeev. This is repeated for each $`N/2`$-particle subsystem. Eventually, after successive divisions, each subsystem consists of two particles, whose collisional effects are described by a pair transition operator. (A minor modification is required if $`N`$ is odd.) This gives a structure as to how the $`N`$-particle transition operator is dependent on the transition operator of fewer and fewer particles. The calculation of the $`N`$-particle transition operator then involves the reverse procedure of successively putting together the transition operators of fewer particles to eventually get an expression for the full transition operator. In going up the chain of partitionings, each step leads to an integral equation of Faddeev type which is connected after one iteration. In this way a set of equations is obtained in which all particles are connected, but the connectedness is organized in sets within sets. A necessary technicality for carrying out this approach is the need to renormalize the potential at each partitioning. Thus the transition operators for smaller sets of particles involve a potential which is generally smaller in magnitude than the true potential. It is believed that the approach presented here has a different basis of approach than those previously presented, has fewer terms and hopefully leads to a more efficient numerical procedure.
This paper is divided into five sections. Section II is devoted to the partitioning of $`N`$-particle systems, with emphasis on partitions having two clusters with an equal number of particles. The notation is similar to that of Redish and Yakubovsky . The distinction between a channel and a partition is stressed. A scheme is developed in Section III to generalize the Faddeev equations to a set of equations for $`N`$-particles. The derivation given for the generalized Faddeev equations closely parallels Faddeev’s original derivation. This emphasizes the simplicity of the new equations. The generalized Faddeev equations for the resolvent operator and the decomposition of the scattering wave functions are given in Section IV. Section V discusses how the proposed procedure would be applied to both 4- and 8-particle scattering. A general discussion in Section VI ends the paper.
## II Properties of Equal Partitioning
Consider a system of $`N`$ labeled particles $`(1,2,3,\mathrm{},N)`$. There are many topologically distinct ways of partitioning the labeled $`N`$ particles into sets of clusters of particles. Each particular way of dividing up the $`N`$ particles will be called a partition, labeled by $`A,B,C`$, etc. The detailed presentation given here partitions the $`N`$ particles into two clusters of nearly equal size (if $`N`$ is odd, one cluster is to have one more particle than the other), but alternate partitioning schemes are also possible, as mentioned in the Discussion. All partitionings of the same size are to be considered, which differ by particle composition in each cluster, to give a set of partitions $`𝒞`$. The objective of this section is to define the potentials and hamiltonians appropriate for such partitions. As used in this work, a partition only becomes a (rearrangement) channel when the clusters composing the partition are bound, thus distinguishing between the mathematical method of solving the $`N`$-particle problem and the physical notion of what asymptotic states arise in a scattering process.
The $`N`$-particle hamiltonian $`𝐇=𝐇_0+𝐕`$ consists of a kinetic energy operator $`𝐇_0`$ and an assumed pairwise additive potential operator
$$𝐕\underset{i<j}{}𝐕_{ij},$$
(1)
where the sum is over all possible ordered pairs, with $`i,j\{1,2,3,\mathrm{},N\}`$. The total hamiltonian is assumed to be self-adjoint and bounded below, and to have spectra that might include a discrete set, associated with the bound eigenstates, as well as the continuum for the scattering states. The objective is to determine the solutions for the Schrödinger equation associated with the total hamiltonian $`𝐇`$,
$$𝐇\mathrm{\Psi }=E\mathrm{\Psi }.$$
(2)
In the present treatment, the $`N`$ particles are partitioned into two clusters of equal size if at all possible. That is, if $`N`$ is even, there are $`N/2`$ particles in each, whereas if $`N`$ is odd, one cluster in a partition has $`(N+1)/2`$ particles and the other $`(N1)/2`$. It follows that the number of partitions in this set $`𝒞`$ of partitions is given by
$$𝒩_N=\{\begin{array}{cc}\frac{1}{2}\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)\hfill & \text{for }N\text{ even;}\hfill \\ \left(\genfrac{}{}{0pt}{}{N}{\left(N1\right)/2}\right)\hfill & \text{for }N\text{ odd.}\hfill \end{array}$$
(3)
The partition potentials $`𝐕_C`$ are to be chosen so that
$$\underset{C𝒞}{}𝐕_C𝐕.$$
(4)
Thus the potential $`𝐕_C`$ of partition $`C`$ is defined as
$$𝐕_C=v_N\underset{\genfrac{}{}{0pt}{}{j<k}{j,kC}}{}𝐕_{jk},$$
(5)
involving the sum over those pairs of particles $`jk`$ that appear in the same cluster in the partition $`C`$. Here $`v_N`$ is a renormalization factor and obviously, $`v_3=1`$. For $`N>3`$, $`v_N`$ is chosen as
$$v_N=\{\begin{array}{cc}\left(\genfrac{}{}{0pt}{}{N}{2}\right)/\left[\left(\genfrac{}{}{0pt}{}{N/2}{2}\right)\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)\right]=\frac{(N/2)!((N/2)2)!}{(N2)!}\hfill & N\text{ even;}\hfill \\ \left(\genfrac{}{}{0pt}{}{N}{2}\right)/\left\{\left(\genfrac{}{}{0pt}{}{N}{(N1)/2}\right)\left[\left(\genfrac{}{}{0pt}{}{(N1)/2}{2}\right)+\left(\genfrac{}{}{0pt}{}{(N+1)/2}{2}\right)\right]\right\}=\hfill & \\ \frac{((N+1)/2)!((N3)/2)!}{(N1)!}\hfill & N\text{ odd,}\hfill \end{array}$$
(6)
so that the total number of potential terms on both sides of Eq. (4) are equal. The special case of $`v_4=1`$ implies that no renormalization is needed for $`N=4`$, whereas $`v_5=1/4`$ and $`v_6=1/4`$ demonstrate how the potential for a particular pair of particles gets distributed among the many partitions in the set $`𝒞`$ when $`N`$ is large. Note that other types of two cluster partitions can also be selected (or included), which however, leads to a larger number of coupled equations for a given $`N`$-particle system.
A partition hamiltonian is defined as the kinetic energy operator $`𝐇_0`$ plus the mutual potential interactions between those particles in each of its clusters, namely the partition potential, thus
$$𝐇_C𝐇_0+𝐕_C.$$
(7)
The residual interaction, $`𝐕^C`$, describing how the clusters in partition $`C`$ interact with each other, is defined as
$$𝐕^C𝐕𝐕_C.$$
(8)
Thus, for each partition $`C`$ the total hamiltonian $`𝐇`$ can be written as
$$𝐇=𝐇_C+𝐕^C.$$
(9)
The resolvent operator $`𝐑(z)`$ for the total hamiltonian is defined as
$$𝐑=\frac{1}{z𝐇},$$
(10)
where $`z`$ is a complex parameter, for scattering theory having a small positive imaginary part. In most equations in this work the dependence on $`z`$ will be implicitly assumed. The free resolvent operator is defined as
$$𝐑_0=\frac{1}{z𝐇_0}$$
(11)
while the resolvent operator for partition $`C`$ is
$$𝐑_C\frac{1}{z𝐇_C}.$$
(12)
The object of this paper is to express the total transition operator $`𝐓`$, as determined by the Lippmann-Schwinger equation
$$𝐓𝐕+\mathrm{𝐕𝐑}_0𝐓,$$
(13)
and the Schrödinger equation (2) in terms of the clustering discussed above.
## III GENERALIZED FADDEEV EQUATIONS FOR $`𝐓`$
In this section a two-step scheme is proposed for the decomposition of the total transition operator $`𝐓`$ in such a manner that its kernel is connected after one iteration. The first step is to decompose $`𝐓`$ according to the selected set of two-cluster partitions $`𝒞`$, see Sec. II. Conceptually, this is exactly what Faddeev did for the three-body system. However for a system of more than three particles, the resulting expression for a partition transition operator $`𝐓_C`$ is governed by a combination of the transition operators for the two clusters. The second step is to identify how this combination is to be carried out, compare Ref. . Finally it is recognized that the separate cluster transition operators are exactly similar to the original $`N`$-particle transition operator but with significantly fewer particles, so the whole procedure can be repeated for each cluster. Care must be exercised when reconstructing the partition transition operators from the separate cluster transition operators so that the energies appearing in the resolvents are consistent with the total energy of the system.
### A The $`N`$-Particle Transition Operator
The full $`N`$-particle transition operator has the equivalent forms
$`𝐓`$ $`=`$ $`𝐕+\mathrm{𝐕𝐑𝐕}`$ (14)
$`=`$ $`𝐕+\mathrm{𝐓𝐑}_0𝐕`$ (15)
$`=`$ $`𝐕+\mathrm{𝐕𝐑}_0𝐓.`$ (16)
Moreover, it can be decomposed according to
$`𝐓`$ $`=`$ $`{\displaystyle \underset{C𝒞}{}}𝐕_C+{\displaystyle \underset{C𝒞}{}}𝐕_C𝐑_0𝐓`$ (17)
$`=`$ $`{\displaystyle \underset{C𝒞}{}}𝐓^C,`$ (18)
where
$$𝐓^C𝐕_C+𝐕_C𝐑_0𝐓.$$
(19)
Subtracting $`𝐕_C𝐑_0𝐓^C`$ from both sides of Eq. (19) gives
$$\left(\mathrm{𝟏}𝐕_C𝐑_0\right)𝐓^C=𝐕_C+𝐕_C𝐑_0\left(𝐓𝐓^C\right),$$
(20)
which, on multiplying by $`(\mathrm{𝟏}𝐕_C𝐑_0)^1`$ gives
$`𝐓^C`$ $`=`$ $`\left(\mathrm{𝟏}𝐕_C𝐑_0\right)^1𝐕_C`$ (22)
$`+\left(\mathrm{𝟏}𝐕_C𝐑_0\right)^1𝐕_C𝐑_0\left(𝐓𝐓^C\right)`$
$`=`$ $`𝐓_C+{\displaystyle \underset{B𝒞}{}}\overline{\delta }_{C,B}𝐓_C𝐑_0𝐓^B.`$ (23)
Here $`\overline{\delta }_{C,B}1\delta _{C,B}`$ and the (2-cluster) partition transition operator $`𝐓_C`$ is
$`𝐓_C`$ $``$ $`\left(\mathrm{𝟏}𝐕_C𝐑_0\right)^1𝐕_C`$ (24)
$`=`$ $`𝐕_C+𝐕_C𝐑_0𝐓_C=𝐕_C+𝐓_C𝐑_0𝐕_C`$ (25)
$`=`$ $`𝐕_C+𝐕_C𝐑_C𝐕_C.`$ (26)
Therefore the whole $`N`$-particle transition operator $`𝐓`$ can be split into $`𝒩_N`$ Faddeev type transition operators $`𝐓^C`$ associated with the selected set $`𝒞`$ of two-cluster partitions. It is interesting to note that Eq. (22) has the precise structure of the Faddeev equation but is valid for an arbitrary number of particles and with $`𝒩_N`$ components. Just like the original Faddeev equation, the kernel in Eq. (22) is connected after one iteration provided the scattering kernels of the subsystems are already connected.
As in the other treatments of the $`N`$-particle system, 2-cluster transition operators $`𝐓_C`$ may be expressed in terms of unconnected subsystems. Thus it is necessary to ensure that each partition transition operator is connected in order for the total transition operator to have a connected kernel. This can be accomplished by assuring that each cluster transition operator is connected. But first, it is necessary to know how to express the partition transition operators in terms of the corresponding single cluster transition operators. This is done in the next subsection.
### B Cluster transition operators
The partition transition operator $`𝐓_C`$ can be calculated from the transition operators for the two clusters $`C_1`$ and $`C_2`$ that constitute the partition $`C`$. Two approaches for carrying out calculations of this nature have been presented in the literature. First is the approach of Sloan for a 4-particle system, see his Sec. III. This is essentially similar to the method of the last subsection, dividing up the partition transition operator into two parts as in Eq. (17) and obtaining coupled equations for the two parts, as in Eq. (22). The second approach is to separate the partition resolvent into cluster resolvents by the use of a convolution. This method is useful since the two clusters are dynamically independent, so all operators for one cluster commute with all operators of the other cluster. This property of the resolvent was pointed out by Bianchi and Favella and emphasized for use in parts of the 4-particle problem by Haberzettl and Sandhas . Quantities that are used by both methods are defined first, then the methods are discussed in turn.
On the basis that the two clusters in a partition are dynamically independent, the partition potential $`𝐕_C`$ is a sum of cluster components
$$𝐕_C=𝐕_{C_1}+𝐕_{C_2}.$$
(27)
Since the kinetic part, $`𝐇_0`$, of the hamiltonian also separates into cluster components, it follows that the partition hamiltonian also separates into commuting cluster hamiltonians. That is, these hamiltonians are related according to $`(j=1,2)`$
$$𝐇_0=𝐊_{C_1}+𝐊_{C_2},𝐇_{C_j}=𝐊_{C_j}+𝐕_{C_j},𝐇_C=𝐇_{C_1}+𝐇_{C_2}.$$
(28)
In a similar manner, cluster transition operators can also be defined according to
$`𝐓_{C_j}(z^{})`$ $``$ $`𝐕_{C_j}+𝐕_{C_j}{\displaystyle \frac{1}{z^{}𝐊_{C_j}}}𝐓_{C_j}(z^{})`$ (29)
$`=`$ $`𝐕_{C_j}+𝐓_{C_j}(z^{}){\displaystyle \frac{1}{z^{}𝐊_{C_j}}}𝐕_{C_j}`$ (30)
$`=`$ $`𝐕_{C_j}+𝐕_{C_j}{\displaystyle \frac{1}{z^{}𝐇_{C_j}}}𝐕_{C_j}.`$ (31)
What complex parameter $`z^{}`$ is to appear in each cluster transition operator depends on how it is to be used, and differs between the two methods.
In the first method, modelled on Sloan’s approach, the partition transition operator is written as a sum,
$$𝐓_C=𝐓^{C_1}+𝐓^{C_2},$$
(32)
whose parts are defined as
$$𝐓^{C_j}𝐕_{C_j}+𝐕_{C_j}𝐑_0𝐓_𝐂.$$
(33)
Then repeating a process analogous to deriving Eq. (22), the components of the partition transition operator can be shown to satisfy the coupled equations
$$\left(\begin{array}{c}𝐓^{C_1}\\ 𝐓^{C_2}\end{array}\right)=\left(\begin{array}{c}𝐓_{C_1}(z_1)\\ 𝐓_{C_2}(z_2)\end{array}\right)+\left(\begin{array}{cc}0& 𝐓_{C_1}(z_1)\\ 𝐓_{C_2}(z_2)& 0\end{array}\right)𝐑_0\left(\begin{array}{c}𝐓^{C_1}\\ 𝐓^{C_2}\end{array}\right),$$
(34)
where the parameters
$$z_1=z𝐊_2,z_2=z𝐊_1$$
(35)
have been chosen so that the energy factors are consistent with the properties of the partition transition operator. According to this, the partition transition operator can be expressed as the series
$$𝐓_C=𝐓_{C_1}(z_1)+𝐓_{C_2}(z_2)+𝐓_{C_1}(z_1)𝐑_0𝐓_{C_2}(z_2)+𝐓_{C_2}(z_2)𝐑_0𝐓_{C_1}(z_1)+\mathrm{}.$$
(36)
The cluster transition operators $`𝐓_{C_j}(z_j)`$ can be evaluated for arbitrary $`z_j`$ entirely in the respective cluster subspace involving the states of $`\frac{N}{2}`$ or $`\frac{N1}{2}`$ particles in $`C_j`$. This can be accomplished in the same manner as that described for the $`N`$ particle system since the cluster is dynamically independent. But it must be evaluated with the appropriate $`z_j`$ parameter when used in calculating the partition transition operator.
In the second method the basic starting point is to express the partition resolvent as the convolution
$$𝐑_C=\frac{1}{2\pi i}_\mathrm{\Gamma }\frac{dz^{}}{(z^{}𝐇_{C_1})(zz^{}𝐇_{C_2})}$$
(37)
of the resolvents of the two clusters. Here the contour $`\mathrm{\Gamma }`$ is to be the straight line from $`\mathrm{}`$ to $`\mathrm{}`$, but lying above the real axis and below $`z`$. Thus $`z^{}`$ must be such that $`0<\mathrm{}(z^{})<\mathrm{}(z)`$. It follows that the partition transition operator is given by
$$𝐓_C=𝐕_{C_1}+𝐕_{C_2}+\frac{1}{2\pi i}_\mathrm{\Gamma }(𝐕_{C_1}+𝐕_{C_2})\frac{dz^{}}{(z^{}𝐇_{C_1})(zz^{}𝐇_{C_2})}(𝐕_{C_1}+𝐕_{C_2}).$$
(38)
This can be expressed in many different ways. The following emphasizes how the integrand can be expressed in terms of the cluster transition operators, Eq. (29), but leaves the contour integral unchanged.
The integrand consists of four terms, according to the different cluster potentials. The diagonal in cluster potential terms can immediately be recognized as related to the corresponding cluster transition operator, whereas the remainder needs the identity
$$𝐕_{C_1}\frac{1}{z^{}𝐇_{C_1}}=𝐓_{C_1}(z^{})\frac{1}{z^{}𝐊_{C_1}}$$
(39)
and its other various combinations. After some calculation, the integrand can be written in the form
$`(𝐕_{C_1}+𝐕_{C_2}){\displaystyle \frac{1}{(z^{}𝐇_{C_1})(zz^{}𝐇_{C_2})}}(𝐕_{C_1}+𝐕_{C_2})`$ (40)
$`=`$ $`[𝐓_{C_1}(z^{})𝐕_{C_1}]{\displaystyle \frac{1}{zz^{}𝐇_{C_2}}}+[𝐓_{C_2}(zz^{})𝐕_{C_2}]{\displaystyle \frac{1}{z^{}𝐇_{C_1}}}`$ (42)
$`+𝐕_{C_1}{\displaystyle \frac{1}{(z^{}𝐇_{C_1})(zz^{}𝐇_{C_2})}}𝐕_{C_2}+{\displaystyle \frac{1}{z^{}𝐇_{C_1}}}𝐕_{C_1}𝐕_{C_2}{\displaystyle \frac{1}{zz^{}𝐇_{C_2}}}`$
$`=`$ $`[𝐓_{C_1}(z^{})𝐕_{C_1}]{\displaystyle \frac{1}{zz^{}𝐊_{C_2}}}+[𝐓_{C_2}(zz^{})𝐕_{C_2}]{\displaystyle \frac{1}{z^{}𝐊_{C_1}}}`$ (47)
$`+[𝐓_{C_1}(z^{})𝐕_{C_1}]{\displaystyle \frac{1}{zz^{}𝐊_{C_2}}}𝐓_{C_2}(zz^{}){\displaystyle \frac{1}{zz^{}𝐊_{C_2}}}`$
$`+[𝐓_{C_2}(zz^{})𝐕_{C_2}]{\displaystyle \frac{1}{z^{}𝐊_{C_1}}}𝐓_{C_1}(z^{}){\displaystyle \frac{1}{z^{}𝐊_{C_1}}}`$
$`+𝐓_{C_1}(z^{}){\displaystyle \frac{1}{(z^{}𝐊_{C_1})(zz^{}𝐊_{C_2})}}𝐓_{C_2}(zz^{})`$
$`+{\displaystyle \frac{1}{z^{}𝐊_{C_1}}}𝐓_{C_1}(z^{})𝐓_{C_2}(zz^{}){\displaystyle \frac{1}{zz^{}𝐊_{C_2}}}.`$
Most terms involve the product of the transition operators for the two clusters, weighted with different combinations of the free particle resolvents for the clusters. It is up to the computational method to decide which way these are to be evaluated. But the first two terms each involves only one cluster transition operator. The contour integral of these two terms can easily be done. For the first term the combination $`𝐓_{C_1}(z^{})𝐕_{C_1}`$ is analytic in the upper half $`z^{}`$-plane and vanishes for $`|z^{}|\mathrm{}`$, so closing the contour at $`\mathrm{}(z^{})+\mathrm{}`$ contributes a pole only when $`z^{}=z𝐊_{C_2}`$. Analogously for the second term, so that the partition transition operator can be written
$$𝐓_C=𝐓_{C_1}(z_1)+𝐓_{C_2}(z_2)+\frac{1}{2\pi i}_\mathrm{\Gamma }\left\{𝐓_{C_1}𝐓_{C_2}\mathrm{terms}\right\}.$$
(48)
While this is only one way of organizing this result there are many other ways in which the cluster transition operator expansion of the partition transition operator could be written and no attempt is made here to catalog all the possibilities. In comparing the two methods, Eq. (36) is an infinite series in cluster transition operators, while Eq. (48) is only quadratic in cluster transition operators, but requires an integration over how the energy is divided up between the two clusters.
In this way the calculation of the partition transition operator has been reduced to the independent computation of the cluster transition operators. The advantage of this decomposition is that the individual cluster transition operators deal with isolated sets of particles whose number is less than $`N`$. Such a cluster of particles can be decomposed into partitions as was the original problem, and the whole process repeated. Specifically, for $`N=2^n`$ even, the successive problems deal with $`2^k`$ particles, $`k=n1,n2,\mathrm{},1`$.
## IV GENERALIZED FADDEEV COMPONENTS
As a function of $`z`$, the resolvent operator $`𝐑`$ has singularities at the spectrum of the system hamiltonian, thus it is of practical use for finding solutions of the Schrödinger equation, for determining the dynamical evolution of the system, and identifying normalizable resonance states. Faddeev’s approach is followed in order to express the $`N`$-particle resolvent operator in terms of selected two-cluster partition transition operators. This representation of the total resolvent operator is then used to decompose a scattering wave function originating from a particular incoming two-cluster partition. The resulting wave functions are here referred to as generalized Faddeev components. The homogeneous system of equations associated with the generalized Faddeev components is shown explicitly to solve the Schrödinger equation. This section first describes the generalized Faddeev equations for the $`N`$-particle resolvent operator and subsequently discusses the associated wave functions.
The total resolvent operator $`𝐑`$ is related to the total transition operator $`𝐓`$ via the Lippmann-Schwinger equation
$`𝐑`$ $`=`$ $`𝐑_0+𝐑_0\mathrm{𝐕𝐑}=𝐑_0+𝐑_0\mathrm{𝐓𝐑}_0.`$ (49)
On taking over the partition expansion (17) of $`𝐓`$, this can be written as
$$𝐑=𝐑_0+\underset{C𝒞}{}𝐑^C,$$
(50)
where the Faddeev type resolvent $`𝐑^C`$ is given by
$$𝐑^C𝐑_0𝐓^C𝐑_0.$$
(51)
The generalized Faddeev equations (22) lead to the coupled set of equations
$`𝐑^C`$ $`=`$ $`𝐑_0𝐓_C𝐑_0+{\displaystyle \underset{B𝒞}{}}\overline{\delta }_{C,B}𝐑_0𝐓_C𝐑_0𝐓^B𝐑_0`$ (52)
$`=`$ $`𝐑_C𝐑_0+{\displaystyle \underset{B𝒞}{}}\overline{\delta }_{C,B}𝐑_0𝐓_C𝐑^B`$ (53)
with partition resolvent operator defined in (12) and related to the partition transition operator by
$$𝐑_C=𝐑_0+𝐑_0𝐕_C𝐑_C=𝐑_0+𝐑_0𝐓_C𝐑_0.$$
(54)
For $`N=3`$, Eq. (52) is Faddeev’s decomposition of the total resolvent operator whose kernel is not connected until after one iteration. For $`N>3`$, the kernel of Eq. (52) requires one iteration in the same manner providing the kernels of the subsystems are already connected. These relations between the transition operators and the resolvent operators are analogous to the original Faddeev equations for the three-particle resolvent operator. The set of resolvent equations can also be used for the evaluation of statistical mechanical virial coefficients.
A scattering wave function for the $`N`$-particle system is determined by the total resolvent operator $`𝐑`$ according to
$$\mathrm{\Psi }_C=\underset{\epsilon 0}{lim}i\epsilon 𝐑(E+i\epsilon )\varphi _C.$$
(55)
This is applied here to the selected set $`𝒞`$ of compatible two-cluster partitions, with $`\varphi _C`$ a stationary solution of the two-cluster partition hamiltonian $`𝐇_C`$ of energy $`E`$,
$$𝐇_C\varphi _C=E\varphi _C,$$
(56)
which may be a distorted wave of the two-cluster subsystem.
The detailed structure of the scattering wave function $`\mathrm{\Psi }_C`$ of Eq. (55) is now discussed. This is begun by first applying the resolvent expansions (50) and (52) to Eq. (55),
$`\mathrm{\Psi }_C`$ $`=`$ $`\underset{\epsilon 0}{lim}i\epsilon 𝐑_0(E+i\epsilon )\varphi _C+{\displaystyle \underset{B𝒞}{}}\underset{\epsilon 0}{lim}i\epsilon 𝐑^B(E+i\epsilon )\varphi _C.`$ (57)
The first term on the right hand side contributes only if $`C`$ is the channel with all particles free, namely
$$\underset{\epsilon 0}{lim}i\epsilon 𝐑_0(E+i\epsilon )\varphi _C=\varphi _C\delta _{C,C_N},$$
(58)
where $`C_N`$ is the $`N`$-cluster partition. The second term can be written as a sum of generalized Faddeev components
$$\underset{\epsilon 0}{lim}i\epsilon 𝐑^B(E+i\epsilon )\varphi _C\psi _{BC}.$$
(59)
According to Eq. (52), the generalized Faddeev component $`\psi _{BC}`$ can be expanded as
$`\psi _{BC}`$ $`=`$ $`\underset{\epsilon 0}{lim}i\epsilon 𝐑_B\varphi _C\underset{\epsilon 0}{lim}i\epsilon 𝐑_0\varphi _C`$ (60)
$`+`$ $`{\displaystyle \underset{A𝒞}{}}\overline{\delta }_{A,B}𝐑_0𝐓_B\underset{\epsilon 0}{lim}i\epsilon 𝐑^A\varphi _C`$ (61)
$`=`$ $`\varphi _C\delta _{B,C}\varphi _C\delta _{C_N,C}+{\displaystyle \underset{A𝒞}{}}\overline{\delta }_{A,B}𝐑_0𝐓_B\psi _{AC}`$ (62)
$`=`$ $`\varphi _C\delta _{B,C}\varphi _C\delta _{C_N,C}+{\displaystyle \underset{A𝒞}{}}\overline{\delta }_{A,B}𝐑_B𝐕_B\psi _{AC}.`$ (63)
Here the identities
$$\underset{\epsilon 0}{lim}i\epsilon 𝐑_B\varphi _C=\varphi _C\delta _{B,C}$$
(64)
and
$$𝐑_0𝐓_B=𝐑_B𝐕_B$$
(65)
have been used as well as the definition in Eq. (59). Thus the scattering wave function, Eq. (57), is expressed in terms of the generalized Faddeev components (60),
$$\mathrm{\Psi }_C=\varphi _C\delta _{C,C_N}+\underset{B𝒞}{}\psi _{BC}.$$
(66)
For $`N`$=3 this expression for the scattering state is exactly the wave function originally derived by Faddeev. Hence Eq. (66) is the generalization of Faddeev’s scattering state based on a chosen set $`𝒞`$ of two-cluster partitions of $`N`$ particles.
It is easily shown that Eq. (66) satisfies the Schrödinger equation (2). A special case is the homogeneous analog. The proof that this formally satisfies the Schrödinger equation is as follows:
$`\mathrm{\Psi }_C`$ $`=`$ $`{\displaystyle \underset{B𝒞}{}}\psi _{BC}`$ (67)
$`=`$ $`{\displaystyle \underset{B𝒞}{}}\left({\displaystyle \underset{A𝒞}{}}\overline{\delta }_{A,B}𝐑_B𝐕_B\psi _{AC}\right)`$ (68)
$`=`$ $`{\displaystyle \underset{B𝒞}{}}\left({\displaystyle \underset{A𝒞}{}}\overline{\delta }_{A,B}𝐑_0𝐕_B\left[\mathrm{𝟏}+𝐑_B𝐕_B\right]\psi _{AC}\right)`$ (69)
$`=`$ $`{\displaystyle \underset{B𝒞}{}}𝐑_0𝐕_B\left({\displaystyle \underset{A𝒞}{}}\overline{\delta }_{A,B}\psi _{AC}+\psi _{BC}\right)`$ (70)
$`=`$ $`{\displaystyle \underset{B𝒞}{}}𝐑_0𝐕_B{\displaystyle \underset{A𝒞}{}}\psi _{AC}`$ (71)
$`=`$ $`{\displaystyle \underset{B𝒞}{}}𝐑_0𝐕_B\mathrm{\Psi }_C`$ (72)
$`=`$ $`𝐑_0𝐕\mathrm{\Psi }_C.`$ (73)
Of course a homogeneous solution of the Schrödinger equation occurs only for bound states energies, so such a solution plays no role when solving a scattering problem.
In general, a cluster can be either a bound, or an unbound but interacting, set of particles, with a free particle included as a (1-particle) bound state. If both clusters in a partition are bound then the partition is a true asymptotic channel. Whether this is or is not the case, the partition wave function has one of two forms, depending on whether one of the clusters is or is not a free particle, namely
$$\varphi _C=\{\begin{array}{cc}\phi _{C^1},𝐪_C\hfill & \\ \phi _{C^1},\phi _{C^2},𝐪_C.\hfill & \end{array}$$
(74)
Here $`\phi _{C^1}`$ is a bound state of the $`N1`$-particle subsystem and $`𝐪_C`$ is the momentum generalized eigenstate for the relative motion of the bound and free states, on the basis that the total center of mass momentum of the $`N`$-particle system has been removed from discussion. In the second form, the $`2`$nd cluster is now also a bound state, which must be explicitly indicated, while $`𝐪_C`$ is again the relative momentum. If a cluster is unbound, but all particles interacting, then the corresponding cluster wavefunction must be replaced by $`\phi _{C^j}^+`$, corresponding to the scattering wavefunction from some initial incoming state. It appears intuitively reasonable that this procedure is correct, though its complete mathematical justification may be needed since some limits must already have been taken in order to define (in general, a product of) scattering states with definite energy as an input, with the next step involving a further limit.
## V The 4- and 8-particle systems
The treatment so far has emphasized how the $`N`$-particle scattering problem can be broken down into a number of problems involving fewer particles. Once a breakdown has been selected, it is then a case of building up the transition operator and scattering wavefunction for the $`N`$-particle system. The procedure is illustrated by discussing systems of 4 and 8 particles. For definiteness it is assumed that there are no bound states for any number of particles.
For the specific case of there being four particles, these will be labelled 1, 2, 3 and 4. Then according to Eq. (3) there are 3 two-cluster partitions, which are the three pairs of pair particles 12,34, 13,24 and 14,23. Moreover, Eq. (6) states that no scaling of the potential is required. Thus the procedure is to first find the two-particle transition operator $`𝐓_{12}(z)`$ for arbitrary complex $`z`$ in the upper half plane. On the basis that all particles are the same species, this transition operator is the same for any other pair, except for the labelling. The second step is then to find the partition transition operator for the two-cluster partition 12,34. This is accomplished according to Sec. IIIB, specifically involving either the first approach, Eq. (36), or the second approach, Eq. (38), with its integrand expressed in terms of the pair particle transition operators, Eq. (40). Again this calculation needs to be done only once since a relabelling immediately gives the partition transition operator for the other two partitions. The last step for getting the four-particle transition operator is to solve the $`3\times 3`$ matrix equations (22) and add the results, Eq. (17). Once the transition operator is known, the various generalized Faddeev components can be calculated, as described in Sec. IV, and the desired scattering amplitude calculated. In this way the four-particle scattering problem is very similar to the three-particle scattering problem, involving only three partitions. The extra complexity is in the added structure of the partition transition operators. In contrast, the method of Sloan involves a $`7\times 7`$ matrix while the Yakubovsky approach organizes the wavefunction into 18 components. The approach presented here would seem to be both simpler and more efficient.
The 8-particle problem is discussed with the view of clarifying how the proposed approach works for more complicated systems. Equation (3) states that there are 35 two-cluster partitions, each cluster having 4 particles. Moreover, Eq. (6) states that $`v_8=1/15`$, so the potential needs to be scaled by this fraction. This scaling is required when finally solving the $`35\times 35`$ matrix, Eq. (22), for the components of the 8-particle transition operator, so enters into the calculation of all (the 8-particle, partition, cluster and subpartition and subcluster) transition operators. Now a typical 8-particle partition transition operator is needed, which is a combination of two 4-particle cluster transition operators. The latter must be found by solving the problem discussed in the previous paragraph, but now with the potential scaled by $`v_8`$. Since there is no further scaling required in solving for the 4-particle transition operator, it is only the $`v_8`$ scaling that must be applied. Thus the detailed procedure that is to be followed is, in sequence: 1) scale the pair potential by $`v_8`$; 2) calculate the 4-particle transition operator for this scaled potential as described in the last paragraph; 3) use either of Eqs. (36) or (38) to obtain a two-cluster partition transition operator for the 8-particle system, with the second cluster transition operator obtained from the first 4-particle cluster transition operator by relabelling; 4) obtain the 34 other partition transition operators by relabelling and solve the generalized Faddeev equations (22).
## VI DISCUSSION
A proposal has been made for the generalization of the Faddeev equations to an arbitrary number, $`N`$, of particles. Essentially this involves selecting a set of partitionings of the $`N`$ particles into pairs of clusters and expressing the scattering properties of the total system in terms of the scattering properties of the clusters. Each cluster is treated in the same way so that after successive treatments an individual cluster contains either only a single particle or a pair of particles. Although mathematically, other (combinations of) sets of two-cluster partitions lead to a connected kernel as well in the present method, the detailed presentation of this paper has stressed the choice of the set of two-cluster partitions in which the pair of clusters in any partition is of nearly equal size. This should be the most efficient procedure since then there are fewer types of clusters that need to be treated.
The proposed method requires, in general, a scaling of the potential and expressing the transition operator, and wave function, in terms of partition transition operators which are in turn expressed in terms of the transition operators for the two individual clusters composing the partition. An essential simplifying feature is that the transition operator for an individual cluster can be calculated independently of the presence of the other particles, so is equivalent to solving the scattering of a system of fewer particles, which can in general be treated in the same manner as the original system.
For most $`N`$-particle integral equation theories, the starting point is the pair particle interaction or the two-particle transition operator, which is then combined in all possible ways with the transition operators of other pairs of particles. In order to assure that all particles are connected, that is, no disconnected diagrams appear, it is necessary to iterate through a whole sequence of processes. This complicates the description of the $`N`$-particle problem as it appears, for example, in the Yakubovsky equations. In contrast, the present method starts by dividing the $`N`$-particle system into a selected set of two-cluster partitions. The treatment assures that all clusters are connected to each other. There remains the possibility that there is a disconnectionness within a cluster. This is eliminated by repeating the procedure for each cluster as if it were a separate scattering system. In this way the formulation has a very simple structure which is exactly similar at each stage to the Faddeev equations, in particular reducing to them if $`N=3`$.
The total resolvent for the $`N`$-particle system is decomposed in a manner analogous to Faddeev’s method. This decomposition determines the decomposition of the scattering wave functions, viz. Eq. (66), into what are here called generalized Faddeev components. The sum of the corresponding homogeneous set of generalized Faddeev components is explicitly shown to be completely equivalent to the Schrödinger equation. Such a property was explicitly demonstrated in Faddeev’s theory and Yakubovsky’s theory but has not been explicitly shown in some other theories. This property is of particular importance for finding the (discrete) eigenvalues for the $`N`$-particle system and for avoiding spurious solutions (for bosons and fermions, an appropriate symmetrization of the Faddeev components may be required). In particular, Federbush first found a spurious solution for the Weinberg equation and a systematic study by Glöckle and coworkers indicated that most few-body equations admit the existence of discrete spurious solutions. However these spurious solutions were not a problem in finding scattering solutions.
The cluster and partition wavefunctions from all the stages are needed before the total wavefunction can be obtained. In general these are nonphysical in that they represent intermediary results for the final calculation. The same is true even more so for the Faddeev components. In contrast, if both clusters in a partition represent bound states, then this is a valid asymptotic condition and the channel wavefunction does represent a physical state for the corresponding scattering system.
## ACKNOWLEDGMENTS
The authors thank the referee for drawing their attention to the method of Haberzettl and Sandhas for combining dynamically independent transition operators, which procedure has been incorporated into the presentation in Sec. IIIB. This work was supported in part by the Natural Sciences and Engineering Research Council of Canada. G.W.W. thanks the Killam Foundation of Canada for a fellowship and the National University of Singapore for research funds.
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