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# Recent results on Charm Physics from FermilabThis work is partially supported by CLAF/CNPq Brazil and CONACyT México To appear, Proc. VII Mexican Workshop on Particles and Fields, Merida Yuc. Méx., Nov.1999
## Introduction
This review contains recent results from three Fermilab fixed target experiments dedicated to study charm physics: E791 and E871 (SELEX) of charm hadroproduction and E831 (FOCUS) of charm photoproduction.
The experiment E791 used a 500 GeV/c $`\pi ^{}`$ beam incident on platinum and diamond target foils and took data in 1991-1992. A loose transverse energy trigger was used to record $`2^{10}`$ interactions. Silicon microstrip detectors (6 in the beam and 17 downstream of the target) provided precision track and vertex reconstruction. The precise location of production and decay vertices of long lived charm particles allowed to reconstruct over $`2^5`$ charm particles.
SELEX used a 600 GeV hyperon beam of negative polarity to make a mixed beam of $`\mathrm{\Sigma }^{}`$ and $`\pi ^{}`$ in roughly equal numbers. The positive beam was composed of 92 $`\%`$ of protons, and 8 $`\%`$ of $`\pi ^+`$. Interactions occurred in a segmented target of 5 foils, 2 Cu and 3 C. The experiment was designed to study charm production in the forward hemisphere, with good mass and decay vertex resolution for charm momenta in the range 100-500 GeV/c. A major innovation was the use of online selection criteria to identify events that had evidence for a secondary vertex. Data taking finished in 1997.
The FOCUS experiment used a photon beam with $`<E_\gamma >=170`$ GeV on a Beryllium Oxide segmented target. Sucessor of Fermilab E687, FOCUS had upgrades in vertexing, Cerenkov identification, electromagnetic calorimetry and muon identification. The target segmentation and the use of additional silicon microstrip detectors after each pair of target segments were major improvements. The experiment took data in 1996-1997 and collected over 7 billion triggers of charm candidates. FOCUS has the largest charm sample available in the world, with over $`10^6`$ fully reconstructed charmed particles.
The data available from these three experiments will allow to make high precision charm physics, to study rare decays and to search for new physics.
## Charm production and cross sections for D mesons
Charm production is a combination of short and long range processes. Perturbative Quantum Chromodynamics (pQCD) can be used to calculate the parton cross section, the short range process. The long range process of parton hadronization has to be modeled from experimental data. The two processes occur at different time scales and should not affect each other, leading to factorization properties. The general expression for production of charm in hadroproductions is crossection
$$\sigma (P_A,P_B)=\underset{i,j}{}𝑑x_A𝑑x_Bf_i^A(x_A,\mu )f_j^B(x_B,\mu )\widehat{\sigma }_{i,j}(\alpha _s(\mu ),x_AP_A,x_BP_B)$$
(1)
where $`f_i^A(x_A,\mu )`$ is the probability of finding parton type $`i`$ inside the hadron $`A`$ with momentum fraction $`x_A`$, and $`\mu `$ is the scale at which the process occurs. The other hadron taking part in the interaction has $`f_j^B(x_B,\mu )`$. The elementary parton cross section $`\widehat{\sigma }_{i,j}(\alpha _s(\mu ),x_AP_A,x_BP_B)`$ term is calculated by QCD, according to parton models. Although QCD is well defined theory, solutions to most problems are quite difficult. So, to calculate the parton cross section we usually do a perturbative series expansion in terms of the strong interaction coupling constant, $`\alpha _s`$, and calculate at leading order (LO), next-to-leading order (NLO) and so on. Nowadays most calculations include NLO. Some Feynman diagrams for hadroproduction of charm at LO and NLO are shown in Fig.1 (more details can be found in reference crossection ).
Usually the differential cross sections is measured as a function of the scaled longitudinal momentum (Feynman $`x_F`$) $`d\sigma /dx_F`$, and as a function of the transverse momentum squared $`d\sigma /dp_t^2`$. A phenomenological parametrization of the double differential cross section often used to describe the experimental data is given by,
$$\frac{d\sigma }{dx_Fdp_t^2}=A(1x_F)^nexp(bp_t^2)$$
(2)
where $`A,n`$ and $`b`$ are parameters used to fit the data. In fact kinematic considerations provide predictions for the power $`n`$ at high $`x_F`$ and QCD calculations give average $`p_t`$ values comparable to the charm quark mass.
E791 has measured the total forward cross section and the differential cross sections as functions of $`x_F`$ and $`p_t^2`$ from a sample of $`88990\pm 460`$ $`D^0`$ mesons. Fig. 2 shows the differential distributions and its comparison to theoretical predictions from QCD calculations at NLO mangano and to the Monte Carlo event generator, Pythia/Jetset pythia for $`c\overline{c}`$ and $`D\overline{D}`$ production. Hadron distributions are softer than $`c\overline{c}`$ distributions due to fragmentation. With suitable choice for the intrinsic transverse momentum $`k_t`$ of the incoming partons ($`k_t=1`$ GeV/c), the Peterson fragmentation function parameter ($`\epsilon =0.01`$), and the charm quark mass $`m_c=1.5`$ GeV/c<sup>2</sup>, NLO $`D`$ meson calculation provides a good match to the $`p_t^2`$ distribution and fair match to the $`x_F`$ distribution. The hadronization scheme implemented in Pythia/Jetset can be adjusted to fit the data. The large number of $`D^0`$ events make it possible to clearly observe a turnover point greater than zero ($`x_c=0.0131\pm 0.0038`$) in the $`x_F`$ distribution. The positive value provide evidence that the gluon distribution in the pion is harder than the gluon distribution in the nucleon cross-D0 . The total forward neutral D meson cross section measured by E791 is $`\sigma (x_F>0)=15.4\pm _{2.3}^{1.8}\mu b`$/nucleon.
SELEX has also preliminary data on $`D^0`$ cross section. Fig. 2 (right) shows the $`x_F`$ distribution for $`D^0`$ \+ c.c. from E791 cross-D0 and the preliminary data from SELEX. Up to $`x_F=0.5`$ data overlap and agree well. However SELEX has no strong evidence for rise at large $`x_F`$ as seen in the E791 data. The continuous line is a fit to the data points with the phenomenological parametrization in Eq. 2.
### Correlations in the production of Charm pairs
Observation of both of the charm particles in a hadron-produced event can give additional information on the production process and allow to test QCD predictions since both longitudinal and transverse momenta of the charm particles and angular correlations are explicitly measured.
In the simplest parton model the charm and anticharm particles are expected to be produced in opposite directions in the transverse plane. However if one assumes that the incoming partons have an intrinsic transverse momentum, $`k_t`$, this will affect the transverse momentum of the heavy quark pair, its azymuthal correlations and the transverse distribution of a single quark. The partons entering the hard interaction are indeed supposed to have a non vanishing primordial $`k_t`$, seen as a nonperturbative Fermi motion of partons inside the incoming hadrons. Typical values of $`k_t`$ should thus be 300-400 MeV. However it has been noted pairs-correlation that much higher values of $`k_t`$ are required, at or above 1 GeV, to reproduce charm data. This could be an indications of the importance of next to leading and higher order effects, by which emitted gluons would further modify the nearly back-to-back Jeff-Sjostrand production of the final charm hadrons.
E791 has measured correlations between D and $`\overline{D}`$ mesons from $`791\pm 44`$ fully reconstructed charm meson pairs correlation-E791 . The main variables used to describe charm pair correlations are the beam direction, $`p_t`$, and either the rapidity $`y`$ or the Feynman scaling variable $`x_F`$, and also the azymuthal distribution $`\varphi `$. The same way the difference and sum between these variables for two charm mesons ($`D,\overline{D}`$) $`\mathrm{\Delta }x_F`$, $`\mathrm{\Sigma }x_F`$ and so on.
The measured distributions are compared to predictions of the fully differential NLO calculation for $`c\overline{c}`$ production mangano , as well as to predictions from the Pythia/Jetset Monte Carlo event generator for $`c\overline{c}`$ and $`D\overline{D}`$ production. Default parameter have been used for the theoretical models.
For the single charm distributions shown in Fig.3, we observe that for the longitudinal momentum distributions $`x_F`$ and $`y`$ the experimental results and theoretical predictions do not agree. In this comparison the experimental distributions are most similar to the NLO and Pythia/Jetset $`c\overline{c}`$ distributions, but are narrower than all three theoretical predictions. The experimental $`p_t^2`$ distribution agrees quite well with all three theoretical distributions. As expected, both the theoretical and experimental $`\varphi `$ distributions are consistent with being flat.
The experimental and theoretical longitudinal distributions for pairs $`\mathrm{\Delta }x_F`$ , $`\mathrm{\Sigma }x_F`$, $`\mathrm{\Delta }y`$ and $`\mathrm{\Sigma }y`$ are shown in Fig. 4. As with the single-charm distributions, the experimental results are much closer to the two $`c\overline{c}`$ predictions than to the Pythia/Jetset $`D\overline{D}`$ predictions, but narrower than all three theoretical predictions. For the transverse distributions for charm pairs any observed discrepancy between data and theory must derive from the theory modeling the correlations between the transverse momentum of the two $`D`$ mesons $`p_{t,D}`$ and $`p_{t,\overline{D}}`$ because the single charm $`p_t^2`$ and $`\varphi `$ experimental distributions agree well with the theory. The $`\mathrm{\Delta }\varphi `$ distribution shows clear evidence of correlations (more details can be found in reference correlation-E791 ).
## Hadronization and Particle - antiparticle asymmetries
The production of a charm hadron can be subdivided in two steps: the production of a $`c\overline{c}`$ pair followed by the hadronization of these quarks. In perturbative QCD, that describes the $`c\overline{c}`$ production, the $`x_F`$ spectra of produced charm/anticharm quarks are identical to leading order and the effects of higher orders are very small in this respect. Therefore any asymmetry between charm and anticharm hadrons is a simple measure of nonperturbative effects coming from the hadronization process.
Particle - antiparticle asymmetries can be quantified by means of the asymmetry parameter
$$A=\frac{N\overline{N}}{N+\overline{N}}$$
(3)
where $`N`$ ($`\overline{N}`$) is the number of produced particles (antiparticles). This parameter is usually measured as a function of $`x_F`$ and $`p_t^2`$.
Several experiments have reported an enhancement in the production rate of charm particles having valence quarks in common with the incident particles, relative to charge conjugate particles which have fewer or no common valence quarks. This effect is known as leading particle effect. Measurements of the asymmetry parameter A can put in evidence leading particle effects, as well as other effects like associated production of a meson and a baryon.
From the theoretical point of view, models which can account for the presence of leading particle effects use some kind of non-perturbative mechanism for hadronization, in addition to the perturbative production of charm quarks. Two examples are:
String fragmentation strings : in this case the parton of the hard interaction and the beam remnants are connected by a string which reflect the confining color field. Successive breaking of the color flux tube stretched between a cluster, when it is kinematic possible, will create light quark-antiquark and hadrons will be produced.
Intrinsic charm model intrinsic : a virtual $`c\overline{c}`$ pair pops from the sea of the beam particle. The $`c\overline{c}`$ pair coalesce with the neighbor valence quarks due to their similar rapidity. This mechanism favors the production of charm particles with valence quarks in common with the beam particle at high $`x_F`$ and low $`p_t`$ region. A similar argument can be drawn with respect to the target.
New results on the asymmetry parameter $`A`$ and evidence for leading particle effects in both meson and baryon production are available from experiments E791, SELEX and FOCUS.
E791 has measured recently the asymmetries of $`D_s^\pm `$ mesons e791-ds-asi . Fig. 5 shows the $`D_s^\pm `$ asymmetry as a function of $`x_F`$ and $`p_t^2`$, compared with previous $`D^\pm `$ results from the same experiment Asim-E791 .
Preliminary results of $`D^0(\overline{D^0})`$ as well as $`D^\pm `$ asymmetries as a function of $`x_F`$ presented by SELEX Lori are shown in Fig. 6 (left), for different incident beam particles ($`\pi ^{}`$, $`\mathrm{\Sigma }^{}`$, $`p`$). The $`\pi ^{}`$ data is compared to $`D^\pm `$ asymmetries from E791. The asymmetry at $`x_F>\mathrm{\hspace{0.25em}0.4}`$ does not rise steeply with $`x_F`$ as previously reported by E791, Fig. 6 (right).
For baryons there are also preliminary results for the asymmetry parameter. Fig.7 (left) shows the E791 results for the $`\mathrm{\Lambda }_c^+`$ asymmetries as function of $`x_F`$ and $`p_t^2`$, compared with predictions from Pythia/Jetset (full lines) magnin . The results show a uniform positive asymmetry of $`12.7\pm 3.4\%`$ over the studied kinematical range but do not exclude a rise in the $`x_F<\mathrm{\hspace{0.25em}0}`$ region as predicted by Pythia/Jetset. For $`x_F>0`$ the observed asymmetry does not agree with Pythia/Jetset predictions.
SELEX has also measured the $`\mathrm{\Lambda }_C^+`$ asymmetry as a function of $`x_F`$ for different incident beam particles ($`\pi ^{},\mathrm{\Sigma }^{},p`$) Fig. 7 (right) Lori . The asymmetry is clearly larger for the baryon beams than for the $`\pi `$ beam. For the protons the only region in which there is $`\overline{\mathrm{\Lambda }}_c`$ production is at very small $`x_F`$. Their preliminary results for the $`\pi ^{}`$ beam are compatible with those from E791, Fig. 7 (left).
FOCUS has also preliminary results on baryon asymmetries. In Fig.8 we see the $`\mathrm{\Lambda }_c`$ asymmetry as functions of $`p_l`$, $`p_t^2`$ and $`x_F`$ obtained from a sample of about 16,000 $`\mathrm{\Lambda }_c^{}s`$, compared with Pythia/Jetset Predictions. As FOCUS has a photon beam, no leading particle effect is expected in the $`x_F>0`$ region. In this case the positive asymmetry observed in all the $`x_F`$ range can be an indication of charm baryon and charm meson associated production, favouring a positive asymmetry.
The high statistic $`\mathrm{\Lambda }_c`$ sample from FOCUS allowed to obtain about 600 $`\mathrm{\Sigma }_c\mathrm{\Lambda }_c\pi `$. It is interesting to compare the asymmetry for charm particles with different light quark content. We present in Fig.9 (right) very preliminary results from FOCUS comparing the $`\mathrm{\Sigma }_c^{++}(uuc)`$ and $`\mathrm{\Sigma }_c^0(ddc)`$ to the $`\mathrm{\Lambda }_c^+(udc)`$ total asymmetry.
In Fig. 9 (left) we show the comparison between the $`\mathrm{\Lambda }^0`$ and the $`\mathrm{\Lambda }_c`$ asymmetries as a function of $`x_F`$ from E791. Their similarity suggests that the $`ud`$ diquark shared by the produced $`\mathrm{\Lambda }^0`$ ($`\mathrm{\Lambda }_c^+`$) and nucleons in the target should play an important role in the measured asymmetry in the $`x_F<0`$ region. However, one expects that $`\mathrm{\Lambda }_c`$ asymmetry grows more slowly than the $`\mathrm{\Lambda }^0`$ asymmetry due to its higher mass.
Leading particle effects were also seen by E791 in hyperon production. Preliminary results on $`\mathrm{\Lambda }`$, $`\mathrm{\Xi }`$ and $`\mathrm{\Omega }`$ asymmetries as a function of $`x_F`$ and $`p_t^2`$ are shown in Fig.10 in comparison with predictions from Pythia/Jetset. The range of $`x_F`$ covered allowed the first simultaneous study of the asymmetry in both the negative and positive $`x_F`$ regions. We can clearly see leading particle effects associated with the target or with the beam particles which qualitatively agree with expectations from recombination models (see table 1) Ahyper . It is interesting to observe, as expected, the crossover of the $`\mathrm{\Xi }`$ asymmetry with respect to the $`\mathrm{\Lambda }`$ asymmetry at $`x_F\mathrm{\hspace{0.25em}0}`$. The positive asymmetry measured in regions $`x_F>0`$ for the $`\mathrm{\Lambda }(udc)`$ and for the $`\mathrm{\Omega }(sss)`$ suggest the associated production of a hyperon and a kaon due to the higher energy threshold imposed by baryon number conservation for the production of an anti-hyperon. Pythia/Jetset does not reproduce the data.
## Double Cabibbo suppressed decays
The Cabibbo suppressed charm decays can provide useful insights into the weak interaction mechanism for nonleptonic decays. The $`D^+K^+\pi ^{}\pi ^+`$ signal obtained from 100 $`\%`$ of FOCUS data set consist of $``$ 300 events and is at least a factor of five larger than two previous observations by E687 and E791 ( E791 observed $`59\pm 13`$ events e791-dcsd ). The preliminary branching ratio relative to $`K^{}\pi ^+\pi ^{}`$ is ($`0.72\pm 0.09`$) $`\%`$, completely consistent with the world average of ($`0.68\pm 0.15`$)$`\%`$ and the E791 values of ($`0.77\pm 0.17\pm 0.08`$)$`\%`$. We note that this is $``$ $`3tan^4(\theta _c)`$, which is roughly the ratio of the $`D^+/D^0`$ lifetime, indicating that the destructive Pauli interference present in the Cabibbo Favored $`D^+`$ decay is absent in the doubly Cabibbo suppressed (DCS) mode.
$`D^+K^{}K^+K^+`$ is an interesting DCS decay, which cannot even occur through a spectator diagram. FOCUS has the first observation of this mode, and reports a preliminary result for the Branching Ratio relative to $`K^{}\pi ^+\pi ^+`$ of ($`0.14\pm 0.02`$) $`\%`$. Several groups have reported observations of a $`D^+\varphi K^+`$ signal, however FOCUS did not find evidence for such decay brian-dcsd .
SELEX announced the first observation of a Cabibbo suppressed decay of a charm baryon through the decay $`\mathrm{\Xi }_c^+pK^{}\pi ^+`$ printSelex . Fig. 11 shows the signal of $`157\pm 22`$ events reported by SELEX and simple spectator diagrams with external $`W`$ emission for $`\mathrm{\Xi }_c^+`$ decaying into a Cabibbo allowed and into a single Cabibbo suppressed (SCS) mode. The other Cabibbo allowed $`\mathrm{\Xi }^{}`$ mode interchanges $`s`$ and $`d`$ quarks lines and produces a $`d\overline{d}`$ pair from the vacuum instead of a $`d\overline{u}`$ pair. FOCUS has also observed the same SCS decay, reporting a signal of $`86\pm 21`$ events from about 70 $`\%`$ of their data.
E791 has published plb423-98-185 results on the singly Cabibbo suppressed decay, $`D^0K^{}K^+\pi ^{}\pi ^+`$. A coherent amplitude analysis of the resonant substructure was used to extract decay fractions. Significant phase angles among different modes indicate very strong interference. The measured branching fractions relative to $`D^0K^{}\pi ^+\pi ^{}\pi ^+`$ are presented in the table 2.
### Hadronic charm decays, Dalitz plot Analysis
With the advent of high statistics experiments, charm meson decay have become a new way to study light meson spectroscopy. The amplitude analysis performed on Dalitz plots gives insight into the decay dynamics, providing direct information about intermediate resonances and relative decay fractions, and allowing to study final state interactions coming from the interference of the amplitudes describing competing resonant channels.
E791 has preliminary results on the decay of $`D^+`$ and $`D_s^+`$ mesons in three pions. A clear signal with $`1240\pm 51`$ $`D^+`$ and $`858\pm 49`$ $`D_s^+`$ was obtained after applying selection criteria aimed at identifying a clearly separated $`3\pi `$ vertex and after carefully estimating the backgrounds coming from possible reflections and three pion combinations. The branching ratios were normalized to $`D^+K^{}\pi ^+\pi ^+`$ ($`34,790\pm 232`$ events) and to $`D_s^+\varphi \pi ^+`$ ($`1038\pm 44`$ events) respectively. Efficiencies were obtained from a full Monte Carlo simulation. The branching ratio of the $`D^+\pi ^+\pi ^{}\pi ^+`$ relative to $`D^+K^{}\pi ^+\pi ^+`$ obtained was $`0.0329\pm 0.0015_{0.0026}^{+0.0016}`$. Similarly the branching ratio of the $`D_s^+\pi ^+\pi ^{}\pi ^{}`$ relative to $`D_s^+\varphi \pi ^+`$ was $`0.247\pm 0.028_{0.012}^{+0.019}`$.
#### $`D_s^+`$ Dalitz plot results from E791
Among the advantages of using charm meson decays to study light I=J=0 states is the fact that, unlike hadron-hadron scattering, in the decays of $`D`$ mesons the initial state is always $`J^P=0^{}`$, limiting the number of possible final states. The decay $`D_s^+\pi ^{}\pi ^+\pi ^+`$ is Cabibbo-favored without a strange meson in the final state. It can proceed via spectator amplitudes producing intermediate resonant states with hidden strangeness like the $`f_0(980)`$ or it can proceed via W-annihilation amplitudes producing intermediate resonant states with no strangeness. The decays like $`D_s^+\rho ^0\pi ^+`$ and the non-resonant $`D_s^+\pi ^{}\pi ^+\pi ^+`$ would proceed via $`W`$-annihilation mechanism. It would also be responsible for the decay $`D_s^+f_0(1370)\pi ^+`$, if the $`f_0(1370)`$ resonance is at least partially a $`u\overline{u}`$ \+ $`d\overline{d}`$ state as predicted by the simple quark model. The scale of the $`W`$-annihilation compared to the $`W`$-radiation amplitude would be indicated by the relative contribution of these channels to the $`\pi ^{}\pi ^+\pi ^+`$ final state albertoreis .
The Dalitz plot of $`D_s^+\pi ^{}\pi ^+\pi ^+`$ and the $`\pi \pi `$ mass projections Reis are shown in Fig. 12.
The $`f_0(980)\pi ^\pm `$ mode is the dominant one, accounting for nearly half of the $`D_s^+\pi ^{}\pi ^+\pi ^+`$ decay width. The $`f_0(980)\pi ^\pm `$ is often supposed to have a large $`s\overline{s}`$ component, indicating a large spectator amplitude in this decay. Significant contributions of $`f_0(1370)\pi ^+`$ and $`f_2(1270)\pi ^+`$ components were also found. The contribution of $`\rho ^0(770)\pi ^+`$ and $`\rho ^0(1450)\pi ^+`$ components correspond to about 10 $`\%`$ of the $`\pi ^{}\pi ^+\pi ^+`$ width. This could indicate either contribution from the annihilation diagram or from inelastic final state interactions. No significant non-resonant component was found.
Preliminary $`f_0`$ masses and widths Reis from E791 and PDG are presented in the table 3.
FOCUS also has preliminary results of this decay mode. Their preliminary Dalitz plot shows the $`D_s^+\pi ^{}\pi ^+\pi ^+`$ based on a very clean signal of $``$ 1300 events reconstructed from 100 $`\%`$ of their data. The preliminary results indicate a negligible contribution from the $`\rho `$ suggesting negligible Weak Annihilation contribution moroni .
#### $`D^+`$ Dalitz plot results and evidence for a light scalar resonance
The Dalitz plot of the single Cabibbo-suppressed decay $`D^+\pi ^{}\pi ^+\pi ^+`$ from E791 data is shown in Fig.13 (left). A coherent amplitude analysis was used to determine the structure of its density distribution. The fit including a non resonant amplitude and amplitudes for $`D^+`$ decaying to a $`\pi ^+`$ and any of the five established $`\pi ^+\pi ^{}`$ resonances $`\rho ^0(770)`$, $`f_0(980)`$, $`f_2(1270)`$, $`f_0(1370)`$, and $`\rho ^0(1450)`$ is shown in Fig. 13 (central). This fit is poor in the low $`\pi ^+\pi ^{}`$ region and has several unsatisfactory features jussara : the NR channel dominates, different from the $`D_s^+`$ decay, and the $`\rho ^0(1450)\pi ^+`$ is more significant than the $`\rho ^0(770)\pi ^+`$ state.
It was found that allowing an additional scalar state, with mass and width unconstrained improves the fit substantially. The mass of the resonance found by this fit is $`486_{26}^{+28}`$ MeV/c<sup>2</sup> and the width $`351_{43}^{+51}`$ MeV/c<sup>2</sup>. Referring to this $`\pi ^+\pi ^{}`$ resonance as the $`\sigma (500)`$, it was found that $`D^+\sigma (500)\pi ^+`$ accounts for about half of the total decay rate, non-resonant decay was very small and the $`\rho ^0(1450)\pi ^+`$ fraction was much less than $`\rho ^0(770)\pi ^+`$. Preliminary results of the fit with this state are shown in Fig. 13 (right side).
Theoretically, light scalar and isoscalar resonances are predicted in models for spontaneous breaking of chiral symmetry, like the $`\sigma `$ linear model sigma-model . These scalar particles have important consequences for the quark model, for understanding low energy $`\pi \pi `$ interactions and also for understanding the $`\mathrm{\Delta }I=1/2`$ rule.
#### Multidimensional analysis of $`\mathrm{\Lambda }_c^+pK^{}\pi ^+`$ from E791
E791 has reported recently the first amplitude analysis of the decay of a charm baryon e791-multia ; brian-dcsd . The study of charm baryon decays can give information regarding the relative importance of spectator and exchange amplitudes. Exchange amplitudes are small in charm meson decays because of helicity suppression. However in charm baryon decays this effect should not inhibit exchange amplitudes due to be three body nature of the interaction. The spectator and W-exchange diagrams can contribute to $`pK^0(890)`$, $`\mathrm{\Lambda }(1520)\pi ^+`$ or $`pK^{}\pi ^+`$ modes. However for the $`\mathrm{\Delta }^{++}(1232)K^{}`$ the $`W`$-exchange is the only diagram possible.
The charm baryon can be produced polarized and its decay products carry spin. These extra quantum numbers require five kinematic variables for a complete description of the decay, and instead of the conventional two dimensional Dalitz plot analysis, a five-dimensional amplitude analysis is required.
A sample of $`946\pm 38`$ $`\mathrm{\Lambda }_c^+pK^{}\pi ^+`$ reconstructed decays was used by E791 to determine relative strengths and phases of resonances in the final state as well as the $`\mathrm{\Lambda }_c`$ production polarization. The fit projections and the polarization as function of $`p_t`$ are shown in Fig.14. The resonant fractions for $`\mathrm{\Lambda }_c^+pK^{}\pi ^+`$ are shown in table 4.
The $`\mathrm{\Delta }^{++}(1232)K^{}`$ and the $`\mathrm{\Lambda }(1520)\pi ^+`$ decay modes are seen as statistically significant for the first time. The observation of a substantial $`\mathrm{\Delta }^{++}(1232)K^{}`$ component provides strong evidence for the W-exchange amplitude in charm baryon decays. It was also observed an increasingly negative polarization for the $`\mathrm{\Lambda }_c`$ as a function of $`p_t`$.
## Rare and forbidden decays
In the charm sector the rare and forbidden dilepton decay modes can be classified mainly into three categories (example of Feynman diagrams are shown in Fig. 15):
1) Flavor Changing Neutral Current decays (FCNC) such as $`D^0l^+l^{}`$ and $`D^+h^+l^+l^{}`$.
2) Lepton Family Number Violating decays (LFNV) such as $`D^+h^+l_1^+l_2^{}`$ and $`D^0l_1^+l_2^{}`$ where the leptons are from different generations.
3) Lepton number Violating decays (LNV) such as $`D^+h^{}l^+l^+`$ where the leptons are of the same generation but have the same sign.
Where $`h`$ stands for $`\pi `$, $`K`$ and $`l`$ for $`e,\mu `$.
The first decay modes (FCNC) are rare, that means a process suppressed via the GIM mechanism which proceeds via an internal quarks loop in the Standard Model scwartz . The FCNC decay mode $`D^0l^+l^{}`$ can proceed via a $`W`$ box diagram and the theoretical estimates Hewett for the branching fraction are of the order of $`10^{19}`$. The predictions for the other FCNC decay modes, $`D^+h^+l^+l^{}`$, are considerably larger, of the order of $`10^9`$ . These decay modes can proceed via penguin diagrams scwartz and from long distance effects.
The decays modes LFNV and LNV are strictly forbidden in the Standard Model as they do not conserve lepton number. However, some theoretical extensions of the Standard Model predict lepton number violation NVL , and then the observation of a signal in these modes would be evidence for new physics beyond the SM.
E791 has recently published e791-rares (see figure 15) a set of new limits in rare or forbidden decay modes that improve the PDG98 numbers by a factor of 10. They searched for 24 different rare and forbidden decay modes and have found no evidence for them. They therefore presented upper limits on their branching fractions. Fourteen of their limits represent a significant improvement over previous results and eight are presented for the first time.
For this study E791 used a blind analysis technique. The mass region where the signal is expected is masked throughout the analysis. Selection criteria are optimised by studying signal events generated by Monte Carlo simulation and background events obtained from data in mass windows above and below the signal region. The criteria were chosen to maximize the ratio $`N_S/\sqrt{(}N_B)`$, where $`N_S`$ and $`N_B`$ are the numbers of signal and background events, respectively. Only after this procedure were events within the signal window unmasked. This blind technique is used so that the presence or absence of signal does not bias the choice of the selection criteria.
FOCUS is also looking for rare and forbidden decays using the same technique. Some examples of rare decays are presented in table 5 where we compare the expected sensitivity from FOCUS to E791 results and PDG values.
## Semileptonic decays, Form-factors
The weak decays of hadrons containing heavy quarks are influenced by strong interaction effects. Semileptonic charm decays such as $`D^+\overline{K}^0e^+\nu _e`$, $`D_s^+\varphi l^+\nu _l`$ are an especially clean way to study these effects because the leptonic and hadronic currents completely factorize in the decay amplitude ,$`A`$, as we can see in the Eq. 4, where $`G_F`$ is the Fermi coupling constant for the weak interaction and $`V_{cs}`$ is the CKM matrix element. $`L^\mu `$ (Eq. 5) and $`H_\mu `$ (Eq. 6) represent the leptonic and hadronic currents key3 .
$$A(D^+\overline{K^{}}^0e^+\nu _e)=\frac{G_F}{\sqrt{2}}V_{cs}L^\mu H_\mu $$
(4)
$$L^\mu =\overline{u}_e\gamma ^\mu (1\gamma _5)v_\nu $$
(5)
$`H_\mu =(m_D+m_K^{})A_1(q^2)ϵ_{mu}{\displaystyle \frac{A_2(q^2)}{m_D+m_K^{}}}(ϵp_D)(p_D+p_K)_\mu `$
$`{\displaystyle \frac{A_3(q^2)}{m_D+m_K^{}}}(ϵp_D)(p_Dp_K)_\mu i{\displaystyle \frac{2V(q^2)}{m_D+m_K^{}}}\epsilon _{\mu \nu \rho \sigma }ϵ^\nu p_D^\rho p_K^\sigma `$ (6)
With a vector meson in the final state, there are four form factors, $`V(q^2)`$, $`A_1(q^2)`$, $`A_2(q^2)`$ and $`A_3(q^2)`$, which are functions of the Lorentz-invariant momentum transfer squared $`q^2`$, the square of the invariant mass of the virtual W key3 . The differential decay rate for $`D^+\overline{K}^0\mu ^+\nu _\mu `$ with $`\overline{K}^0K^{}\pi ^+`$ is a quadratic homogeneous function of the four form factors. Unfortunately, the limited size of current data samples does not allow precise measurement of the $`q^2`$-dependence of the form factors; we thus assume the dependence to be given by the nearest-pole dominance model: $`F(q^2)=F(0)/(1q^2/m_{pole}^2)`$ where $`m_{pole}=m_V=2.1\mathrm{GeV}/c^2`$ for the vector form factor $`V`$ (which correspond to $`J^P=\mathrm{\hspace{0.25em}1}^+`$ state, $`D_{s1}^{}`$), and $`m_{pole}=m_A=2.5\mathrm{GeV}/c^2`$ for the three axial- vector form factors A key4 (corresponding to $`1^{}`$ state, $`D_s^{}`$).
The third form factor $`A_3(q^2)`$, which is unobservable in the limit of vanishing lepton mass, probes the spin-0 component of the off-shell $`W`$. Additional spin-flip amplitudes, suppressed by an overall factor of $`m_{\mathrm{}}^2/q^2`$ when compared with spin no-flip amplitudes, contribute to the differential decay rate. Because $`A_1(q^2)`$ appears among the coefficients of every term in the differential decay rate, we can factor out $`A_1(0)`$ and measure the ratios:
$`r_V=V(0)/A_1(0)`$, $`r_2=A_2(0)/A_1(0)`$ and $`r_3=A_3(0)/A_1(0)`$. The values of these ratios can be extracted without any assumption about the total decay rate or the weak mixing matrix element $`V_{cs}`$.
We report E791 measurements of the form factor ratios $`r_v`$ and $`r_2`$ for the muon channel factor1-791 and combined with electron channel factor2-791 (see Fig. 16). This is the first set of measurements in both muon and electron channels from a single experiment. We also report the first measurement of $`r_3=A_3(0)/A_1(0)`$, which is unobservable in the limit of vanishing charged lepton mass.
The measurements of the form factor ratios for $`D^+\overline{K}^0\mu ^+\nu _\mu `$ presented here and for the similar decay channel $`D^+\overline{K}^0e^+\nu _e`$ key5 follow the same analysis procedure except for the charged lepton identification. Both results in the electron and muon channels are consistent within errors, supporting the assumption that strong interaction effects incorporated in the values of form factor ratios do not depend on the particular $`W^+`$ leptonic decay.
The combined results of electronic and muonic decay modes produce $`r_V=1.87\pm 0.08\pm 0.07`$ and $`r_2=0.73\pm 0.06\pm 0.08`$. The combination of all systematic errors is ultimately close to that which one would obtain assuming all the errors are uncorrelated. The third form factor ratio $`r_3=0.04\pm 0.33\pm 0.29`$ was not measured in the electronic mode.
Table 6 compares the values of the form factor ratios $`r_V`$ and $`r_2`$ measured by E791 in the electron, muon and combined modes with previous experimental results. The size of the data sample and the decay channel are listed for each case. All experimental results are consistent within errors. The comparison between the E791 combined values of the form factor ratios $`r_V`$ and $`r_2`$ and other experimental results is also shown in Fig. 16 together with theoretical predictions.
The FOCUS experiments anticipate measuring $`r_2`$ and $`r_V`$ to better precision than previous experiments BrianSL
## Charm Lifetimes
The study of the charm particle lifetimes is motivated by two main goals: to extract partial decay rates and to study decay dynamics. The total decay width can be expressed as a sum of the three possible classes of decays, so the lifetime of a particle can be written as
$$\tau =\frac{\mathrm{}}{\mathrm{\Gamma }_{tot}}=\frac{\mathrm{}}{\mathrm{\Gamma }_{leptonic}+\mathrm{\Gamma }_{SL}+\mathrm{\Gamma }_{nonleptonic}}$$
(7)
The leptonic partial width is normally very small ($`\mathrm{\Gamma }_{leptonic}10^310^4`$) due to helicity suppression. The observed $`\mathrm{\Gamma }_{SL}(D^+)`$ and $`\mathrm{\Gamma }_{SL}(D^0)`$ are equal to within $`10\%`$, as expected from isospin invariance. So, the large difference observed in the $`D^+`$ and $`D^0`$ lifetimes, $`\tau (D^+)/\tau (D^0)=2.55\pm 0.04`$, is due to a large difference in the hadronic decay rates ($`\mathrm{\Gamma }_{nonleptonic}`$) for the $`D^+`$ and the $`D^0`$. If there were no other diagram but the spectator and no QCD effects we would have the same value for $`\tau `$. Thus, the different lifetimes are an indication that we need to take into account contributions from diagrams where the W interact with two valence quarks, such as W-annihilation (WA) and W-exchange (WX), and any interference between them, as shows the Fig. 17.
A systematic approach now exists for the treatment of inclusive decays, based on QCD and consists of an Operator Product Expansion (OPE) in the Heavy Quark mass. In this approach the interaction is factorized into three parts: weak interaction between quarks, perturbative QCD corrections and non-perturbative QCD effects. The decay rate is given by
$$\mathrm{\Gamma }_{HQ}=\frac{G_F^2m_Q^2}{192\pi ^3}\mathrm{\Sigma }f_i|V_{Qq_i}|^2[A_1+\frac{A_2}{\mathrm{\Delta }^2}+\frac{A_3}{\mathrm{\Delta }^3}+\mathrm{}]$$
(8)
where $`\mathrm{\Delta }`$ is taken as the heavy quark mass and $`f_i`$ is a phase space factor. $`A_1=1`$ contains the spectator diagram contribution; $`A_2`$ is the spin interaction of the heavy quark with light quark degrees of freedom inside the hadron and produces differences between the baryon and meson lifetimes; $`A_3`$ includes the non-spectator $`W`$-annihilation, $`W`$-exchange and Pauli Interference (PI) of the decay and the spectator quarks contributions.
New results on charm lifetime measurements are shown in table 7. The most relevant information is the $`D_s^+`$ lifetime from E791 tau-Ds ; taud0-791 and FOCUS tau-Ds-focus which is now conclusively above the $`D^0`$ lifetime, with a ratio
$$R_\tau =\tau (D_s^+)/\tau (D^0)=1.22\pm 0.02$$
(9)
It is important to note that $`R_\tau `$ is now ten standard deviations away from unity, indicating that although not dominant, the WA/WX contributions are significant. In fact the OPE model predicts $`R_\tau `$ = 1.00 - 1.07 without WA/WX contributions, allowing a variation of $`\pm 20\%`$ if the WA/WX operators are included, in agreement with the new result.
In the baryon sector SELEX and FOCUS have preliminary measurements of the $`\mathrm{\Lambda }_c^+`$ lifetime, as shown in table.7 and Fig. 18. Using 100$`\%`$ of their data in $`\mathrm{\Lambda }_c^+pK^{}\pi ^+`$, SELEX find $`\tau =177\pm 10(stat)fs`$ tau-lambda-selex , in a $`23\sigma `$ disagreement with PDG98 and FOCUS.
### Lifetime differences and $`D^0\overline{D^0}`$ mixing
E791 has published searches for a lifetime difference between the $`CPeven`$ and $`CPodd`$ eigenstates of the $`D^0`$ taud0-791 . To do so they compared the lifetimes of the decays $`D^0K^{}K^+`$ ($`CP=+1`$) and $`D^0K^{}\pi ^+`$, (CP mixed) cpvs , shown in Fig. 19.
Defining
$$\frac{\mathrm{\Gamma }(K^{}K^+)\mathrm{\Gamma }(K^{}\pi ^+)}{\mathrm{\Gamma }(K^{}\pi ^+)}=y_{CP}$$
(10)
The time integrated ratio of mixed to non mixed decay rates in charm meson is given by
$$R_{mix}=\frac{\mathrm{\Gamma }(D^0\overline{D^0}\overline{f})}{D^0f}=\frac{x^2+y^2}{2}$$
(11)
where
$`x={\displaystyle \frac{\mathrm{\Delta }m}{\overline{\mathrm{\Gamma }}}},y={\displaystyle \frac{\mathrm{\Delta }\mathrm{\Gamma }}{2\overline{\mathrm{\Gamma }}}}`$
with
$`\mathrm{\Delta }m=m_1m_2,\mathrm{\Delta }\mathrm{\Gamma }=\mathrm{\Gamma }_1\mathrm{\Gamma }_2,\overline{\mathrm{\Gamma }}=(\mathrm{\Gamma }_1+\mathrm{\Gamma }_2)/2`$
where $`\mathrm{\Gamma }_1`$ is for CP even states and $`\mathrm{\Gamma }_2`$ for CP-odd states.
$`\mathrm{\Gamma }_1`$ applies to $`D^0K^{}K^+`$ and $`\mathrm{\Gamma }`$ applies to $`D^0K^{}\pi ^+`$ if CP is conserved.
Mixing can appear if there is either a difference in the masses of the CP eigenstates $`\mathrm{\Delta }m`$ or if there is a difference in the decay rates $`\mathrm{\Delta }\mathrm{\Gamma }`$.
E791 observed no difference in lifetimes and quoted $`\tau (K\pi )`$ = $`0.413\pm 0.003\pm 0.004`$ ps $`\tau (KK)`$ = $`0.410\pm 0.011\pm 0.006`$ ps, $`y_{CP}`$ = $`0.008\pm 0.029\pm 0.010`$ or $`0.04<y_{CP}<0.06`$ ( $`90\%`$ CL) or $`\mathrm{\Delta }\mathrm{\Gamma }`$= $`2(\mathrm{\Gamma }_{KK}\mathrm{\Gamma }_{K\pi })`$ = $`(0.04\pm 0.14\pm 0.05)ps^1`$
## ACKNOWLEDGEMENTS
One of us, J. C. Anjos would like to thank the conference organizers for their invitation to attend the Symposium. E. Cuautle would like to thank the Centro Brasileiro de Pesquisas Físicas (CBPF) for its kind hospitality during his postdoctoral stay. The authors would like to thank CLAF/CNPq (Brazil) and CONACyT (México) for financial support of this work.
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# Covariant two point function for minimally coupled scalar field in de Sitter space-time
## 1 Introduction
The de Sitter metric appears as a logical geometrical background for inflationary cosmological scenario. We in particular insist on the fact that recent observational data are strongly in favor of a positive acceleration of the present universe. In a first approximation, this acceleration may be described by a constant $`\mathrm{\Lambda }`$-term in the Einstein’s equation of gravitation, even if this $`\mathrm{\Lambda }`$-term may be permitted to slowly vary with time . Therefore, the background space-time is very similar to a de Sitter space-time, and a quantization of the linear gravitational field without infrared divergence in such a space-time model may reveal to be of extreme importance for further developments. In a previous paper we have shown that one can construct a covariant quantization of the “massless” minimally coupled scalar field in de Sitter space-time, which is free of any infrared divergence. In a forthcoming paper , we prove that this is true for linear gravity (the traceless rank-2 “massless” tensor field), which has been also demonstrated by . Antoniadis, Iliopoulos and Tomaras have also shown that the pathological large-distance behavior of the graviton propagator on a dS background does not manifest itself in the quadratic part of the effective action in the one-loop approximation. This means that the pathological behavior of the graviton propagator may be gauge dependent and so should not appear in an effective way as a physical quantity. Recently, de Vega and al. have shown that, in flat coordinates on de Sitter space-time, the infrared divergence does not appear in the “massless” minimally coupled scalar field. They use this result for studying the generation of the gravitational waves in de Sitter space-time. They have calculated the retarded Green’s function in terms of flat coordinates of de Sitter space.
This paper is a completion of our previous work . We make explicit here the two-point function for the “massless” minimally coupled scalar field in de Sitter space-time in a fully gauge invariant way (coordinate independent), which is free of any infrared divergence. In Section $`2`$ we briefly recall the notations used in this paper. In Section $`3`$, we give an explicit construction of the covariant two-point function and the related Schwinger commutator function. Then the retarded Green’s function is calculated. In section $`4`$, we compare this result with the results obtained by de Vega and al., and we make precise in what way the latter are de Sitter invariant. Finally, section $`5`$ contains a brief conclusion and outlook.
## 2 Notation
The de Sitter space is an elementary solution of the cosmological Einstein equation. It is conveniently seen as a hyperboloid embedded in a five-dimensional Minkowski space
$$X_H=\{x\mathrm{IR}^5;x^2=\eta _{\alpha \beta }x^\alpha x^\beta =H^2\},\alpha ,\beta =0,1,2,3,4,$$
(1)
where $`\eta _{\alpha \beta }=`$diag$`(1,1,1,1,1)`$. The de Sitter metrics reads
$$ds^2=\eta _{\alpha \beta }dx^\alpha dx^\beta =g_{\mu \nu }^{dS}dX^\mu dX^\nu ,\mu =0,1,2,3,$$
where the $`X^\mu `$’s are the $`4`$ space-time coordinates in dS hyperboloid. Different coordinate systems can be chosen . A so-called flat coordinatization satisfying $`(1)`$ and covering the half of the de Sitter hyperboloid is given by
$$\{\begin{array}{cccc}x^0& =H^1\mathrm{sinh}Ht+\frac{1}{2}He^{Ht}\stackrel{}{X}^2\hfill & & \\ x^i& =e^{Ht}X^i,i=1,2,3\hfill & & \\ x^4& =H^1\mathrm{cosh}Ht\frac{1}{2}He^{Ht}\stackrel{}{X}^2.\hfill & & \end{array}$$
(2)
Let us give details on this (non-global) coordinate system in order to compare our results with those of . The line element is given by
$$ds^2=dt^2e^{2Ht}(dX^1dX^1+dX^2dX^2+dX^3dX^3)$$
$$=dt^2e^{2Ht}[d\rho ^2+\rho ^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)],$$
(3)
where $`0\rho <\mathrm{}`$. In the conformal time notation, $`\eta =H^1e^{Ht},\mathrm{}\eta 0`$, and so we have
$$ds^2=\frac{1}{H^2\eta ^2}[d\eta ^2d\rho ^2\rho ^2(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2)].$$
(4)
Given two points on the de Sitter hyperboloid $`x`$ and $`x^{}`$, the quadratic form
$$𝒵(x,x^{})=H^2x.x^{}=H^2\eta _{\alpha \beta }x^\alpha x^\beta ,$$
(5)
is invariant under the isometry group $`O(4,1)`$. Moreover, any invariant function is function of this basic invariant $`𝒵(x,x^{})`$. The latter can be written in terms of the invariant geodesic distance $`\sigma `$ between the points $`x`$ and $`x^{}`$. In the flat coordinate system $`(2)`$, we have
$$𝒵(x,x^{})=H^2x.x^{}=1+\frac{H^2}{2}(xx^{})^2=1+\frac{1}{2\eta \eta ^{}}[(\eta \eta ^{})^2(\stackrel{}{X}\stackrel{}{X}^{})^2]=\mathrm{cosh}H\sigma .$$
We now briefly recall the quantization of the minimally coupled massless scalar field. The latter is defined by
$$\mathrm{}_H\varphi (x)=0,$$
where $`\mathrm{}_H`$ is the Laplace-Beltrami operator on de Sitter space. As proved by Allen , the covariant canonical quantization procedure with positive norm states must fail in this case. The Allen’s result can be reformulated in the following way: the Hilbert space generated by any complete set of modes is not de Sitter invariant,
$$=\{\underset{k}{}\alpha _k\varphi _k;\underset{k}{}|\alpha _k|^2<\mathrm{}\},$$
where $`k=\{(L,l,m)\text{IN}\times \text{IN}\times \text{ZZ};\mathrm{\hspace{0.33em}0}lL,lml\}`$. This means that it is not closed under the action of the de Sitter group. Nevertheless, one can obtain a fully covariant quantum field by adopting a new construction . In order to obtain a fully covariant quantum field, we added all the modes $`L<1`$. In consequence, we have to deal with an orthogonal sum of a positive definite inner product space and a negative one. This sum is closed under the de Sitter group. The negative values of the inner product are precisely produced by the conjugate modes: $`\varphi _k^{},\varphi _k^{}=1`$. We do insist on the fact that the space of states contains unphysical states, which have a negative norm. Now, we can decompose the field into positive and negative norm parts (eq. $`(30`$) ref. )
$$\varphi (x)=\frac{1}{\sqrt{2}}\left[\varphi _p(x)\varphi _n(x)\right],$$
(6)
where
$$\varphi _p(x)=\underset{k}{}a_k\varphi _k(x)+H.C.,\varphi _n(x)=\underset{k}{}b_k\varphi ^{}(x)+H.C..$$
(7)
The positive mode $`\varphi _p(x)`$ is the scalar field as was used by Allen. Then we have the following operatorial relations
$$a_k|0>=0,[a_k,a_k^{}^{}]=\delta _{kk^{}},b_k|0>=0,[b_k,b_k^{}^{}]=\delta _{kk^{}}.$$
(8)
## 3 Two-point function
The explicit knowledge of Wightman two-point function $`𝒲`$ allows us to make the quantum field formalism work. It also allows defining the various Green’s functions of the wave equation. This function is defined by
$$𝒲(x,x^{})=<0\varphi (x)\varphi (x^{})0>.$$
For the scalar field $`(6)`$ we have
$$𝒲(x,x^{})=\frac{1}{2}[<0\varphi _p(x)\varphi _p(x^{})0>+<0\varphi _n(x)\varphi _n(x^{})0>]$$
$$=\frac{1}{2}[𝒲_p(x,x^{})+𝒲_n(x,x^{})],$$
(9)
where $`𝒲_n(x,x^{})=𝒲_p^{}(x,x^{})`$ and $`𝒲_p(x,x^{})`$ is the two-point function for the positive modes as it was calculated by Allen and Folacci . Its expression reads
$$𝒲_p(x,x^{})=\frac{H^2}{8\pi ^2}[\frac{1}{1𝒵}\mathrm{ln}(1𝒵)+\mathrm{ln}2+f_{AB}(\eta ,\eta ^{})],$$
(10)
where $`f_{AB}(\eta ,\eta ^{})`$ is a function of the conformal time and is not de Sitter invariant. Therefore the two-point function $`(9)`$ is de Sitter invariant and reads:
$$𝒲(x,x^{})=\frac{iH^2}{8\pi }ϵ(x^0x^0)[\delta (1𝒵(x,x^{}))\theta (𝒵(x,x^{})1)],$$
(11)
where
$$ϵ(x^0x^0)=\{\begin{array}{cccc}1& x^0>x^0\hfill & & \\ 0& x^0=x^0\hfill & & \\ 1& x^0<x^0.\hfill & & \end{array},$$
(12)
and $`\theta `$ is the Heaviside step function. This expression in contrast to the previous calculation is de Sitter invariant, i.e. coordinate independent and also free of any infrared divergence. Now we can write the various Green functions from $`(11)`$. The anticommutation function is given by
$$G^1(x,x^{})=\{\varphi (x),\varphi (x^{})\}=2e𝒲(x,x^{})=0.$$
(13)
For the Schwinger commutator function, we have
$$G(x,x^{})=i[\varphi (x),\varphi (x^{})]=2i𝒲(x,x^{})$$
$$=\frac{H^2}{4\pi }ϵ(x^0x^0)[\delta (1𝒵(x,x^{}))\theta (𝒵(x,x^{})1)].$$
(14)
Then corresponding retarded Green’s function is given by
$$G_R(x,x^{})=\theta (x^0x^0)G(x,x^{})$$
$$=\frac{H^2}{4\pi }\theta (x^0x^0)ϵ(x^0x^0)[\delta (1𝒵(x,x^{}))\theta (𝒵(x,x^{})1)].$$
(15)
Thus, the retarded Green’s function propagates on the past light cone and in its interior. It is also de Sitter invariant and free of any infrared divergence.
## 4 Comparison to previous works
The Schwinger commutator function $`(14)`$ is de Sitter invariant and so coordinate independent. In this section, we show that it is the same as the Schwinger commutator function which was calculated in the flat coordinate system $`(2)`$ by de Vega and al. . The latter used the corresponding flat modes solutions. The field operator was written in terms of the Fourier transformation in the flat space
$$\varphi (\eta ,\stackrel{}{X})=d^3\stackrel{}{k}[u_\stackrel{}{k}(\eta ,\stackrel{}{X})a(\stackrel{}{k})+u_\stackrel{}{k}^{}(\eta ,\stackrel{}{X})a^{}(\stackrel{}{k})],$$
(16)
where
$$u_\stackrel{}{k}(\eta ,\stackrel{}{X})=\frac{H}{(2\pi )^{3/2}\sqrt{2w}}(\eta \frac{i}{w})\mathrm{exp}(iw\eta +i\stackrel{}{k}\stackrel{}{X}),$$
(17)
and $`w=\stackrel{}{k}`$. There results the Schwinger commutator function
$$G(x,x^{})=\frac{H^2}{4\pi \stackrel{}{X}\stackrel{}{X}^{}}\times $$
$$\left[\eta \eta ^{}(\delta (y_{})\delta (y_+))+(\eta \eta ^{})(\theta (y_{})\theta (y_+))+\frac{1}{2}(y_+y_{})\right],$$
(18)
where
$$y_\pm =(\eta \eta ^{})\pm \stackrel{}{X}\stackrel{}{X}^{}.$$
For calculating this function the authors have not used the two point function, which is divergence in their construction. They calculated directly the commutator function in which the infrared divergence disappears due to the sign of the divergence term. If we use the following relations,
$$\theta (t^2x^2)=\theta (tx)\theta (t+x)+1,$$
$$ϵ(t)\theta (t^2x^2)=\frac{1}{x}\{t[\theta (tx)\theta (t+x)]+\frac{1}{2}(t+xtx)\},$$
$$ϵ(\eta \eta ^{})\theta (\frac{1}{2\eta \eta ^{}}y_+y_{})=\frac{\eta \eta ^{}}{\stackrel{}{X}\stackrel{}{X}^{}}[\theta (y_{})\theta (y_+)]+\frac{y_+y_{}}{2\stackrel{}{X}\stackrel{}{X}^{}},$$
(19)
$$\delta (1𝒵)=\delta (\frac{1}{2\eta \eta ^{}}y_+y_{})=2\eta \eta ^{}\delta (y_+y_{})=\frac{\eta \eta ^{}}{\stackrel{}{X}\stackrel{}{X}^{}}(\delta (y_{})+\delta (y_+)),$$
(20)
it is easy to show that Eq. $`(18)`$ is the same as Eq. $`(14)`$, which was calculated from the finite two point function $`(11)`$. The result of de Vega and al., which was claimed by them as free of infrared divergence, can be considered as a direct consequence of a requirement of de Sitter fully invariance. Note that the latter is irremediably broken if commutations are worked out with non-global coordinates.
## 5 Conclusion
In this paper we have presented a covariant two-point function for minimally coupled scalar field in de Sitter space-time, which is free of any infrared divergence and is also de Sitter invariant (coordinate independent). Using this result in a forthcoming paper , we are able to get a covariant two-point function for linear gravity (the traceless rank-2 “massless” tensor field) in de Sitter space-time, which is free of infrared divergence. This means that pathological behavior may be gauge dependent. It should not appear in an effective way in a physical quantity whereas it exists in an irreducible way in the pure-trace part (conformal sector) . So far the physical meanings of these results have not been clarified. There are different possibilities. Let us mention two of them. The negative norm states are totally part of the structure of the renormalized field and it is not needed to remove them. Then we have de Sitter invariance for linear gravity (the traceless rank-2 “massless” tensor field). But in the inflationary model one introduces an inflaton scalar field. Because of this field, the conformal sector (pure trace) of the metric becomes dynamical and it must be quantized . This sector produces a gravitational instability and the breaking of de Sitter invariance. In this point of view, one may use a scalar-tensor theory for gravity $`(g_{\mu \nu },\varphi )`$ and this paves the way to the existence of an universe with a positive acceleration . The other week possibility is the following. Like for covariant QED à la Gupta-Bleuler in Minkowski space negative norm states appear in a covariant quantization of the minimally coupled scalar field in de Sitter space. They are eliminated under the constraint of primordial fluctuations in the inflationary model. The latter indeed selects the physical states (the positive norm states) . The price to pay is precisely the breaking of de Sitter invariance. One can say that the negative norm states play the role ghost states. Now, breaking the de Sitter invariance is responsible of infrared divergence and this produces a gravitational instability. This gravitational instability and the primordial quantum fluctuations of the inflaton scalar field define the inflationary model. The latter can explain the formation of the galaxies, clusters of galaxies and the large scale structure of the universe .
Acknowlegements: The author would like to thank J. P. Gazeau, J. Iliopoulos, D. Polarski and J. Renaud for very useful discussions.
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# Magnetoelastic coupling in epitaxial magnetic films: An ab-initio study
## Abstract
A method is developed which allows to determine the first-order and the second-order magnetoelastic coefficients of a magnetic bulk material from the ab-initio calculation of the magnetocrystalline anisotropy energy as function of a prestrain $`ϵ_0`$. Explicit results are given for bcc Fe, and they agree well with experimental data obtained from the magnetostrictive stress measurements for epitaxial Fe films.
PACS: 75.70.-i, 75.80.+q, 71.15.-m
In recent years magnetic devices based on magnetic films technologies have attracted a considerable interest, e.g., magnetooptical recording media or magnetoresistive devices based on the giant magnetoresistive and the tunnel magnetoresistive effect designed for sensors or magnetostrictive random access memories. Thereby the magnetic anisotropy plays an important role, for instance, the issue of perpendicular anisotropy for the magnetooptical recording or the demand for soft magnetic layers with weak anisotropy as part of the magnetoresistive devices. It has been shown by numerous investigations that the magnetic anisotropies of magnetic films grown epitaxially on a substrate may strongly deviate from those of the respective bulk materials. The reason for this deviation is in general ascribed to several different effects. First, there are contributions to the anisotropy originating from the free surface of the magnetic layer and from the interface between layer and substrate, as well as from the morphology of the film due to a heterogeneous film growth. The influence of all these effects must decrease with increasing thickness of the film. What remains for a film of thickness larger than typically $`10\mathrm{nm}`$ is the effect of the magnetoelastic coupling to the film strain induced by the lattice mismatch between film and substrate. Because the epitaxial film strain may be of the order of several % which is much larger than typical magnetostrictive strains of $`10^6`$ to $`10^4`$ and because of the dependence of the magnetoelastic coupling energy on the direction of the magnetization (see below), this may result in new magnetic anisotropies different from that of the unstrained bulk. The numerous experiments on the effect of epitaxial strain on the magnetic properties of magnetic films are reviewed in Ref. . An ab-initio study of this effect within the framework of density functional theory is the purpose of the present letter.
To be more specific we consider a material which is cubic in the unstrained state (Fe, for instance). Then the density of the magnetoelastic coupling energy may be written (up to the second order in the strain $`ϵ_{ij}`$, omitting the terms including the shear strain $`ϵ_{ij},ij`$, which are not required for the situation discussed below) as
$`f_{\mathrm{me}}=B_1\left(ϵ_{11}\alpha _1^2+ϵ_{22}\alpha _2^2+ϵ_{33}\alpha _3^2\right)+`$ (1)
$`{\displaystyle \frac{1}{2}}D_{11}\left(ϵ_{11}^2\alpha _1^4+ϵ_{22}^2\alpha _2^4+ϵ_{33}^2\alpha _3^4\right)+D_{12}\left(ϵ_{11}ϵ_{22}\alpha _1^2\alpha _2^2+ϵ_{22}ϵ_{33}\alpha _2^2\alpha _3^3+ϵ_{33}ϵ_{11}\alpha _3^2\alpha _1^2\right).`$ (2)
Here $`B_1`$ and $`D_{11}`$, $`D_{12}`$ represent magnetoelastic coupling coefficients of the first and the second order, and the $`\alpha _i`$ denote the direction cosines of the magnetization referred to the cubic axes. For Fe grown epitaxially on a cubic (100) surface which is a prototype system experimentally investigated intensively there are epitaxial strains $`ϵ_{11}=ϵ_{22}`$ which may be different for various atomic layers of the film . In the following we adopt a simple model where we consider only the average strain of the film which depends on the thickness of the film, i.e., we assume $`ϵ_{11}=ϵ_{22}=ϵ_0`$. The strain $`ϵ_{33}`$ then may be obtained by minimizing the total energy density $`f=f_{\mathrm{me}}+f_{\mathrm{el}}`$, with the elastic energy density (again omitting the terms containing the shear strains):
$$f_{\mathrm{el}}=\frac{1}{2}C_{11}\left(ϵ_{11}^2+ϵ_{22}^2+ϵ_{33}^2\right)+C_{12}\left(ϵ_{11}ϵ_{22}+ϵ_{22}ϵ_{33}+ϵ_{33}ϵ_{11}\right)$$
(3)
and with the cubic elastic stiffness constants $`C_{ij}`$. This yields
$$ϵ_{33}=\frac{2C_{12}ϵ_0+B_1\alpha _3^2}{C_{11}+D_{11}\alpha _3^4}.$$
(4)
Because the magnetostrictive contribution to the strain originating from $`f_{\mathrm{me}}`$ is much smaller than $`ϵ_0`$ we find $`ϵ_{33}2C_{12}/C_{11}`$.
From (1), (3) we obtain the strain dependent part, $`f_{\mathrm{mca}}(ϵ_0)`$, of the magnetocrystalline anisotropy energy density:
$$f_{\mathrm{mca}}(ϵ_0)=f_{\mathrm{me}}(ϵ_0,\alpha _1=1)f_{\mathrm{me}}(ϵ_0,\alpha _3=1)=k_0+k_1ϵ_0+k_2ϵ_0^2,$$
(5)
with
$`k_0=`$ $`{\displaystyle \frac{1}{2}}B_1^2{\displaystyle \frac{1}{C_{11}+D_{11}}},`$ (6)
$`k_1=`$ $`B_1\left(1+2C_{12}{\displaystyle \frac{1}{C_{11}+D_{11}}}\right),`$ (7)
$`k_2=`$ $`{\displaystyle \frac{1}{2}}D_{11}\left(14{\displaystyle \frac{C_{12}^2}{C_{11}}}{\displaystyle \frac{1}{C_{11}+D_{11}}}\right).`$ (8)
The term $`k_0`$ is negligible and arises from the fact that we fix the strains $`ϵ_{11}=ϵ_{22}=ϵ_0`$ independently on the direction of the magnetization. Furthermore, when changing the direction of the magnetization there is a change in the magnetostrictive stress $`\tau _1=\frac{f_{\mathrm{me}}}{ϵ_{11}}`$ according to
$$\mathrm{\Delta }\tau _1=\tau _1(\alpha _1=1)\tau _1(\alpha _2=1)=B_1+D_{11}ϵ_0.$$
(9)
An experimental determination of $`\mathrm{\Delta }\tau _1`$ (exploiting the change of the bending moment that is created by the film onto the substrate ) as function of the layer thickness and hence as function of $`ϵ_0`$ then yields the two magnetoelastic coupling coefficients $`B_1`$ and $`D_{11}`$. For Fe on MgO (100) Koch et al. obtained $`B_1=3.2\mathrm{MJ}/\mathrm{m}^3`$, $`D_{11}=1.1\mathrm{GJ}/\mathrm{m}^3`$ $`(\pm 10\%)`$, and for Fe on W(100) Enders et al. found $`B_1=3\mathrm{MJ}/\mathrm{m}^3`$, $`D_{11}=1\mathrm{GJ}/\mathrm{m}^3`$. The values for $`B_1`$ extracted from the film experiments agree rather well with the bulk value of $`B_1=3.44\mathrm{MJ}/\mathrm{m}^3`$. There are no values of $`D_{11}`$ obtained from bulk measurements for comparison and therefore the two experiments were considered as the first determination of the second-order magnetoelastic coupling coefficient of bulk Fe by a film experiment.
It was pointed out that due to the large strains accessible by epitaxial film growth the effective first-order coefficient $`B_{\mathrm{eff}}`$ defined as $`B_{\mathrm{eff}}=B_1+D_{11}ϵ_0`$ changes sign from negative to positive for Fe with decreasing film thickness, i.e., increasing $`ϵ_0`$, and this clearly demonstrates that the magnetic anisotropy energy depends dramatically on the film thickness, a result which is most relevant for the design of the magnetic film devices (see introduction). The linear strain dependence of $`B_{\mathrm{eff}}`$ failed to describe the experimental data for film thicknesses below $`10\mathrm{nm}`$, most probably because then the effects of the surface, interface and film morphology become relevant.
In the present paper we determine the magnetoelastic coupling coefficients $`B_1`$ and $`D_{11}`$ for cubic Fe by the ab-initio density functional theory. To do this, we calculate the magnetocrystalline anisotropy energy density $`f_{\mathrm{mca}}(ϵ_0)`$ as function of the strain $`ϵ_{11}=ϵ_{22}=ϵ_0`$ imposed to the bulk material, represent the data by a quadratic polynomial in $`ϵ_0`$ according to eq. (4) and determine $`B_1`$ and $`D_{11}`$ from eqs. (6,7), inserting the elastic stiffness constants $`C_{12}`$ and $`C_{11}`$ which we have also obtained ab initio.
We have performed the calculations using the WIEN97 code which adopts the full-potential linearized augmented plane-wave (FLAPW) method . For the exchange-correlation potential the local-spin-density (LSDA) functional by Perdew and Wang and the generalized-gradient-approximation (GGA) functional by Perdew et al. were used. The total energy minimizations on the non-strained bcc Fe gave us the equilibrium lattice parameters $`a=5.2\mathrm{a}_0`$ for LSDA and $`a=5.34\mathrm{a}_0`$ for GGA where $`\mathrm{a}_0`$ denotes Bohr’s radius. The calculated ratio $`2C_{12}/C_{11}`$ is 1.08 for LSDA and 1.13 for GGA. The experimental values are $`a=5.42\mathrm{a}_0`$ and $`2C_{12}/C_{11}=1.17`$.
Numerically the most difficult step is the calculation of $`f_{\mathrm{mca}}`$ which is due to the spin-orbit coupling (SOC). First, we calculate the self-consistent electronic structure in the scalar-relativistic approximation using $`N_𝐤^3`$ $`𝐤`$ vectors with $`N_𝐤=21`$ in the total Brillouin zone (BZ) which correspond to the 762 $`𝐤`$ vectors in the irreducible part of the Brillouin zone (IBZ). The criterion for the self-consistency is the difference in the charge densities after the last two iterations being less than $`2\times 10^6e/(\mathrm{a}.\mathrm{u}.)^3`$. The contribution of the SOC is determined perturbatively using the second variational method . The quantity $`f_{\mathrm{mca}}`$ is calculated by applying the force theorem as the difference between the sums of the perturbed eigenvalues for the different magnetization directions. Fig. 1 represents the convergency test for the calculation of $`f_{\mathrm{mca}}`$ with respect to $`N_𝐤`$. The data are for the case with $`a=5.4\mathrm{a}_0`$ and $`c=5.2\mathrm{a}_0`$, using GGA. The modified tetrahedron and the Gaussian smearing integration schemes were used. The proper convergency with the Gaussian smearing was achieved by setting the smearing parameter as $`\mathrm{\Gamma }/N_𝐤`$. The suitable values of $`\mathrm{\Gamma }`$ for the particular case are roughly from the interval between $`6.8\mathrm{eV}`$ and $`10.1\mathrm{eV}`$ which follows from the curves in Fig. 1. All the final calculations of $`f_{\mathrm{mca}}`$ were performed with $`N_𝐤=51`$ (17576 $`𝐤`$ vectors in the IBZ) using both the tetrahedron and the Gaussian smearing method in order to minimize the numerical uncertainties. The estimated accuracy is $`\pm 1\mu \mathrm{eV}/\mathrm{unit}\mathrm{cell}`$ marked by the horizontal lines in Fig. 1.
Fig. 2 shows the calculated magnetocrystalline anisotropy energy density $`f_{\mathrm{mca}}`$ with respect to the lateral strain $`ϵ_0`$. The numerical data are well fitted by quadratic polynomials as predicted by eq. (4). From the calculated parameters $`k_1`$ and $`k_2`$ the magnetoelastic coefficients $`B_1`$ and $`D_{11}`$ are determined according to eqs. (6,7). Table I summarizes the theoretical results in comparison with the experimental data from . There is a big discrepancy between the LSDA result for $`B_1`$ and the experimental result, whereas the GGA result is much closer to the experiment. This is in line with the calculation of the magnetoelastic coefficient $`\lambda _{100}`$ of unstrained bulk Fe by Wu et al. who also obtained a strong deviation from the experiment when using LSDA but a satisfactory agreement when using GGA. The calculated second-order magnetoelastic coupling coefficient $`D_{11}`$ for bulk bcc Fe matches the experimental value obtained from the measurements on epitaxial thin films very well, especially the value from the GGA calculation. The agreement represents the direct proof that the experimental results can be really ascribed to the pure strain effect on the magnetoelastic properties and that the measurements of the magnetostrictive film stress as function of the film thickness can provide the second-order coupling constant $`D_{11}`$ of the bulk which is hard to obtain by bulk measurements.
We close with an important warning. The experiments and the present theory demonstrate that the magnetoelastic properties of thin epitaxial films may deviate significantly from that described by the first-order magnetoelastic coupling coefficients of the bulk and that one has to take into account the second-order terms of the bulk. The change of the magnetostrictive stress $`\mathrm{\Delta }\tau _1`$ obtained when switching the magnetization from to as function of the epitaxial strain $`ϵ_0`$ then may be expressed by an effective first-order coefficient $`B_{\mathrm{eff}}=B_1+D_{11}ϵ_0`$. However, the correct result for the coefficients $`k_0`$, $`k_1`$ and $`k_2`$ of the polynomial expansion (4) of $`f_{\mathrm{mca}}`$ may not be obtained by neglecting the second-order terms in eq. (1) and instead replacing $`B_1`$ in the first-order term by $`B_{\mathrm{eff}}=B_1+D_{11}ϵ_0`$. This would yield $`f_{\mathrm{mca}}=k_0+k_1ϵ_0+k_2ϵ_0^2`$ with $`k_1=B_1(1+2C_{12}/C_{11})`$ and $`k_2=D(1+2C_{12}/C_{11})`$ instead of eqs. (6,7). Inserting the value of $`2C_{12}/C_{11}`$ for Fe it becomes obvious that $`D`$ and $`D_{11}`$ even have a different sign; the sign of $`D`$ being opposite to the one of the experimentally determined $`D_{11}`$ according to eq. (8). Guo et al. have performed an ab-initio calculation of the magnetoelastic properties of epitaxial Co and Ni films, and they indeed proceeded on this line, i.e. they neglected the second-order term in eq. (1) and instead replaced the constant $`B_1`$ by a strain-dependent term $`B_1(ϵ_0)`$. It should be cautioned that the coefficient $`D`$ obtained from the linearization of $`B_1(ϵ_0)`$ is not identical to the second order coefficient $`D_{11}`$ of the bulk material but it may deviate strongly.
Acknowledgement: The authors are indebted to O. Grotheer, P. Novak and R.Q. Wu for helpful discussions.
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# Non-abelian statistics of half-quantum vortices in 𝑝-wave superconductors
## Abstract
Excitation spectrum of a half-quantum vortex in a $`p`$-wave superconductor contains a zero-energy Majorana fermion. This results in a degeneracy of the ground state of the system of several vortices. From the properties of the solutions to Bogoliubov-de-Gennes equations in the vortex core we derive the non-abelian statistics of vortices identical to that for the Moore-Read (Pfaffian) quantum Hall state.
Certain types of superconductors with triplet pairing allow half-quantum vortices . Such vortices appear if the multi-component order parameter has extra degrees of freedom besides the overall phase, and the vortex involves both a rotation of the phase by $`\pi `$ and a rotation of the “direction” of the order parameter by $`\pi `$, so that the order parameter maps to itself on going around the vortex. The magnetic flux through such a vortex is one half of the superconducting flux quantum $`\mathrm{\Phi }_0`$.
Another signature of this unusual flux quantization is a Majorana fermion level at zero energy inside the vortex core . This energy level has a topological nature and from the continuity considerations must be stable to any local perturbations. In terms of the energy levels, the Majorana fermions in vortex cores imply a $`2^n`$-fold degeneracy of the ground state of a system with $`2n`$ isolated vortices. If we let vortices adiabatically move around each other, this motion may result in a unitary transformation in the space of ground states (non-abelian statistics). We shall see that it is indeed the case.
The non-abelian statistics for half-quantum vortices has been originally derived for the Pfaffian quantum Hall state proposed by Moore and Read . The Pfaffian state is of Laughlin type and may be possibly realized for filling fractions with even denominator. The excitations in the Pfaffian state are half-quantum vortices, and their non-abelian statistics has been derived in the field-theoretical framework .
On the other hand, recently Read and Green suggested that the Pfaffian state belongs to the same topological class as the BCS pairing state and thus the latter must have the same non-abelian statistics . In our paper we verify this directly in the BCS framework as the property of solutions to Bogoliubov-de-Gennes equations. Our derivation provides an alternative (and possibly more transparent) point of view on the non-abelian statistics of half-quantum vortices as well as an additional verification of topological equivalence between Pfaffian and BCS states.
Let us begin our discussion with reviewing the properties of a half-quantum vortex. To be specific, we consider a chiral two-dimensional superconductor with the order parameter of $`A`$ phase of <sup>3</sup>He. The order parameter is characterized by the direction $`\widehat{𝐝}`$ of the spin triplet (the projection on which of the spin of the Cooper pair is zero) and by the overall phase $`\phi `$. The wave function of the condensate is
$$\mathrm{\Psi }_\pm =e^{i\phi }[d_x\left(\right|+|)+id_y\left(\right||)+d_z\left(\right|+|)](k_x\pm ik_y).$$
(1)
The $`\pm `$ signs denote the two possible chiralities of the condensate. The chirality breaks the time-reversal symmetry and means a non-zero angular momentum of the Cooper pairs. In a physical chiral superconductor there must exist domain walls separating domains of opposite chirality. Experimentally, domain walls may possibly be expelled from the sample by an external field which makes one of the chiralities energetically favorable. In our discussion we do not consider interaction of vortices with domain walls, but assume that the chirality is fixed over the region where the vortex braiding occurs (and takes positive sign in eq.(1)).
For the half-quantum vortex to exist, the vector $`\widehat{𝐝}`$ must be able to rotate (either in a plane or in all three dimensions). An important observation is that the order parameter maps to itself under simultaneous change of sign of the vector $`\widehat{𝐝}`$ and shift of the phase $`\phi `$ by $`\pi `$: $`(\phi ,\widehat{𝐝})(\phi +\pi ,\widehat{𝐝})`$. The half-quantum vortex then combines rotations of the vector $`\widehat{𝐝}`$ by $`\pi `$ and of the phase $`\phi `$ by $`\pi `$ on going around the vortex core (Fig. 1). This vortex is topologically stable, i.e. it cannot be removed by a continuous (homotopic) deformation of the order parameter.
Without loss of generality, we consider the vector $`\widehat{𝐝}`$ rotating in the $`x`$-$`y`$ plane. The direction of the rotation of the phase $`\phi `$ may either coincide or be opposite to the chirality of the condensate, which defines either a positive ($`\mathrm{\Phi }=1/2`$) or a negative ($`\mathrm{\Phi }=1/2`$) vortex respectively.
There are also two possible directions of rotating the vector $`\widehat{𝐝}`$. If the vector $`\widehat{𝐝}`$ is confined to a plane (i.e. takes values on a one-dimensional circle) by an anisotropy interaction, this gives two possible winding numbers of the vector $`\widehat{𝐝}`$ ($`m=\pm 1/2`$). If the vector $`\widehat{𝐝}`$ is not confined to a plane, but takes values on a two-dimensional sphere, the two directions of winding vector $`\widehat{𝐝}`$ are topologically equivalent, and there is no additional quantum number characterizing the vortex.
We further restrict our discussion to $`\mathrm{\Phi }=1/2`$ vortices ($`\mathrm{\Phi }=1/2`$ vortices have a slightly different structure of the quasiparticle eigenfunctions, but their spectrum and braiding statistics are the same as for $`\mathrm{\Phi }=1/2`$ vortices). At such a vortex, the condensate wave function (1) takes the form:
$$\mathrm{\Psi }(r,\theta )=\mathrm{\Delta }(r)\left[e^{i\theta }\right|+|](k_x+ik_y).$$
(2)
Here $`r`$ and $`\theta `$ are the polar coordinates in the vortex, the windings of the phase $`\phi `$ and of the vector $`\widehat{𝐝}`$ have been taken into account.
The Bogoliubov-de-Gennes Hamiltonian in this case decouples into the two Hamiltonians for spin-up and spin-down electrons. The spin-down sector contains no vortex and no subgap states. The Hamiltonian of the spin-up sector is
$$H=d^2r[\mathrm{\Psi }_{}^{}(\frac{p^2}{2m}\epsilon _F)\mathrm{\Psi }_{}+e^{i\theta }\mathrm{\Delta }(r)\mathrm{\Psi }_{}^{}(_x+i_y)\mathrm{\Psi }_{}^{}+\mathrm{h}.\mathrm{c}.].$$
(3)
Thus it may be considered as a single-quantum vortex in a superconductor of spinless fermions. The quasiparticles do not have a definite spin projection, but are mixtures of spin-up electrons and spin-down holes:
$$\gamma ^{}=u\mathrm{\Psi }_{}^{}+v\mathrm{\Psi }_{}.$$
(4)
Bogoliubov-de-Gennes equations for $`(u,v)`$ are obtained as $`[H,\gamma ^{}]=E\gamma ^{}`$. They are identical to those for a single-quantum vortex (with the vector $`\widehat{𝐝}`$ constant in space) and were solved by Kopnin and Salomaa in the context of superfluid <sup>3</sup>He vortices . The low-energy spectrum is
$$E_n=n\omega _0,$$
(5)
where $`\omega _0\mathrm{\Delta }^2/\epsilon _F`$ is the level spacing. The quantum number $`n`$ takes integer values (which distinguishes $`p`$-wave vortex states from Caroli-de-Gennes-Matricon states in $`s`$-wave vortices ) and has the meaning of the angular momentum of the quasiparticle.
In contrast to the single-quantum vortex considered by Kopnin and Salomaa, in the half-quantum vortex there is an additional relation between positive- and negative-energy eigenstates, namely $`\gamma ^{}(E)=\gamma (E)`$. In other words, the solutions with positive and negative energies are the creation and annihilation operators for the same fermionic level. Therefore the number of degrees of freedom in a half-quantum vortex is one half of that in a single-quantum vortex. The zero-energy level becomes a Majorana fermion:
$$\gamma ^{}(E=0)=\gamma (E=0).$$
(6)
It is worth mentioning that this Majorana fermion is stable with respect to any local perturbation including external potential, electromagnetic vector potential, local deformations of the order parameter, spin-orbit interaction and Zeeman splitting (in a single-quantum vortex, only the first three of those perturbations preserve the zero-energy level ). We can easily prove it with continuity considerations. Indeed, suppose that we gradually increase perturbation to the vortex Hamiltonian (which includes both the spin-up and spin-down sectors). The levels will shift and mix, but they must do it continuously, and therefore the number of levels is preserved. Since it is half-integer without perturbation, it must remain half-integer for the perturbed Hamiltonian, i.e. the Majorana fermion survives the perturbation. This argument is valid as long as the perturbation is sufficiently small so that the low-lying states remain localized in the vortex.
Before we turn to discussing the non-abelian statistics of vortices, let us see how the Majorana fermion $`\gamma (E=0)`$ transforms under $`U(1)`$ gauge transformations. If the overall phase of the superconducting gap shifts by $`\varphi `$, it is equivalent to rotating electronic creation and annihilation operators by $`\varphi /2`$: $`\mathrm{\Psi }_\alpha e^{i\varphi /2}\mathrm{\Psi }_\alpha `$, $`\mathrm{\Psi }_\alpha ^{}e^{i\varphi /2}\mathrm{\Psi }_\alpha ^{}`$. The solution $`(u,v)`$ transforms accordingly: $`(u,v)(ue^{i\varphi /2},ve^{i\varphi /2})`$. The important consequence of this transformation rule is that under change of the phase of the order parameter by $`2\pi `$ the Majorana fermion in the vortex changes sign: $`\gamma \gamma `$. This is an obvious consequence of the fact that the quasiparticle is a linear combination of fermionic creation and annihilation operators carrying charge $`\pm 1`$.
Now consider a system of $`2n`$ vortices, far from each other (at distances much larger than $`\xi _0v_F/\mathrm{\Delta }`$). To each vortex there corresponds one Majorana fermion (further we shall denote them by $`c_i`$, $`i=1,\mathrm{},2n`$) commuting with the Hamiltonian. They can be combined into $`n`$ complex fermionic operators and therefore give rise to the degeneracy of the ground state equal to $`2^n`$ (each fermionic level may be either filled or empty). If the vortices move adiabatically slowly so that we can neglect transitions between subgap levels, the only possible effect of such vortex motion is a unitary evolution in the space of ground states.
Let us fix the initial positions of vortices. Consider now a permutation (braiding) of vortices which returns vortices to their original positions (possibly in a different order). Such braid operations form a braid group $`B_{2n}`$ (multiplication in this group corresponds to the sequential application of the two braid operations) . This group may be described formally in the following way.
Suppose we order all particles along a fixed non-self-intersecting path. Then braid operations are generated by elementary interchanges of two neighboring particles. Denote such an elementary operation interchanging particles $`i`$ and $`i+1`$ by $`T_i`$ ($`i=1,\mathrm{},2n1`$). Then the group $`B_{2n}`$ is generated by operators $`T_i`$ modulo the following relations (see Fig. 2):
$`T_iT_j`$ $`=`$ $`T_jT_i,|ij|>1,`$ (7)
$`T_iT_jT_i`$ $`=`$ $`T_jT_iT_j,|ij|=1.`$ (8)
The braiding statistics is defined by the unitary operators in the space of ground states representing the braid operations from $`B_{2n}`$. Here an important reservation has to be made. When a single vortex moves along a closed loop, the multi-particle state acquires a phase proportional to the area inside the loop (every electron inside the loop effectively moves around the vortex). We shall disregard this effect and, as a consequence, loose information about the overall phase of the wave function. In other words, we shall speak about only a projective representation of the braid group $`B_{2n}`$. However, since the representation is multi-dimensional, the resulting projective representation is still nontrivial and transforms different states into each other — which implies the non-abelian statistics of vortices.
Since the Majorana fermions $`c_i`$ change sign under a shift of the superconducting phase by $`2\pi `$, we introduce cuts connecting vortices to the left boundary of the system (Fig. 3). We take the superconducting phase single-valued away from the cuts and jumping by $`2\pi `$ across the cuts. From examining Fig. 3 one easily obtains that the transformation exchanging the two vortices $`i`$ and $`i+1`$ (with no vortices between them) changes the phase of the order parameter at one of the vortices by $`2\pi `$, which results in the following transformation rule:
$$T_i:\{\begin{array}{cccc}c_i& & c_{i+1}& \\ c_{i+1}& & c_i& \\ c_j& & c_j& \text{for }ji\text{ and }ji+1\end{array}$$
(9)
This defines the action of $`T_i`$ on Majorana fermions. One easily checks that this action obeys the commutation relations (8).
Now the action of operators $`T_i`$ may be extended from operators to the Hilbert space. Since the whole Hilbert space can be constructed from the vacuum state by fermionic creation operators, and the mapping of the vacuum state by $`T_i`$ may be determined uniquely up to a phase factor, the action (9) of $`B_{2n}`$ on operators uniquely defines a projective representation of $`B_{2n}`$ in the space of ground states.
The explicit formulas for this representation may be written in terms of fermionic operators. Namely, we need to construct operators $`\tau (T_i)`$ obeying $`\tau (T_i)c_j[\tau (T_i)]^1=T_i(c_j)`$, where $`T_i(c_j)`$ is defined by (9). If we normalize the Majorana fermions by
$$\{c_i,c_j\}=2\delta _{ij},$$
(10)
then the expression for $`\tau (T_i)`$ is
$$\tau (T_i)=\mathrm{exp}\left(\frac{\pi }{4}c_{i+1}c_i\right)=\frac{1}{\sqrt{2}}\left(1+c_{i+1}c_i\right)$$
(11)
(up to a phase factor).
This formula presents the main result of our calculation. On inspection, this representation coincides with that described by Nayak and Wilczek for the statistics of the Pfaffian state (our Majorana fermions correspond to the operators $`\gamma _i`$ in section 9 of their paper).
The two simplest examples of the representation (11) are the cases of two and four vortices. These examples were previously discussed to some extent in the Pfaffian framework in refs., and we review them here for illustration purposes.
In the case of two vortices, the two Majorana fermions may be combined into a single complex fermion as $`\mathrm{\Psi }=(c_1+ic_2)/2`$, $`\mathrm{\Psi }^{}=(c_1ic_2)/2`$. The ground state is doubly degenerate, and the only generator of the braid group $`T`$ is represented by
$$\tau (T)=\mathrm{exp}\left(\frac{\pi }{4}c_2c_1\right)=\mathrm{exp}\left[i\frac{\pi }{4}(2\mathrm{\Psi }^{}\mathrm{\Psi }1)\right]=\mathrm{exp}\left(i\frac{\pi }{4}\sigma _z\right),$$
(12)
where $`\sigma _z`$ is a Pauli matrix in the basis ($`|0`$, $`\mathrm{\Psi }^{}|0`$).
In the case of four vortices, the four Majorana fermions combine into two complex fermions $`\mathrm{\Psi }_1`$ and $`\mathrm{\Psi }_2`$ by $`\mathrm{\Psi }_1=(c_1+ic_2)/2`$, $`\mathrm{\Psi }_2=(c_3+ic_4)/2`$ (and similarly for $`\mathrm{\Psi }_1^{}`$ and $`\mathrm{\Psi }_2^{}`$). The ground state has degeneracy four, and the three generators $`T_1`$, $`T_2`$, and $`T_3`$ of the braid group are represented by
$`\tau (T_1)`$ $`=`$ $`\mathrm{exp}\left(i{\displaystyle \frac{\pi }{4}}\sigma _z^{(1)}\right)=\left(\begin{array}{cccc}e^{i\pi /4}& & & \\ & e^{i\pi /4}& & \\ & & e^{i\pi /4}& \\ & & & e^{i\pi /4}\end{array}\right),`$ (13)
$`\tau (T_3)`$ $`=`$ $`\mathrm{exp}\left(i{\displaystyle \frac{\pi }{4}}\sigma _z^{(2)}\right)=\left(\begin{array}{cccc}e^{i\pi /4}& & & \\ & e^{i\pi /4}& & \\ & & e^{i\pi /4}& \\ & & & e^{i\pi /4}\end{array}\right),`$ (14)
$`\tau (T_2)`$ $`=`$ $`\mathrm{exp}\left({\displaystyle \frac{\pi }{4}}c_3c_2\right)={\displaystyle \frac{1}{\sqrt{2}}}(1+c_3c_2)={\displaystyle \frac{1}{\sqrt{2}}}\left[1+i(\mathrm{\Psi }_2^{}+\mathrm{\Psi }_2)(\mathrm{\Psi }_1^{}\mathrm{\Psi }_1)\right]={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{cccc}1& 0& 0& i\\ 0& 1& i& 0\\ 0& i& 1& 0\\ i& 0& 0& 1\end{array}\right),`$ (15)
where the matrices are written in the basis ($`|0`$, $`\mathrm{\Psi }_1^{}|0`$, $`\mathrm{\Psi }_2^{}|0`$, $`\mathrm{\Psi }_1^{}\mathrm{\Psi }_2^{}|0`$).
There are two important properties of the representation (11) . The first one is that $`\tau (T)`$ are even in fermionic operators and therefore preserve the parity of the number of fermions (physically, this simply means that the superconducting Hamiltonian creates and destroys electrons only in pairs). Therefore the representation may be restricted to odd or even sector of the space of ground states, each of them containing $`2^{n1}`$ states (this degeneracy was also found for the Pfaffian state in refs. ). Still, in each of these subspaces the representation operators are non-trivial and non-commuting.
The second property of the representation (11) is that $`T_i^4`$ is represented by a scalar matrix (projectively equivalent to the unity matrix, since we disregard the overall phase). That is, an elementary interchange of two vortices repeated four times produces an identity operator (up to an overall phase).
Quite remarkably, our derivation of the non-abelian statistics only relies on the two facts: first, the flux quantization (half-quantum for spin-1/2 electrons or, equivalently, single-quantum for spinless fermions) and, second, that the Majorana fermions carry odd charge with respect to the vortex gauge field, i.e. they transform as $`c_ic_i`$ when the phase of the order parameter changes by $`2\pi `$. But these are quantization properties that depend only on the presence of the Majorana fermion in the vortex spectrum, but not on the exact form of the Hamiltonian. Therefore, if we introduce disorder or other local perturbation in the BCS Hamiltonian (such as electromagnetic vector potential, spin-orbit scattering or local deformation of the order parameter), then not only the Majorana fermions survive, but also the braiding statistics (11) remains unchanged (provided the Majorana fermions stay localized in vortices). Thus we may speak of the topological stability of the non-abelian statistics (11).
Finally, we mention that the operators $`\tau (T_i)`$ have also been discussed in the context of quantum computation as part of a universal set of operators . Also, non-abelian anyons provide a topologically stable realization of unitary operators for quantum computing . Thus, should $`p`$-wave superconductors with sufficiently large $`T_c`$ (or, equivalently, large $`\omega _0`$) be discovered, they may provide a promising hardware solution for quantum computation.
The author thanks M. V. Feigelman for suggesting this problem and for many fruitful discussions. Useful discussions with G. E. Volovik, C. Nayak, I. Gruzberg, and M. Zhitomirsky are gratefully acknowledged. The author thanks Swiss National Foundation for financial support.
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# Effects of Boson Dispersion in Fermion-Boson Coupled Systems
## I Introduction
Interacting fermion-boson systems are very important in condensed matter physics and have been studied intensively. They are directly relevant to the description of electron-lattice interaction. Other problems can be mapped onto interacting fermions and bosons by means of the Hubbard-Stratonovich transformation. While the problem of a single fermion interacting with a boson field, i.e., the polaron problem, is well understood, a lot less is known about the many-fermion problem in interaction with a boson field; it is a full interacting many-body problem which is only tractable analytically in the adiabatic and the atomic limits.
In this paper we revisit the interacting and dispersive fermion-boson problem using dynamical mean-field (DMT) theory. This method reduces the quantum many-body problem to a quantum impurity model obeying a self-consistency condition. This method has been useful in describing strong coupling problems such as the Mott transition. There are several motivations for our work.
First, a DMF treatment of the bosonic and fermionic degrees of freedom taking into account the boson dispersion, requires an extension of the DMF equations where the bosonic propagator degrees of freedom are determined self-consistently. This represents a new type of self-consistent DMF equation, which so far has not been investigated. These equations are relevant to many problems, electron-phonon interactions, fermions interacting with spin fluctuations or among themselves via the long-ranged Coulomb interactions, and to the boson-fermion model.
Second, while the Mott transition in the Hubbard model is well understood using DMF methods, it is interesting to understand how it is modified by the variation of the frequency of the mode that mediates the interaction, or how the results are changed by the electron-phonon interactions. The approach discussed in this paper is a first step in this direction.
Finally, phonon dispersion effects are relevant to many systems. The Jahn-Teller or breathing-type phonons, for instance, seen in manganese oxides should be dispersive due to intersite coupling. A distortion of a MnO<sub>6</sub>-octahedron affects distortions of the neighbor octahedra, since the MnO<sub>6</sub>-octahedra share their oxygen atoms which leads to an intersite coupling. This may be relevant to fascinating orderings of lattice and charge in doped manganites.
We study the mutual feedback of fermionic and bosonic degrees of freedom in a very simple system of fermions interacting with one branch of bosons at half filling. However the methodology can be extended to other problems where similar DMF equations occur such as electron problems with long-ranged Coulomb interactions and the competition of magnetic order and the heavy fermion state, and to the boson-fermion mixture of high temperature superconductivity.
This paper is organized as follows. In the next section we discuss how DMF theory needs to be extended to fully include the feedback effects through fermion-boson interaction. Quantum Monte Carlo (QMC) method is introduced to solve the DMF equations for all parameter regions. We also discuss some technical points of the QMC relevant to this problem. The formalism is applied to demonstrate effects of boson dispersion in a wide region of parameters and the results are summarized in Sec. III. In Sec. IV, we discuss our main result: The existence of two distinct regimes of the DMF solutions. In the first regime, the feedback effects increase the fermion-boson coupling. In the second regime, which is strongly fluctuating, the boson dispersion accelerates the delocalization of fermions. Complete softening characterizes the crossover between these regimes. Section V is devoted to summary.
## II Dynamical Mean-Field Formalism and Hamiltonian
In this work, we discuss such feedback effects caused by the fermion-boson interaction using DMF theory. DMF theory provides a local view of a many-body problem in terms of an impurity model which satisfies a self-consistency condition. For general fermion-boson problems with a local interaction, the local action has the form
$`S_{\mathrm{eff}}`$ $`=`$ $`{\displaystyle 𝑑\tau 𝑑\tau ^{}\underset{\alpha }{}c_\alpha ^{}(\tau )𝒢_{0\alpha }^1(\tau \tau ^{})c_\alpha (\tau ^{})}`$ (1)
$`+`$ $`{\displaystyle 𝑑\tau 𝑑\tau ^{}\underset{\nu }{}x_\nu (\tau )𝒟_{0\nu }^1(\tau \tau ^{})x_\nu (\tau ^{})}`$ (2)
$`+`$ $`{\displaystyle 𝑑\tau \underset{\alpha _1\alpha _2\nu }{}\lambda _{\alpha _1\alpha _2\nu }c_{\alpha _1}^{}(\tau )c_{\alpha _2}(\tau )x_\nu (\tau )},`$ (3)
where $`𝒢_0`$ and $`𝒟_0`$ are the bare impurity Green’s functions for fermion and boson, respectively, which contain the dynamical information of the integrated other sites. Here $`c_\alpha `$ is the fermion annihilation operator and $`x_\nu `$ is the boson field. $`\lambda _{\alpha _1\alpha _2\nu }`$ denotes the coupling between fermions and bosons. The index $`\alpha `$ ($`\nu `$) denotes internal degrees of freedom of fermions (bosons) such as spins or orbitals of electrons (normal modes of phonons). We do not explicitly write the contribution from fermion interactions such as the Coulomb interaction since we focus on the effects of boson dispersions in this paper. However the action (3) is quite general which contains such fermion interactions through the Hubbard-Stratonovich transformation with continuous fields. Of course, alternatively, one can include additionally the fermion interactions according to the DMF theory for the Hubbard-type models.
The full Green’s functions are related to the bare ones by
$`𝒢_\alpha ^1(i\omega _n)`$ $`=`$ $`𝒢_{0\alpha }^1(i\omega _n)\mathrm{\Sigma }_\alpha (i\omega _n),`$ (4)
$`𝒟_\nu ^1(i\omega _n)`$ $`=`$ $`𝒟_{0\nu }^1(i\omega _n)\mathrm{\Pi }_\nu (i\omega _n),`$ (5)
at each Matsubara frequency $`\omega _n=(2n+1)\pi /\beta `$ for fermion and $`\omega _n=2n\pi /\beta `$ for boson, respectively ($`n`$ is an integer). $`\beta `$ is the inverse temperature. $`\mathrm{\Sigma }`$ and $`\mathrm{\Pi }`$ are the self-energy for fermion and boson, respectively. The Green’s functions for both fermion and boson are determined in a self-consistent way. This is achieved by the following set of self-consistency conditions,
$`𝒢_\alpha `$ $`=`$ $`{\displaystyle \underset{𝒒}{}}\left[i\omega _n+\mu ϵ_{𝒒\alpha }\mathrm{\Sigma }_\alpha (i\omega _n)\right]^1,`$ (6)
$`𝒟_\nu `$ $`=`$ $`{\displaystyle \underset{𝒒}{}}\left[(i\omega _n)^2\omega _{𝒒\nu }^2\mathrm{\Pi }_\nu (i\omega _n)\right]^1,`$ (7)
where $`ϵ_{𝒒\alpha }`$ and $`\omega _{𝒒\nu }`$ give the dispersion relations for fermions and bosons as a function of the wave number $`𝒒`$, respectively. $`\mu `$ is the chemical potential to control the density of fermions. Here the bosons are described as harmonic oscillators. The condition (7) is modified according to the property of boson degrees of freedom.
Previous studies of other models have indicated that the results of the DMF theory can give useful insights into three-dimensional systems. We have therefore taken the dispersions $`ϵ_𝒒`$ and $`\omega _𝒒`$ that correspond to a semicircular density of states, see the details in Sec. III A. These DMF equations are exact for a model where the fermions and the bosons have random hopping on lattice sites.
The self-consistency loop is closed as follows: The effective action (3) is solved for given bare impurity Green’s functions $`𝒢_0`$ and $`𝒟_0`$ to obtain the full Green’s functions $`𝒢`$ and $`𝒟`$. The self-energy $`\mathrm{\Sigma }`$ and $`\mathrm{\Pi }`$ are calculated by the relations (4) and (5), and used to obtain the Green’s functions through the self-consistency conditions (6) and (7). New bare impurity Green’s functions are calculated by the relations (4) and (5) again. This loop is iterated until all the quantities are converged. In this way, both fermionic and bosonic dispersions are renormalized through the fermion-boson interaction, and the mutual feedback effects are fully included.
The above DMF equations assume that no symmetry breaking is present in the system although the extension to phases with broken symmetry is straightforward. And they can be derived from an electron-phonon model:
$$=_\mathrm{F}+_\mathrm{B}+_\mathrm{I},$$
(8)
where
$`_\mathrm{F}`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}{\displaystyle \underset{𝒒}{}}\left(ϵ_{𝒒\alpha }\mu \right)c_{𝒒\alpha }^{}c_{𝒒\alpha }+_{\mathrm{int}},`$ (9)
$`_\mathrm{B}`$ $`=`$ $`{\displaystyle \underset{\nu }{}}{\displaystyle \underset{𝒒}{}}\omega _{𝒒\nu }\left(a_{𝒒\nu }^{}a_{𝒒\nu }+{\displaystyle \frac{1}{2}}\right),`$ (10)
$`_\mathrm{I}`$ $`=`$ $`{\displaystyle \underset{\alpha _1\alpha _2\nu }{}}{\displaystyle \underset{𝒑,𝒒}{}}\stackrel{~}{\lambda }_{𝒒\alpha _1\alpha _2\nu }c_{𝒑+𝒒\alpha _1}^{}c_{𝒑\alpha _2}\left(a_{𝒒\nu }+a_{𝒒\nu }^{}\right),`$ (11)
where $`a_𝒒`$ is the boson annihilation operator which is related with the boson field by $`x_𝒒=(2M\omega _𝒒)^{1/2}(a_𝒒+a_𝒒^{}).`$
The model (8) has been intensively studied using DMF methods in the limit of zero boson dispersion, i.e., in the Holstein model: Bosons with the same index $`\nu `$ have a same frequency (Einstein phonons) and the fermion-boson coupling is local as
$`_\mathrm{B}`$ $`=`$ $`{\displaystyle \underset{\nu }{}}\omega _{0\nu }{\displaystyle \underset{𝒒}{}}\left(a_{𝒒\nu }^{}a_{𝒒\nu }+{\displaystyle \frac{1}{2}}\right)`$ (12)
$`=`$ $`{\displaystyle \underset{\nu }{}}{\displaystyle \frac{M_\nu }{2}}{\displaystyle \underset{i}{}}\left(\dot{x}_{i\nu }^2+\omega _{0\nu }^2x_{i\nu }^2\right),`$ (13)
$`_\mathrm{I}`$ $`=`$ $`{\displaystyle \underset{\alpha _1\alpha _2\nu }{}}{\displaystyle \underset{i}{}}\lambda _{\alpha _1\alpha _2\nu }c_{i\alpha _1}^{}c_{i\alpha _2}x_{i\nu },`$ (14)
where the index $`i`$ denotes a lattice site. In the ground state, the possibility of charge-ordered or superconducting states has been intensively discussed for this model. Above the critical temperatures of these states, the crossover behavior is observed from the Fermi liquid with a mass enhancement in the weak coupling region to the so-called polaron which is a combined object between fermion and boson in the strong coupling region.
It is instructive to compare the present framework with the DMF theory for the problem without the boson dispersions such as the Holstein model. If bosons have no dispersion, that is, all $`\omega _𝒒`$ take the same value $`\omega _0`$ independent of $`𝒒`$, Eq. (7) is rewritten as
$$𝒟=\left[(i\omega _n)^2\omega _0^2\mathrm{\Pi }(i\omega _n)\right]^1.$$
(15)
Although the full Green’s function $`𝒟`$ contains a feedback effect in the self-energy $`\mathrm{\Pi }`$, the bare impurity Green’s function $`𝒟_0`$ is fixed at the noninteracting Green’s function given by
$$𝒟_0^{\mathrm{free}}=\left[(i\omega _n)^2\omega _0^2\right]^1$$
(16)
throughout the self-consistency iterations when we start from the solution $`𝒟_0=𝒟_0^{\mathrm{free}}`$. This is equivalent to the ordinary DMF theory for the Holstein model which does not need Eqs. (5) and (7). Compared to this, for the cases with finite bosonic dispersions, the bare impurity Green’s functions $`𝒟_0`$ is renormalized from $`𝒟_0^{\mathrm{free}}`$ in the iterations in our formalism.
The renormalization of $`𝒟_0`$ plays a crucial role because $`𝒟_0`$ is related to the effective interaction between fermions. If we integrate out the boson variables $`x`$ in the Hamiltonian, the effective interaction between fermions takes the form
$$\underset{\alpha _1\alpha _2\alpha _3\alpha _4}{}𝑑\tau 𝑑\tau ^{}c_{\alpha _1}^{}(\tau )c_{\alpha _2}(\tau )U_{\mathrm{eff}}(\tau \tau ^{})c_{\alpha _3}^{}(\tau ^{})c_{\alpha _4}(\tau ^{}),$$
(17)
where
$$U_{\mathrm{eff}}(\tau )=\lambda ^2𝒟_0(\tau ).$$
(18)
In the absence of the boson dispersion, since $`𝒟_0`$ is unchanged through the self-consistency loop as mentioned above, the effective interaction (18) is also unrenormalized. On the other hand, $`𝒟_0`$ is renormalized in our formalism for finite dispersions, which means that the effective interaction between fermions is renormalized by the mutual feedback of the fermion-boson coupling.
There are several techniques to solve the effective impurity problem with the action (3). In this work, we employ QMC method because it is an unbiased calculation and suitable to investigate all the parameter regions beyond perturbative regimes. In the QMC approach, the imaginary time is discretized into $`L`$ slices with the width $`\mathrm{\Delta }\tau `$ ($`\mathrm{\Delta }\tau =\beta /L`$). Continuous variables $`x_{\nu l}=x_\nu (\tau _l)`$ ($`\tau _l=l\mathrm{\Delta }\tau `$, $`l=1,2,\mathrm{},L`$) are randomly updated to $`x_{\nu l}^{}`$ with the probability
$$\underset{\alpha }{}\frac{det𝒢_\alpha }{det𝒢_\alpha ^{}}\frac{\mathrm{exp}\left[\mathrm{\Delta }\tau B(x_{\nu l}^{})\right]}{\mathrm{exp}\left[\mathrm{\Delta }\tau B(x_{\nu l})\right]},$$
(19)
where $`B(x_{\nu l})=_{j=1}^Lx_{\nu j}𝒟_{0\nu jl}^1x_{\nu l}`$ with $`𝒟_{0\nu jl}=𝒟_{0\nu }(\tau _j\tau _l)`$. The fermion Green’s functions $`𝒢`$ and $`𝒢^{}`$ are calculated by the standard algorithm for the configurations with $`x_{\nu l}`$ and $`x_{\nu l}^{}`$, respectively.
In actual QMC samplings, we consider both local and global updates for the continuous fields $`x_{\nu l}`$. The local update consists of sequential updates of the fields on each discretized point; a change from $`x_{\nu l}`$ to $`x_{\nu l}^{}=x_{\nu l}+r\delta `$ is attempted where $`r`$ is a random number between $`1`$ and $`1`$ and $`\delta `$ is a given amplitude. The global one is a simultaneous movement of all the fields at a same amount $`r\delta `$. The latter becomes important especially in the strong coupling region and/or at low temperatures where the fields $`x`$ show some orderings or are nearly ordered. The update amplitude $`\delta `$ is chosen to give an appropriate value of the acceptance ratio which is defined as the ratio of the number of accepted samples to the total number of trials.
The QMC calculations generally have the negative sign problem; the MC weight (19) can be negative for the general action (3), which leads to numerical instability in the QMC measurements. However, if fermions couple to bosons only in the diagonal form, that is, the coupling parameter $`\lambda _{\alpha _1\alpha _2\nu }`$ is nonzero only for the case of $`\alpha _1=\alpha _2`$, the MC weight (19) becomes positive definite. In this case, there is no negative sign problem for all parameters.
There are two sources of errors in the QMC calculations. One is a systematic error due to the discretization of the imaginary time, and the other is a statistical error from the random sampling. The former error is known to be proportional to $`(\mathrm{\Delta }\tau ^2)`$. Measurement is divided into several blocks to estimate the latter statistical error by the variance among the blocks. The size of each error depends on a specific form of models and parameters.
## III Results
### A Model and Parameters
We apply the new DMF framework proposed in the previous section to a case that the general Hamiltonian (8) contains two species of fermions and one branch of bosons. We set the mass $`M=1`$. The model is a straightforward extension of the Holstein model to include dispersive bosons, whose fermion-boson interaction is given by
$$_\mathrm{I}=\lambda \underset{i\alpha }{}\left(c_{i\alpha }^{}c_{i\alpha }\frac{1}{2}\right)x_i,$$
(20)
where the index $`\alpha `$ takes two values like spin degrees of freedom of electrons. The interaction is diagonal in the fermion index $`\alpha `$ so that the QMC does not suffer from the negative sign problem as mentioned in Sec. II. We take the coupling parameter $`\lambda `$ to be positive, which favors a doubly-occupied or an empty state on each site. Note that the model has the particle-hole symmetry at $`\mu =0`$.
The boson dispersion is taken into account through Eq. (7) in the present framework. We replace the summations over the wave number $`𝒒`$ in Eqs. (6) and (7) by the energy integrations as
$`𝒢(i\omega _n)={\displaystyle \frac{D_\mathrm{F}(\epsilon )d\epsilon }{i\omega _n+\mu \epsilon \mathrm{\Sigma }(i\omega _n)}},`$ (21)
$`𝒟(i\omega _n)={\displaystyle \frac{D_\mathrm{B}(\epsilon )d\epsilon }{(i\omega _n)^2\epsilon ^2\mathrm{\Pi }(i\omega _n)}},`$ (22)
where $`D_\mathrm{F}`$ and $`D_\mathrm{B}`$ are the the density of states for fermion and boson, respectively. In the following calculations, we assume semicircular density of states as
$`D_\mathrm{F}(\epsilon )`$ $`=`$ $`{\displaystyle \frac{2}{\pi W^2}}\sqrt{W^2\epsilon ^2},`$ (23)
$`D_\mathrm{B}(\epsilon )`$ $`=`$ $`{\displaystyle \frac{2}{\pi \omega _1^2}}\sqrt{\omega _1^2(\epsilon \omega _0)^2},`$ (24)
where $`W`$ is the half-bandwidth of the fermion density of states which is taken as unity hereafter ($`W=1`$); $`\omega _0`$ and $`\omega _1`$ are the center and the half-bandwidth of the boson density of states, respectively ($`\omega _0>0`$, $`\omega _0\omega _1>0`$). For the semicircular density of states, the integrations (21) and (22) are performed analytically which give
$`𝒢`$ $`=`$ $`{\displaystyle \frac{\zeta \sqrt{\zeta ^24t^2}}{2t^2}},`$ (25)
$`𝒟`$ $`=`$ $`{\displaystyle \frac{1}{\xi }}\left[{\displaystyle \frac{1}{\xi _{}+\sqrt{\xi _{}^2\omega _1^2}}}+{\displaystyle \frac{1}{\xi _++\sqrt{\xi _+^2\omega _1^2}}}\right],`$ (26)
where $`\zeta =i\omega _n+\mu `$ and $`\xi _\pm =\xi \pm \omega _0`$ with $`\xi ^2=(i\omega _n)^2\mathrm{\Pi }`$.
The shape of the boson density of states near the bottom is important because bosons at the band edge can be easily excited and strongly interact with fermions. The semicircular density of states (24) has an $`\epsilon ^{1/2}`$-singularity which is expected for bosons with ordinary cosine dispersions in three dimensions. Therefore we believe that the following results are qualitatively unchanged in realistic three-dimensional models. Results would be different for the two-dimensional density of states which has a step discontinuity at the band edges and results in very different DMF solutions.
In the absence of the boson dispersion ($`\omega _1=0`$), the model with the interaction (20) (the ordinary Holstein model) shows a charge ordering around half-filling ($`\mu =0`$) and superconductivity in doped regions at very low temperatures. In the following, we examine effects of boson dispersions in the low temperature region above and around these transition temperatures at half-filling ($`\mu =0`$) assuming no symmetry breaking. The calculations are mainly performed at $`\beta =8`$. We take $`\mathrm{\Delta }\tau =1/4`$ for which all the measured quantities are converged to the limit of $`\mathrm{\Delta }\tau 0`$ within the statistical errors. We have typically run $`1,000,000`$ MC steps for measurements; one MC sampling means a set of a sweep of local updates over the whole discretized points and a global update. Convergence in the self-consistency loop is usually rapid; typically 10 iterations are required to converge within the statistical errorbars when we start from the noninteracting Green’s functions. However in the strong coupling case, the iteration often suffers from an oscillation between two solutions. To avoid the oscillation, we make the iteration proceed by mixing the previous solutions.
### B Dispersionless Boson
First, we reconsider the limit without the boson dispersion, that is, $`\omega _1=0`$. In this case, we use the two parameters $`\omega _0`$ and $`U=\lambda ^2/M\omega _0^2`$ to characterize basic properties of the system. The first parameter $`\omega _0`$ describes the adiabaticity. In the adiabatic limit of $`\omega _00`$, the boson fields do not change in the imaginary time, that is, they behave as classical fields. In the opposite limit of $`\omega _0\mathrm{}`$, the bosons react instantaneously to fermion motions. Between these two limits, bosons with a finite $`\omega _0`$ mediate a retarded effective interaction which is given by $`U_{\mathrm{eff}}`$ in Eq. (18). The second parameter $`U`$ describes the magnitude of the effective interaction between fermions. Note that $`U=|U_{\mathrm{eff}}(\omega _n=0)|`$ in this dispersionless case, since the bare impurity Green’s function is given by the noninteracting one (16).
For a fixed value of $`\omega _0`$, the system behaves quite differently in the regions with $`U1`$ and $`U1`$. For small values of $`U`$, fermions are nearly free and each lattice site is in an empty, a singly-occupied, or a doubly-occupied state with almost equal probability at half-filling ($`\mu =0`$). If we define the probability $`P(x)`$ that the boson field $`x`$ lies in the interval between $`x`$ and $`x+\mathrm{\Delta }x`$, $`P(x)`$ shows a single broad peak centered at $`x=0`$. Compared to this, if $`U`$ becomes large, fermions strongly interact with each other, and a combined state between fermion and boson may be formed, which is called a small polaron. The polaron consists of double occupancy of fermions for the model with the interaction (20) (bipolaron). Then, the probability $`P(x)`$ displays a double peak at $`x=\pm \lambda /M\omega _0^2`$ which corresponds to the doubly-occupied and empty states. Figure 1 shows this behavior by changing the value of $`U`$ for the case of $`\omega _0=0.5`$. The single peak of the probability $`P(x)`$ appears for small $`U`$, while the double peaks are developed for U
>1
>𝑈1U\mathrel{\mathchoice{\lower 2.5pt\vbox{\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}1 as shown in Fig. 1(a). At the same time, in Fig. 1(b), the probability of the double occupancy $`P_\mathrm{D}`$ increases from $`1/4`$ for the noninteracting case to $`1/2`$ for the situation in which the system consists of only empty and doubly-occupied sites. The self-energy for fermion $`\mathrm{\Sigma }`$ is also enhanced by the effective interaction between fermions $`U`$. Figure 1(c) shows that the absolute value of the imaginary part of the self-energy as a function of Matsubara frequency is strongly enhanced by $`U`$. Note that the data for $`\omega _n>1/\mathrm{\Delta }\tau `$ contain no unbiased information. These clearly indicate the crossover from the weakly-correlated fermions in the small $`U`$ region to the small polarons in the large $`U`$ region.
The similar crossover is found for other values of $`\omega _0`$. Figures 2 and 3 show the results for $`\omega _0=2`$ and $`8`$, respectively. The value of $`U`$ for the crossover, which we call $`U^{}`$ hereafter, depends on the value of $`\omega _0`$. For example, for the case of $`\omega _0=0.5`$ in Fig. 1, the double-peak structure of $`P(x)`$ appears at $`U1`$, on the other hand, it does not appear up to $`U3`$ for $`\omega _0=8`$. This can be understood as follows: In the limit of $`\omega _0\mathrm{}`$, since the effective interaction becomes spontaneous, $`U_{\mathrm{eff}}(\tau )=U\delta (\tau )`$, the model maps onto an attractive Hubbard model in which the boson field corresponds to the continuous Hubbard-Stratonovich field. In the Hubbard model, it is known that the continuous field develops a double-peak distribution at $`U3`$, which corresponds to the opening of the Hubbard gap in the case of a repulsive interaction. On the other hand, in the opposite limit of $`\omega _00`$, the effective interaction becomes constant in the imaginary time, $`U_{\mathrm{eff}}(\tau )=U`$. This case is identical to an attractive Falicov-Kimball model in the limit of a continuous number of configurations for the static fields. In the adiabatic limit, fermions are localized at a smaller value of $`U`$ since fluctuations of the boson field is smaller in this case than in the anti-adiabatic limit. Then, the splitting of the distribution of $`x`$ should appear at a lower value of $`U`$. In the Falicov-Kimball model with a discrete static field, the critical value of $`U`$ is estimated to be $`1`$. The finite value of $`\omega _0`$ can interpolate these two limits. Thus, the value of $`U^{}`$ may change smoothly from $`U^{}3`$ in the limit of $`\omega _0\mathrm{}`$ to $`U^{}1`$ in the limit of $`\omega _00`$.
### C Dispersive Boson: Weak Coupling Limit
Now we discuss the cases with a finite bosonic dispersion; $`\omega _10`$. First, we study the weak coupling limit of $`W\omega _0`$ and $`U`$ which has been studied by a perturbation theory.
In this region, the finite width of the boson dispersion $`\omega _1`$ enhances the effective interaction between fermions. Figure 4(a) shows the bare impurity Green’s function for boson $`𝒟_0`$ as a function of Matsubara frequency for various values of $`\omega _1`$ for the case of $`\omega _0=0.5`$ and $`U=0.16`$ ($`\lambda =0.2`$). $`𝒟_0`$ is enhanced by the width of the dispersion $`\omega _1`$, which indicates that through the relation (18), the effective interaction between fermions $`U_{\mathrm{eff}}`$ is enhanced by $`\omega _1`$. This enhancement is also observed in the imaginary part of the fermion self-energy as shown in Fig. 4(b). At the same time, the probability of the double occupancy becomes large as shown in Fig. 4(c). These features are similar to those in Figs. 1-3 when the parameter $`U`$ increases in the small $`U`$ region. These results can be understood using a perturbative argument in Sec. IV.
### D Dispersive Boson: Atomic Limit
Next, we consider the limit of $`W\omega _0`$ and $`U`$ which has been studied based on the so-called small-polaron theory. In this limit, the coherent band motion of fermions in Eq. (9) is a perturbation on other terms of (10) and (11). The small-polaron theory is a perturbative approach from the atomic limit. The strong interaction between fermions and bosons leads to the formation of the small-polaron state as mentioned in the dispersionless case in Sec. III B.
In this region, as the weak coupling case in Sec. III C, the effective interaction between fermions is enhanced by the finite width of the boson dispersion. Figure 5(a) plots the bare impurity Green’s function for boson at zero Matsubara frequency for $`\omega _0=8`$ and $`U=9`$ ($`\lambda =24`$). A finite width of the boson dispersion $`\omega _1`$ enhances $`𝒟_0(\omega _n=0)`$. $`𝒟_0`$ shows the largest change at zero frequency as in Fig. 4(a). At the same time, the absolute value of the imaginary part of the fermion self-energy increases as shown in Fig. 5(b). We plot here the data at the smallest Matsubara frequency to show the behavior clearly. The double-peak structure of the probability function $`P(x)`$ shown in Fig. 3(a) at $`\omega _1=0`$, which indicates the formation of the small-polaron state, does not change for $`\omega _1`$ within statistical errorbars. This suggests that the finite width of the boson dispersion enhances the effective interaction while the small-polaron state remains stable. These features will be discussed based on the small-polaron theory in Sec. IV.
### E Dispersive Boson: Strong Fluctuation Regime
Here we go beyond the perturbative regimes studied in Sec. III C and III D. We consider the strong coupling case away from the anti-adiabatic limit, that is, $`U>W`$ and $`\omega _0W`$. It is difficult to study this regime by any perturbative and analytical approach because of strong fluctuations. Our DMF method including the fluctuation effects is applied to this regime without any difficulty.
Figure 6 shows the results for $`\omega _0=0.5`$ and $`U=2.56`$ ($`\lambda =0.8`$). As shown in Fig. 6(a), the absolute value of the bare impurity Green’s function for boson $`𝒟_0`$ decreases as $`\omega _1`$ increases. The imaginary part of the self-energy for fermion also decreases its absolute value as shown in Fig. 6(b). At the same time, the probability of the double occupancy $`P_\mathrm{D}`$ decreases from $`1/2`$ as shown in Fig. 6(c). Figure 6(d) shows that the double-peak structure of the probability $`P(x)`$ becomes unclear to merge into a single peak. All these features exhibit that the effective interaction between fermions $`U_{\mathrm{eff}}`$ is weakened and the small-polaron state becomes unstable for $`\omega _1`$. This is a striking contrast to the previous results in Sec. III C and III D. We will discuss a physical picture for this behavior in Sec. IV.
In the intermediate region, we find a crossover as the value of $`\omega _1`$ increases. Figure 7 shows this crossover for $`\omega _0=0.5`$ and $`U=0.64`$ ($`\lambda =0.4`$). For small values of $`\omega _1`$, we find a similar behavior as seen in Fig. 4; the bare impurity Green’s function for boson is enhanced and both the absolute value of the self-energy and the double occupancy increase as $`\omega _1`$. However, for ω1
>0.2
>subscript𝜔10.2\omega_{1}\mathrel{\mathchoice{\lower 2.5pt\vbox{\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}0.2, the behavior is reversed; all the three quantities turn to decrease as in Fig. 6. Therefore in this intermediate region, as the value of $`\omega _1`$ increases, the effective interaction between fermions is enhanced for small values of $`\omega _1`$, however, turns to be weakened for large values of $`\omega _1`$.
This crossover is closely related with complete softening of the boson field. Figure 8 shows the effective frequency of the boson field $`\omega ^{}`$ which is given by a pole of the Green’s function for boson as
$$\omega ^{}=\sqrt{(\omega _0\omega _1)^2+\mathrm{\Pi }(i\omega _n=0)},$$
(27)
where $`\mathrm{\Pi }`$ is the self-energy for boson. The frequency $`\omega ^{}`$ goes to zero at the value of $`\omega _1`$ where the crossover from the enhancement to the weakening of the effective interaction exhibited in Fig. 7.
### F Phase Diagram
We systematically investigate the crossover found in the previous section changing the parameter $`\omega _0`$ and $`U`$. Figure 9 shows the values of $`\omega ^{}`$ as a function of $`U`$ for the cases of (a) $`\omega _0=0.5`$ and (b) $`\omega _0=2.0`$. For finite values of the width $`\omega _1`$, the frequency $`\omega ^{}`$ goes to zero in both cases. We determine the crossover values of $`U`$ by this complete softening of $`\omega ^{}`$ for $`\omega _0`$ and $`\omega _1`$.
Figure 10 summarizes the phase diagram for the crossovers determined by the above criterion. This indicates the boundary between the weak-fluctuation and the strong-fluctuation regimes as discussed in Sec. IV. The most important point in this phase diagram is that especially for large $`\omega _0`$, the energy scale of $`U`$ for this crossover is quite different from $`U^{}`$ determined in Sec. III B. This suggests that there is another parameter which controls the onset of strong fluctuations as discussed in the next section IV.
## IV Discussion
In this section, we discuss the results obtained in Sec. III. Perturbative arguments are applied to discuss the enhancement of the effective interaction between fermions in the weak-coupling and the atomic regions. In the strong coupling regime away from the anti-adiabatic limit, the weakening of the effective interaction and the instability of the small-polaron state are discussed as a consequence of the strong fluctuations of the boson fields accompanied by complete softening. The phase diagram is examined to clarify the parameters which control the onset of the strong fluctuations.
In the weak coupling region, the dispersion width $`\omega _1`$ enhances the effective interaction between fermions in our DMF solutions. The absolute value of the self-energy for fermion is enhanced. A first-order perturbation in the coupling parameter $`\lambda `$ concludes that the self-energy $`\mathrm{\Sigma }`$ becomes larger as the width $`\omega _1`$ increases since $`𝒟_0`$ increases as $`\omega _1`$. Thus, perturbation theory suggests that the boson dispersion increases the effective interaction between fermions in the weak coupling limit. This enhancement can be understood intuitively as follows: In the weak coupling region, the density of states for both fermion and boson are not altered drastically by the fermion-boson interaction; a rigid-band picture should be justified. For a finite $`\omega _1`$, the band edge of the boson density of states is lowered linearly. Then, the effective interaction (18) is mainly mediated by the bosons near the band edge as
$$U_{\mathrm{eff}}(i\omega _n)\frac{\lambda ^2}{(i\omega _n)^2(\omega _0\omega _1)^2}.$$
(28)
Thus the absolute value $`|U_{\mathrm{eff}}|`$ becomes large as the value of $`\omega _1`$ increases. Therefore, the enhancement of the effective interaction in the weak coupling region can be understood as the decrease of the effective boson frequency. Our results in Sec. III C are consistent with this perturbative argument.
We now turn to the atomic limit, $`W\omega _0`$ and $`U`$. Now the fermion hopping term in Eq. (9) is a perturbation to the terms (10) and (11). If we apply the canonical transformation to diagonalize the unperturbed terms according to the small-polaron theory, we obtain the expression of the Hamiltonian as
$``$ $`=`$ $`{\displaystyle \underset{𝒒}{}}\omega _𝒒\left(a_𝒒^{}a_𝒒+{\displaystyle \frac{1}{2}}\right){\displaystyle \underset{i\alpha }{}}c_{i\alpha }^{}c_{i\alpha }\mathrm{\Delta }`$ (29)
$`+`$ $`{\displaystyle \underset{ij,\alpha }{}}t_{ij}c_{i\alpha }^{}c_{j\alpha }X_i^{}X_j,`$ (30)
where $`\mathrm{\Delta }`$ is the stabilization energy of polarons given by
$$\mathrm{\Delta }=\underset{𝒒}{}\frac{\lambda ^2}{\omega _𝒒},$$
(31)
and the operator $`X_i`$ takes the form
$$X_i=\mathrm{exp}\left[\underset{𝒒}{}e^{\mathrm{i}𝒒𝒓_i}\frac{\lambda }{\omega _𝒒}(a_𝒒a_𝒒^{})\right].$$
(32)
The third term of the Hamiltonian (30) indicates that the hopping occurs not as a bare fermion but as a combined object between fermion and boson. Each fermion is associated by a local boson. This is called the small-polaron state.
In the unperturbed state ($`t_{ij}=0`$), the polarons are almost localized in real space with the stabilization energy $`\mathrm{\Delta }`$ given by Eq. (31). When one increases the width of the boson dispersion $`\omega _1`$, the stabilization energy $`\mathrm{\Delta }`$ increases. This makes the polarons more strongly localized. The delocalization of the polarons is a second-order perturbation in $`t`$-term in Eq. (30) since the polarons consist of the double occupancy of fermions in this model. The hopping matrix by the second-order process is suppressed by $`\omega _1`$ because the intermediate state in the perturbation costs the energy $`\mathrm{\Delta }`$. The operator $`X_i`$ does not change the result in this limit of $`\omega _0W`$. Therefore a finite width of the boson dispersion suppresses the band motion of the polarons and enhances the effective interaction in this limit. Our results in Sec. III D are in good agreement with this argument based on the small-polaron theory.
We now discuss the strong coupling regime away from the anti-adiabatic limit studied in Sec. III E. In this regime, the polaron state is formed by the strong coupling, however the boson fields are loosely bound to the fermions due to a finite $`\omega _0`$, compared to the previous atomic limit with $`\omega _0W`$ where the boson fields react instantaneously to fermion motions. This leads to large fluctuations in the boson fields by the hopping of fermions. Thus this regime is characterized by these strong fluctuations, which is the reason why any perturbation cannot be applied. These fluctuations may in turn accelerate the delocalization of fermions through the mutual feedback effects of the fermion-boson coupling.
A finite width of the boson dispersion $`\omega _1`$ introduced in our calculations increases the fluctuations of the boson fields. The bosons are not localized and gain their kinetic energy through the dispersion. By tuning the width $`\omega _1`$, we can control the fluctuations of the boson fields by hand. Our results in Sec. III E clearly exhibited that the effective interaction between fermions is weakened by $`\omega _1`$. This is considered to be a consequence of the strong fluctuations of the boson fields enhanced by $`\omega _1`$ which tend to make fermions more delocalized. This behavior is elucidated for the first time by our method which fully includes the mutual feedback in many-body systems.
In the intermediate coupling region, a sharp crossover was found by changing the value of $`\omega _1`$ in Sec. III E. For small $`\omega _1`$, the effective interaction is enhanced by $`\omega _1`$. Since the fluctuations are small there, this may be smoothly connected to the behavior discussed in the weak-coupling or the atomic regions. The effective frequency $`\omega ^{}`$ defined by Eq. (27) becomes small but remains finite as in the perturbative regime although the reduction of $`\omega ^{}`$ is large and nonlinear in this nonperturbative regime. When the value of $`\omega _1`$ becomes large enough to soften the boson field completely ($`\omega ^{}0`$), fluctuations play a crucial role to enhance the delocalization of fermions. The boundaries in the phase diagram in Fig. 10 are the crossovers between the weak-fluctuation and the strong-fluctuation regimes.
In the dispersionless case in Sec. III B, we have found another crossover by the formation of the small polaron which is, for instance, characterized by the development of the double-peak structure in the probability $`P(x)`$. The critical value of $`U`$ for this crossover, $`U^{}`$, changes from $`U^{}3`$ in the limit of $`\omega _0W`$ (anti-adiabatic limit) to $`U^{}1`$ in the limit of $`\omega _0W`$ (adiabatic limit) as shown in the schematic phase diagram in the plane ($`1/\omega _0,U`$) in Fig. 11 (dotted gray line). On the other hand, the crossover to the strong fluctuation regime in Fig. 10 appears at much larger values of $`U`$ than $`U^{}`$ especially in the anti-adiabatic regime with large but finite $`\omega _0`$. This strongly suggests the importance of another energy scale to characterize the strong fluctuation regime which is not clearly found in the dispersionless case.
The importance of such characteristic energy scale has been pointed out also in a previous mean-field study. The parameter is defined by the ratio of the fermion-boson interaction to the spring constant of boson fields, $`\eta =\lambda /M\omega _0^2`$. In the case of $`\eta <1`$, since the fermion-boson interaction is weak compared to the stored energy in the boson field, the single-boson process should be important. In the case of $`\eta >1`$, the fermion-boson interaction is strong enough to excite a large numbers of bosons, i.e., multiboson processes become important. In the previous study, the importance of fluctuations of the boson fields has been suggested in the latter multiboson regime.
In the plane ($`1/\omega _0,U`$), the crossover between the single-boson and the multiboson regimes occurs at $`\eta =1`$, i.e., for $`U=\omega _0^2`$ and is shown in Fig. 11 (solid line). The line of $`\eta =1`$ becomes much larger than $`U^{}`$ in the anti-adiabatic regime. If we plot these values of $`U(\eta =1)`$ on the axis of $`\omega _1=0`$ in Fig. 10, the crossover boundaries seem to be smoothly connected to these values in the anti-adiabatic region. We demonstrate this behavior in Fig. 12 for $`\omega _0=4`$ and $`2`$ (gray lines). This indicates that in the anti-adiabatic regime, the line of $`\eta =1`$ corresponds to the crossover between the weak-fluctuation and the strong fluctuation regimes in the dispersionless case.
On the other hand, in the adiabatic regime with small but finite $`\omega _0`$, the value of $`U`$ for the boundary in Fig. 10 is not smoothly connected to the value of $`U`$ for $`\eta =1`$. For instance, in the case of $`\omega _0=0.5`$, $`\eta =1`$ gives $`U=1/4`$ which is much smaller than the boundary. This suggests that the condition of $`\eta >1`$ does not characterize the strong fluctuation regime in this adiabatic region.
To understand this behavior in the adiabatic regime, let us discuss the adiabatic limit of $`\omega _00`$. In this limit, the boson fields behave as classical fields which do not fluctuate in the imaginary time direction. Boson fluctuations only come from the fluctuation of the value of the field $`x`$ which is constant in time. Then even if the system is in the region of $`\eta >1`$, the fluctuations of the boson fields are small when $`U`$ is smaller than $`U^{}`$ since the fields feel a deep single-well potential as indicated in the probability $`P(x)`$. The boson fields begin to fluctuate when $`U`$ becomes comparable to $`U^{}`$ where the potential for $`x`$ softens around $`x=0`$. Therefore in the adiabatic limit, the strong fluctuation regime should be characterized not by $`\eta >1`$ but by $`U>U^{}`$.
Based on this argument, if we plot the value of $`U^{}`$ on the axis of $`\omega _1=0`$, the crossover boundaries in the adiabatic regime seem to be smoothly connected to these values for the cases of $`\omega _0=0.5`$ and $`0.25`$ in Fig. 12 (dotted lines).
These results reveal that the boson fluctuations which are strong enough to delocalize fermions appear in a different way in the anti-adiabatic and the adiabatic regimes. In the anti-adiabatic regime ($`\omega _0>W`$), the small-polaron state is formed at U
>U
>𝑈superscript𝑈U\mathrel{\mathchoice{\lower 2.5pt\vbox{\halign{$\m@th\displaystyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\m@th\textstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\m@th\scriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}{\lower 2.5pt\vbox{\halign{$\m@th\scriptscriptstyle\hfil#\hfil$\cr>\crcr\sim\crcr}}}}U^{*}. In the region with $`U>U^{}`$ and $`\eta <1`$, however, the fluctuations of the boson fields are small in the sense that the single-boson process contributes mainly and that the effective boson frequency is finite. If $`\eta `$ becomes larger than $`1`$, the boson field is softened and the boson fluctuations play a crucial role through the multiboson process. On the other hand, in the adiabatic regime ($`\omega _0<W`$), the fluctuations do not become large until $`UU^{}`$ even if $`\eta `$ is larger than $`1`$. The fluctuations are mainly the classical origin there.
To summarize, the strong fluctuations of the boson fields become important only when the conditions $`U>U^{}`$ and $`\eta >1`$ are both satisfied. These conditions are shown as the hatched area in Fig. 11. This area is strongly-correlated region for both fermions and bosons. The criterion for the formation of small polarons, $`UU^{}`$, corresponds to a competition between the kinetic energy of fermions $`W`$ and the effective interaction $`U`$. The line of $`U^{}`$ may be modified according to a specific form of the fermion-boson coupling. On the other hand, the criterion $`\eta 1`$ corresponds to a competition between the stored energy of the boson field and the coupling energy to fermions. Thus the hatched area in Fig. 11 is the region where correlations become strong in both standpoints of fermions and bosons.
A subtle problem remains open about the boson softening in the dispersionless case. As shown in Fig. 9, when $`\omega _1`$ is zero, the effective frequency $`\omega ^{}`$ becomes very small but remains finite for large values of $`U`$. We note that there are finite-temperature effects; $`\omega ^{}`$ is suppressed more strongly for lower temperatures. Unfortunately we cannot conclude in this study whether $`\omega ^{}`$ goes to zero even when $`\omega _1=0`$. In this dispersionless case, the boson density of states is a delta function at $`\epsilon =\omega _0`$, which is special since the shape of the boson density of states at the bottom is important as mentioned in Sec. III A. For instance, the step-like singularity in the two-dimensional density of states might prevent the boson field from complete softening at finite temperatures. Though further studies are necessary for the property of DMF equations for various types of density of states, we believe from the results in Figs. 11 and 12 that the nonlinear suppression of $`\omega ^{}`$ is relevant to strong fluctuations of bosons and that there are two important energy scales even when $`\omega _1=0`$.
## V Summary and Concluding Remarks
We have investigated the effects of the boson dispersion in a system of dynamical mean-field equations describing coupled fermion-boson systems. The analysis of the equations revealed that the boson dispersion plays a crucial role in a wide region of parameters. By introducing a parameter for the width of the dispersion in the model, we can control the fluctuations of the boson fields. To handle the boson fluctuations and the feedback effects, we have extended the dynamical mean-field theory to determine the Green’s functions for both fermion and boson in the self-consistent way. In the ordinary framework for the dispersionless case, the channel for the boson Green’s function is frozen in the sense that the bare impurity Green’s function is fixed and unrenormalized from the noninteracting one. The renormalization of the bare impurity Green’s function for boson is very important since the bare impurity Green’s function is directly related to the effective interaction between fermions. The equations in the extended dynamical mean-field theory are solved by using quantum Monte Carlo technique.
The main result in the models with dispersive bosons is that in the strong coupling regime away from the anti-adiabatic limit, the fluctuations of the boson fields become relevant to accelerate the delocalization of fermions. The effective interaction between fermions is weakened as the width of the boson dispersion increases in this regime. This behavior is explicitly shown for the first time by our method which fully includes the mutual feedback effects. The crossover to this nonperturbative regime is closely correlated with softening of the boson field. We have examined the phase diagram where this strong fluctuation occurs by tuning the coupling parameter and the width of the dispersion. The strong fluctuations to delocalize fermions become relevant when the small-polaron state is formed and the multiboson processes become important. The small polarons become stable when the effective interaction between fermions overcomes the fermion band energy. The multiboson regime is characterized by a coupling parameter larger than the boson energy. Thus the strong fluctuation regime is the strong correlated region for both fermions and bosons. As the coupling parameter increases, the boson fluctuations appear in a different way between in the adiabatic and the anti-adiabatic regimes. In the adiabatic regime, the fluctuations are mainly classical which are enhanced by the softening of the potential for the boson fields in the formation of the small-polaron state. On the other hand, in the anti-adiabatic regime, the small polarons are formed in the single-boson regime, where the dynamical fluctuations are small and the effective boson frequency is finite. The boson fluctuations do not play a crucial role until the system enters in the multiboson regime by complete softening of bosons.
The onset of the strong fluctuations occurs near the region where the boson degrees of freedom soften. In this paper we have studied the DMF equations in the absence of any freezing of the boson degrees of freedom. These effects together with generalizations to states with different symmetries and other generalizations are currently under investigation.
Our results imply that the behavior of the boson fluctuations may depend on the specific form of the boson density of states. Different forms of the density of states should be tested in the present DMF framework in a future study. Especially we are interested in the two-dimensional case with a step-like singularity at the edge which might be free from complete softening at finite temperatures. Boson fluctuations in this case may lead to light-mass bipolaronic states, which would give some insights into the high-temperature superconductivity in Cu-oxide materials where fermions strongly couple with spin fluctuations. We plan to understand this two-dimensional case in a later publication.
The dynamical mean-field equations allow us to vary the width of the boson dispersion in the calculations. This reveals the interesting properties in the strong fluctuation regime. The tuning of the electronic bandwidth has been the subject of a great deal of theoretical and experimental work . Our work suggests the possible interest of varying the boson dispersion experimentally even though this may be easier in systems where the bosons are spin fluctuations whose dispersion determined by exchange interactions can be controlled more easily than optical phonon dispersions. Another possibility may be the realization of the dynamical mean-field theory in a random model.
There are many materials which satisfy the above conditions for the strong fluctuation regime. In many physical situations, the fermion bandwidth $`W`$ is large or comparable to $`\omega _0`$, which makes possible to access to the strong fluctuation regime by a relatively weak coupling. Our method provides a powerful theoretical tool to examine the physical properties in this regime. We can apply it to more realistic models including orbital degrees of freedom of electrons, different normal modes of phonons, or interactions between fermions. Such extensions are now under investigation.
## Acknowledgement
Y. M. acknowledges the financial support of Research Fellowships of Japan Society for the Promotion of Science for Young Scientists. G. K. is supported by the NSF under DMR 95-29138.
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# Scaling behavior for finite 𝒪(𝑛) systems with long-range interaction
## I introduction
The theory of continuous phase transitions is based on the hypotheses that at temperatures close to the critical $`T_c`$, there is only one dominating length scale related with the critical behavior of the system. Because of the divergent nature of the correlation length as the critical point is approached, the microscopic details of the system becomes irrelevant for the critical exponents describing the singular dependence of the thermodynamic functions. This intuitive picture is based on the grounds of the renormalization group treatment of second order phase transitions.
Scaling is a central idea in critical phenomena near a continuous phase transition and in the field theory when we are interested in the continuum limit . In both cases we are interested in the singular behavior emerging from the overwhelming large number of degrees of freedom, corresponding to the original cutoff scale, which need to be integrated out leaving behind long-wave length which vary smoothly. Their behavior is controlled by a dynamically generated length scale: the correlation length $`\xi _b`$. Such a fundamental idea is difficult to test theoretically because it requires a study of a huge number of interacting degrees of freedom. Experimentally, however, one hopes to be able to study scaling in finite systems near a second order phase transition. Namely the system is confined to a finite geometry and the finite-size scaling theory is expected to describe the behavior of the system near the bulk critical temperature (for a review on the finite-size scaling theory see Ref. ).
The $`𝒪(n)`$-symmetric vector models are extensively used to explore the finite-size scaling theory, using different methods and techniques both analytically and numerically. The most thoroughly investigated case is the particular one corresponding to the limit $`n=\mathrm{}`$ (this limit includes also the mean spherical model) . In this limit, these models are exactly soluble for arbitrary dimensions and in a general geometry. These investigations were devoted exclusively to systems with short (including nearest neighbors) as well as long-range forces decaying with the interparticle distance in a power law. For finite $`n`$ the most frequently used analytical method is that of renormalization group . However this is limited to the case of short-range interaction. The crossover from long to short-range forces was discussed in Ref. , where it has been found that renormalized values of the temperature and the coupling constant are continuous functions of the parameter controlling the range of the interaction, when this approaches the value $`2`$ characterizing the short-range force potential. The case of pure long-range interaction was investigated very recently in Ref. (a comment on the method and the results obtained there is presented in Section IV). In the mean time, a special attention was devoted to the investigation of finite size-scaling for the mean-spherical model with long-range interaction (for a review see Ref. and references therein).
In recent years there has been an increasing interest in the numerical investigation of the critical properties of systems with long-range interaction decaying at large distances $`r`$ by a power-law as $`r^{d\sigma }`$, where $`d`$ is the space dimensionality and $`\sigma `$ is the parameter controlling the range of the interaction. The mostly used technique for this achievements is the Monte Carlo method. This method was used to investigate the critical properties of Heisenberg ferromagnetic systems as well as Ising models . Nevertheless all the analysis there was concentrated on systems with classical critical behavior in the sense that the critical exponents are given by Landau theory.
In this paper we present a detailed investigation of the finite-size scaling properties of the field theoretic $`𝒪(n)`$ vector $`\phi ^4`$ model with long-range interaction. We will also check the influence of the interaction range on the critical behavior. These interactions enter the exact expressions for the free energy only through their Fourier transform, which leading asymptotic is $`U(q)q^\sigma ^{}`$, where $`\sigma ^{}=\mathrm{min}(\sigma ,2)`$ . As it was shown for bulk systems by renormalization group arguments $`\sigma 2`$ corresponds to the case of finite (short) range interactions, i.e. the universality class then does not depend on $`\sigma `$ . Values satisfying $`0<\sigma <2`$ correspond to long-range interactions and the critical behavior depends on $`\sigma `$. With the renormalization group treatment it has been found that the critical behavior depend on the small parameter $`\epsilon =2\sigma d`$, where $`2\sigma `$ corresponds to the upper critical dimension . According to the above reasoning one usually considers the case $`\sigma >2`$ as uninteresting for critical effects, even for the finite-size treatments . So, here we will consider only the case $`0<\sigma 2`$.
Here, we will provide a systematic and controlled approach to the quantitative computation of the thermodynamic momenta, usually used in numerical analysis. These momenta are related to the Binder’s cumulant and to various thermodynamic functions like the susceptibility. We will concentrate on the scaling properties of the coupling constants defining the system in the vicinity of the critical point. Our method is quite general and should apply to a large extend to the investigation of finite-size scaling in systems with long-range interaction in the vicinity of the critical point.
The plan of the paper is as follows. In Section II we review, briefly, the $`\phi ^4`$-model with long-range interaction and discuss its bulk critical behavior. Section III is devoted to the explanation of the methods used here to achieve our analysis. We end the section with the computation of some thermodynamic quantities of interest. In Section IV we discuss our results briefly. In the remainder of the paper we present some details of the calculations of some formula used throughout the paper.
## II Finite-size scaling for systems with long-range interactions
In the vicinity of its critical point the Heisenberg model, with long-range interaction decaying as power-law, is equivalent to the $`d`$-dimensional $`𝒪(n)`$-symmetric model
$$\beta \left\{\phi \right\}=\frac{1}{2}_Vd^dx\left[\left(^{\sigma /2}\phi \right)^2+r_0\phi ^2+\frac{1}{2}u_0\phi ^4\right],$$
(1)
where $`\phi `$ is a short hand notation for the space dependent $`n`$-component field $`\phi (x)`$, $`r_0=r_{0c}+t_0`$ ($`t_0TT_c`$) and $`u_0`$ are model constants. $`V`$ is the volume of the system. In equation (1), we assumed $`\mathrm{}=k_B=1`$ and the size scale is measured in units in which the velocity of excitations $`c=1`$. We note that the first term in the model denotes $`𝒌^\sigma |\phi (𝒌)|^2`$ in the momentum representation where the parameter $`0<\sigma 2`$ takes into account short-range as well as long-range interactions. $`\beta `$ is the inverse temperature. The nature of the spectrum suggests that the critical exponent $`\eta =2\sigma `$ . Here we will consider periodic boundary conditions. This means
$$\phi (x)=\frac{1}{\sqrt{V}}\underset{𝒌}{}\phi (𝒌)\mathrm{exp}\left(i𝒌x\right),$$
(2)
where $`𝒌`$ is a discrete vector with components $`k_i=2\pi n_i/L`$ $`\left(n_i=0,\pm 1,\pm 2,\mathrm{},i=1,\mathrm{},d\right)`$ and a cutoff $`\mathrm{\Lambda }a^1`$ ($`a`$ is the lattice spacing). In this paper, we are interested in the continuum limit i.e. $`a0`$. As long as the system is finite we have to take into account the following assumptions $`L/a\mathrm{}`$, $`\xi _b\mathrm{}`$ while $`\xi _b/L`$ is finite.
Fisher et al. and Yamazaki et al. have shown that for the model under consideration the Landau theory holds for $`d>2\sigma `$. In the opposite case i.e. $`d<2\sigma `$ an expansion in powers of $`\epsilon =2\sigma d=4d2\eta `$ takes place, where $`2\sigma `$ plays the role of the upper critical dimension. We will present the renormalized parameters which characterize the bulk critical behavior and appear in the scaling functions. Since the computations are standard , we will be quite brief.
The application of the renormalized theory, above the critical temperature, to the model Hamiltonian requires a scaling field amplitude $`Z`$, a coupling constant renormalization $`Z_g`$ and a renormalization of the $`\phi ^2`$ insertions in the critical theory $`Z_t`$. In term of these, we define as usual
$$t=ZZ_t^1(r_0r_{0c})\text{and }g=\mathrm{}^\epsilon Z^2Z_g^1u_0.$$
(3)
In the remainder we will work in units where the reference length $`\mathrm{}`$ is set to unity. To one loop order the renormalization constants in the minimal subtraction scheme are given by
$$Z=1+𝒪(\widehat{g}^2)$$
(5)
$$Z_t=1+\frac{n+2}{\epsilon }\widehat{g}+𝒪(\widehat{g}^2)$$
(6)
$$Z_g=1+\frac{n+8}{\epsilon }\widehat{g}+𝒪(\widehat{g}^2)$$
(7)
In Eqs.(3),
$$\widehat{g}=g\frac{2}{(4\pi )^{d/2}\mathrm{\Gamma }(d/2)}=\frac{2g}{(4\pi )^\sigma \mathrm{\Gamma }(\sigma )}\left(1+\frac{\epsilon }{2}\left[\mathrm{ln}(4\pi )+\psi (\sigma )\right]+𝒪(\epsilon ^2)\right),$$
(8)
where $`\psi (x)`$ is the digamma function.
The fixed point of the $`\beta `$ function is at $`\widehat{g}=\widehat{g}^{}`$ with
$$\widehat{g}^{}=\frac{\epsilon }{n+8}+𝒪(\epsilon ^2).$$
(9)
Before starting to investigate the finite-size scaling in the field theoretical model under consideration, we shall recall briefly the corresponding renormalization group formalism. In the continuum limit, the lattice spacing completely disappears. The integration over wave vectors of the fluctuations are evaluated without cutoff and are convergent. When some dimensions of the system are finite the integrals over the corresponding momenta are transformed into sums. Since the lattice spacing is taken to be zero, the limits of the sums still tend to infinity.
From general renormalization group considerations an observable $`X`$, the susceptibility for example, will scale like :
$$X[t,g,\mathrm{},L]=\zeta (\rho )X[t(\rho ),g(\rho ),\mathrm{}\rho ,L],$$
(10)
where $`t`$ is the reduced temperature, $`g`$ a dimensionless coupling constant and $`L`$ the finite-size scale. The length scale $`\mathrm{}`$ is introduced in order to control the renormalization procedure.
It is known that in the bulk limit, when $`g(\rho )`$ approaches its stable fixed point $`g^{}`$ then we have
$$t(\rho )t\rho ^{1/\nu }\mathrm{and}\zeta (\rho )\rho ^{\gamma _x/\nu },$$
(11)
where $`\gamma _x`$ and $`\nu `$ are the bulk critical exponents measuring the divergence of the observable $`X`$ and the correlation length, respectively, in the vicinity of the critical point and $`\rho `$ is a scaling parameter. Using dimensional analysis together with equation (10) one gets
$$X[t,g,\mathrm{},L]=\zeta (\rho )X[t(\rho )(\rho \mathrm{})^2,g(\rho ),1,L/\mathrm{}\rho ].$$
(12)
Choosing the arbitrary parameter $`\rho =L/\mathrm{}`$, we obtain the well known finite-size scaling result
$$X[t,g,\mathrm{},L]=L^{\gamma _x/\nu }f\left(tL^{1/\nu }\right).$$
(13)
Here the function $`f(x)`$ is a universal function of its argument. In the remainder of this paper we will verify the scaling relation (13) in the framework of model (1).
## III Finite-size scaling below the upper critical dimension
### A Method
The method, we shall use here to analyze the finite-size scaling of the model under consideration, is originally due to Lüscher in his study on the quantum $`𝒪(n)`$ nonlinear $`\sigma `$ model in $`1+1`$ dimensions. An extension of the method was employed by Brezin and Zinn-Justin and by Rudnick, Guo, and Jasnow in their works on the finite-size scaling in systems with short-range potentials. Very recently it was used in the investigation of crossovers in quantum $`𝒪(n)`$ systems near their upper critical dimension . We will see here that the problem related to finite-size scaling in systems with long-range forces can be successfully analyzed by the same approach. Nevertheless, here we will observe the emergence of some subtleties, which need to be discussed.
The central idea of the method is that at finite linear size $`L`$ of the system, one can treat the $`𝒌=0`$ mode of the field $`\phi (x)`$, playing the role of the magnetization, separately from the non zero $`𝒌`$ modes. The non-zero modes are treated perturbatively using the loop expansion. They are integrated out to yield an effective Hamiltonian for the lowest mode only. All the modes being integrated out are regulated in the infrared by $`|𝒌|^\sigma `$ and consequently the process is necessarily free of infrared divergences. On the other hand the renormalizations of the bulk theory control the ultraviolet divergences at finite size. In other words if we define by
$$\varphi =\frac{1}{V}_Vd^dx\phi (x)$$
(14)
the total spin by unit volume, then, from $``$, we can get an effective Hamiltonian function of $`\varphi `$ after entirely integrating out the $`\phi (k0)`$ fields:
$$_{\mathrm{eff}}=\frac{L^d}{2}\left(R\varphi ^2+\frac{U}{2}\varphi ^4\right).$$
(15)
The coupling constants $`R`$ and $`U`$ are computed in powers of $`\epsilon `$, with the initial coupling constants renormalized as in their bulk critical theory. This approach will rule out all the ultraviolet divergences of the bulk critical point. The new coupling constants are necessarily free of all ultraviolet divergences since the theory is superrenormalizable . They are also free of infrared divergences as we are only integrating out finite modes. Obviously, these constants must obey the scaling forms,
$$R=L^{\eta 2}f_R\left(tL^{1/\nu }\right)\text{and}U=L^{d4+2\eta }f_U\left(tL^{1/\nu }\right)$$
(16)
for $`t0`$, where $`f_R`$ and $`f_U`$ are scaling functions which are properties of the bulk critical point. They are analytic at $`t=0`$. This is a consequence of the fact that only finite modes have been integrated out.
Once the scaling functions $`f_R`$ and $`f_U`$ are known one can attack the problem of computing observables in the $`\phi ^4`$ theory with the action $`_{\mathrm{eff}}`$. This theory is in dimension $`d`$ close to the upper critical dimension $`2\sigma `$ (not in $`d`$ close to the usual $`4`$), and the problem seems to be unsolvable. In the next section we will show that it is not the case.
In order to investigate the long distance physics of the finite system, one has to calculate thermal averages with respect to the new effective Hamiltonian defined in (15). They are related to the thermodynamic functions of the system under consideration. The averages of the field $`\varphi `$ are defined by
$$_{2p}=\left(\varphi ^2\right)^p=\frac{d^n\varphi \varphi ^{2p}\mathrm{exp}\left(_{\mathrm{eff}}\right)}{d^n\varphi \mathrm{exp}\left(_{\mathrm{eff}}\right)},$$
(17)
Using an appropriate rescaling of the field $`\varphi `$: $`\mathrm{\Phi }=\left(UL^d\right)^{1/4}\varphi `$, we can transform the effective Hamiltonian into
$$_{\mathrm{eff}}=\frac{1}{2}z\mathrm{\Phi }^2+\frac{1}{4}\mathrm{\Phi }^4,$$
(18)
where the scaling variable $`z=RL^{d/2}U^{1/2}`$ is an important quantity in the investigations of finite-size scaling in critical statics as well as in critical dynamics . With the effective Hamiltonian (18), we obtain the general scaling relation
$$_{2p}=L^{p(d2+\eta )}\frac{L^{p(d4+2\eta )/2}}{U^{p/2}}f_{2p}\left(RL^{2\eta }\frac{L^{(d4+2\eta )/2}}{U^{1/2}}\right)$$
(19)
for the momenta of the field $`\varphi `$. Having in mind Eqs. (16), we can write down Eq. (19) in the following scaling form
$$_{2p}=L^{p(d2+\eta )}_{2p}(tL^{1/\nu }),$$
(20)
in agreement with the finite-size scaling predictions of (13). In eq. (20), the function $`_{2p}(x)`$ are universal.
All the measurable thermodynamic quantities can be obtained from the momenta $`_{2p}`$. For example the susceptibility is obtained from
$$\chi =\frac{1}{n}_Vd^dx\phi (x)\phi (0)=L^{2\eta }_2(tL^{1/\nu }).$$
(21)
Another quantity of importance for numerical analysis of the finite-size scaling theory is the Binder’s cumulant defined by
$$B=1\frac{1}{3}\frac{_4}{_2^2}.$$
(22)
In the remainder of this section we concentrate on the computation of the coupling constants $`R`$ and $`U`$ of the effective Hamiltonian (15) for the system with long-range interaction decaying with the distance as a power law. As a consequence we will deduce results for the characteristic variable $`z=RU^{1/2}L^{2\eta \epsilon /2}`$, the susceptibility $`\chi `$ and the amplitude ratio $`r=_4/_2^2`$ entering the definition of the Binder’s cumulant.
### B Computation of the coupling constants $`R`$ and $`U`$
As we explained above, loop corrections will be treated perturbatively on the non-zero $`𝒌`$ modes. At the tree level (lowest order in $`\epsilon =2\sigma d`$) this procedure generates a shift of the critical temperature $`T_c`$ and a change of the coupling constant $`u_0`$ and additional operators involving powers of $`\phi `$ larger than $`4`$. The calculations will be performed in the renormalized theory. The renormalized coupling constant $`u_R`$ is expressed in terms of the dimensionless coupling constant $`g=\mathrm{}^\epsilon u_R`$ in which the parameter $`\mathrm{}`$ is an arbitrary length scale. Here we will work in system in which $`\mathrm{}=1`$. Throughout these calculations we use the minimal subtraction scheme. In this scheme, the counterterms of the massless theory including the $`\phi ^2`$ insertions are introduced. The one-loop counterterm for the coupling constant and the $`\phi ^2`$ insertion will be the only one relevant in the lowest corrections.
The finite-size correction to the renormalized coupling constant $`t`$ is given by
$$𝒲_{d,\sigma }^t(t,g,L)=(n+2)g\frac{1}{L^d}\underset{𝒌}{}^{}\frac{1}{t+|𝒌|^\sigma }$$
(23)
to one-loop order.
In order to investigate the finite-size scaling of the model under consideration one can use a suitable approach allowing to simplify the analytical calculations. In the case $`\sigma =2`$ it is possible to replace the summand by its Laplace transform. This is the so called Schwinger representation. The aim of this approach is to reduce the $`d`$-dimensional sum in the r.h.s of equation (23) to the one-dimensional effective problem. In the general case of arbitrary $`\sigma `$, one cannot just use the Schwinger transformation or at least in its familiar form. So we have to solve the problem by introducing some kind of generalization for it. In the spirit of the same problem a method to investigate the finite-size scaling in the framework of the mean spherical model was suggested in Ref. . The method is based upon the following genius identity
$$\frac{1}{1+z^\alpha }=_0^{\mathrm{}}𝑑x\mathrm{exp}\left(xz\right)x^{\alpha 1}E_{\alpha ,\alpha }\left(x^\alpha \right),$$
(25)
where the functions
$$E_{\alpha ,\beta }(z)=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{z^{\mathrm{}}}{\mathrm{\Gamma }\left(\alpha \mathrm{}+\beta \right)}$$
(26)
is the so called Mittag-Leffler type functions. For a more recent review on these functions and other related to them, and their application in statistical and continuum mechanics see Ref. . See also Ref. and Appendix A.
Using the identity (III B), one gets, after some algebra,
$$𝒲_{d,\sigma }^t(t,g,L)=(n+2)g\frac{L^{\sigma d}}{(2\pi )^\sigma }_0^{\mathrm{}}𝑑xx^{\frac{\sigma }{2}1}E_{\frac{\sigma }{2},\frac{\sigma }{2}}\left(x^{\sigma /2}\frac{tL^\sigma }{(2\pi )^\sigma }\right)\left[𝒜^d(x)1\right],$$
(28)
where
$$𝒜(x)=\underset{\mathrm{}=\mathrm{}}{\overset{\mathrm{}}{}}e^{x\mathrm{}^2}.$$
(29)
The analytic properties of the function $`𝒜(x)`$ are known very well. For large $`x`$, $`𝒜(x)1`$ decreases exponentially and the integral in the r.h.s of Eq. (28) converges at infinity. For small $`x`$, the Poisson transformation $`𝒜(x)=\left(\frac{\pi }{x}\right)^{\frac{1}{2}}𝒜\left(\frac{\pi ^2}{x}\right)`$ shows that $`𝒜(x)`$ converges.
For small $`x`$ the integral in the r.h.s of Eq. (28) has ultra violet divergence for $`\text{Re}d>\sigma `$. So, an analytic continuation in $`d`$ is required to give a meaning to the integral. Adding and subtracting the small asymptotic behavior of the function $`𝒜(x)`$, we get after some algebra
$`𝒲_{d,\sigma }^t(t,g,L)`$ $`=`$ $`(n+2)g{\displaystyle \frac{L^{\sigma d}}{(2\pi )^\sigma }}F_{d,\sigma }\left(tL^\sigma \right)`$ (32)
$`+2\pi (n+2)gL^{\sigma d}\left[(4\pi )^{d/2}\mathrm{\Gamma }\left({\displaystyle \frac{d}{2}}\right)\sigma \mathrm{sin}{\displaystyle \frac{d\pi }{\sigma }}\right]^1\left(tL^\sigma \right)^{d/\sigma 1},`$
where
$$F_{d,\sigma }\left(y\right)=_0^{\mathrm{}}𝑑xx^{\frac{\sigma }{2}1}E_{\frac{\sigma }{2},\frac{\sigma }{2}}\left(\frac{yx^{\sigma /2}}{(2\pi )^\sigma }\right)\left[𝒜^d(x)1\left(\frac{\pi }{x}\right)^{d/2}\right].$$
(33)
In the particular case $`\sigma =2`$, from Eq. (III B) we recover the result of Ref. .
By introducing the $`\phi ^2`$ counterterm insertion the renormalized coupling constant $`t`$ is replaced by $`tZ_t`$, where $`Z_t`$ is given by (6). Hence to one loop order we have
$$R=t\left(1+\widehat{g}\frac{n+2}{\epsilon }\right)+𝒲_{d,\sigma }^t(t,g,L).$$
(34)
At $`d=2\sigma `$, $`𝒲_{d,\sigma }^t(t,g,L)`$ has a simple pole. An expansion about this pole leads to the final expression
$$R=t+\frac{n+2}{\sigma }\widehat{g}t\mathrm{ln}t+2^{\sigma 1}(n+2)\mathrm{\Gamma }(\sigma )\widehat{g}L^\sigma F_{2\sigma ,\sigma }\left(tL^\sigma \right)+𝒪\left(\widehat{g}^2\right).$$
(35)
This result shows that, at the critical point, $`R`$ has the required scaling properties of Eq. (16), since
$$\nu ^1=\sigma \frac{n+2}{n+8}\epsilon +𝒪(\epsilon ^2).$$
(36)
For the finite system the renormalized coupling constant $`g`$, to one-loop order, is shifted by a quantity expressed in the form
$$𝒲_{d,\sigma }^g(t,g,L)=(n+8)g^2\frac{1}{L^d}\underset{𝒌}{}^{}\frac{1}{\left(t+|𝒌|^\sigma \right)^2}.$$
(37)
As one can see the summand here can be expressed as the first derivative of the summand of Eq. (23) with respect to $`t`$. So, the result for $`U`$ can be derived from that of $`R`$. Using this fact one gets
$$𝒲_{d,\sigma }^g(t,g,L)=(n+8)g^2\left[\frac{L^{2\sigma d}}{(2\pi )^\sigma }F_{d,\sigma }^{}\left(tL^\sigma \right)L^{2\sigma d}\frac{2}{\sigma (4\pi )^{d/2}}\frac{\mathrm{\Gamma }\left(2d/\sigma \right)\mathrm{\Gamma }\left(d/\sigma \right)}{\mathrm{\Gamma }(d/2)}\left(tL^\sigma \right)^{d/\sigma 2}\right],$$
(38)
where the prime indicates that we have the derivative of the function $`F`$ with respect to its argument.
At the fixed point one ends up with
$`U`$ $`=`$ $`g\left[1+\widehat{g}{\displaystyle \frac{n+8}{\sigma }}(1+\mathrm{ln}t)+\widehat{g}{\displaystyle \frac{n+8}{2^{1\sigma }}}\mathrm{\Gamma }(\sigma )F_{2\sigma ,\sigma }^{}\left(tL^\sigma \right)+𝒪\left(\widehat{g}^2\right)\right]`$ (39)
for the renormalized coupling constant $`U`$. Eq. (39) is obtained using the fact that at one-loop order the coupling constant is renormalized by $`Z_g`$ form Eq. (7). The obtained expression (39) shows that the coupling constant $`U`$ obeys the scaling law of Eq. (16). Note that $`U`$ has a finite limit as $`t0`$, i.e. it is analytic at the bulk critical temperature. Indeed as $`t0`$ one can use the expansion of the function $`F_{2\sigma ,\sigma }(y)`$ for small $`y`$ given by (see Appendix A)
$$F_{2\sigma ,\sigma }(y)=F_{2\sigma ,\sigma }(0)+2^\sigma y𝒞_\sigma \frac{2^{1\sigma }}{\sigma \mathrm{\Gamma }(\sigma )}y\mathrm{ln}y+𝒪(y^2),$$
(41)
where
$$𝒞_\sigma =\frac{1}{\mathrm{\Gamma }(\sigma )}_0^{\mathrm{}}\frac{du}{u}\left[E_{\frac{\sigma }{2},1}\left(\frac{u^{\sigma /2}}{(2\pi )^\sigma }\right)\frac{u^\sigma }{\pi ^\sigma }𝒜^{2\sigma }(u)+\frac{u^\sigma }{\pi ^\sigma }\right].$$
(42)
After substitution of (III B) in (39) the terms proportional to $`\mathrm{log}y`$ cancel, which shows that the coupling constant $`U`$ is finite at $`t=0`$. Whence, one gets
$$U=gL^\epsilon \left[1+\widehat{g}\frac{n+8}{\sigma }\left(1+\frac{\sigma }{2}\mathrm{\Gamma }(\sigma )𝒞_\sigma \right)+𝒪(\widehat{g}^2)\right]$$
showing that $`U`$ is analytic, as it should be, at the critical point.
### C Some thermodynamic quantities
#### 1 Shift of the critical point
It is obvious that the coupling constant $`R`$ in the effective Hamiltonian (15) is just the deviation of the temperature of the system from its ‘critical’ value. By setting $`t=0`$ in (35), we obtain an expression for the finite-size shift of the bulk critical temperature $`T_c`$. This is given by
$$T_cT_c(L)=\epsilon 2^{\sigma 1}\frac{n+2}{n+8}\mathrm{\Gamma }(\sigma )L^\sigma F_{2\sigma ,\sigma }(0),$$
(43)
where the coefficient $`F_{2\sigma ,\sigma }(0)`$, appearing in the right hand side of (43) can be evaluated for some particular values of the interparticle interaction range $`\sigma `$
$$F_{2\sigma ,\sigma }\left(0\right)=\{\begin{array}{cc}2\zeta \left(1/2\right),\hfill & \sigma =1/2,\hfill \\ 4\zeta \left(1/2\right)\beta \left(1/2\right),\hfill & \sigma =1,\hfill \\ 4.82271993,\hfill & \sigma =3/2,\hfill \\ 8\mathrm{ln}2,\hfill & \sigma =2.\hfill \end{array}$$
(44)
Here $`\zeta (x)`$ is the Riemann zeta function with $`\zeta \left(\frac{1}{2}\right)=1.460354508\mathrm{}`$ and $`\beta (x)`$ is the analytic continuation of the Dirichlet series:
$$\beta (x)=\underset{\mathrm{}=0}{\overset{\mathrm{}}{}}\frac{(1)^{\mathrm{}}}{\left(2\mathrm{}+1\right)^x},$$
with $`\beta \left(\frac{1}{2}\right)=0.667691457\mathrm{}`$. Remark that the function $`F_{2\sigma ,\sigma }(0)`$ increases as the parameter $`\sigma `$ vanishes.
In fact, since there is no true phase transition in the finite system, the critical temperature is shifted to a ‘pseudocritical’ temperature, $`T_c(L)`$, corresponding to the rounding of the thermodynamic singularities holding in the bulk limit. From (43) one remarks that $`T_c(L)`$ is larger than $`T_c`$, confirming previously obtained results in the framework of the spherical model . Notice also that for the shift exponent $`\lambda `$, we get $`\lambda =\sigma `$ to lowest order in $`\epsilon `$.
#### 2 Binder’s cumulant
In this subsection we are interested in the calculation of the amplitude ratio $`r=_4/_2^2`$ instead of the Binder’s cumulant from definition (22). This quantity can be expressed in power series of the scaling variable $`z=RL^{2\eta \epsilon /2}U^{1/2}`$ as
$`r`$ $`=`$ $`{\displaystyle \frac{n}{4}}{\displaystyle \frac{\mathrm{\Gamma }^2\left(\frac{1}{4}n\right)}{\mathrm{\Gamma }^2\left(\frac{1}{4}(n+2)\right)}}\{1z[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{4}(n+6)\right)}{\mathrm{\Gamma }\left(\frac{1}{4}(n+4)\right)}}+{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{4}(n+2)\right)}{\mathrm{\Gamma }\left(\frac{1}{4}n\right)}}2{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{4}(n+4)\right)}{\mathrm{\Gamma }\left(\frac{1}{4}(n+2)\right)}}]`$ (46)
$`+z^2[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{4}(n+6)\right)\mathrm{\Gamma }\left(\frac{1}{4}(n+2)\right)}{\mathrm{\Gamma }\left(\frac{1}{4}(n+4)\right)\mathrm{\Gamma }\left(\frac{1}{4}n\right)}}+3{\displaystyle \frac{\mathrm{\Gamma }^2\left(\frac{1}{4}(n+4)\right)}{\mathrm{\Gamma }^2\left(\frac{1}{4}(n+2)\right)}}n1]+𝒪\left(z^3\right)\}.`$
So, in order to obtain a result for $`r`$ it is enough to evaluate $`z`$ at the fixed point $`g^{}`$ and to deduce the value for the Binder’s cumulant. As we mentioned before this parameter appear in all thermodynamic functions through the momenta defined earlier in this paper.
At the fixed point $`g^{}`$ in the vicinity of the upper critical dimension, we obtain
$`z^{}`$ $``$ $`{\displaystyle \frac{RL^{2\eta }}{\sqrt{UL^\epsilon }}}|_{\mathrm{fixedpoint}}`$ (47)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{g^{}}}}[y{\displaystyle \frac{\epsilon }{2\sigma }}y(1{\displaystyle \frac{n4}{n+8}}\mathrm{ln}y)+2^{\sigma 1}\epsilon {\displaystyle \frac{n+2}{n+8}}\mathrm{\Gamma }(\sigma )F_{2\sigma ,\sigma }\left(y\right)`$ (49)
$`\epsilon 2^{\sigma 2}y\mathrm{\Gamma }(\sigma )F_{2\sigma ,\sigma }^{}\left(y\right)].`$
This result is obtained by using (36) and the fact that up to one loop order the terms proportional to $`\mathrm{ln}L`$ cancel. In Eq. (47), we introduce the scaling variable $`y=tL^{1/\nu }`$. Finally let us notice that from this equation one can see easily that $`z^{}`$ verifies the finite-size scaling hypotheses and consequently all the thermodynamic functions do.
At the critical temperature $`T_c`$ (i.e. $`t=0`$, and so $`y=0`$), we obtain
$$z_0^{}=\sqrt{\epsilon }\left[\frac{n+2}{\sqrt{n+8}}\sqrt{\frac{\mathrm{\Gamma }(\sigma )}{2\pi ^\sigma }}F_{2\sigma ,\sigma }(0)+𝒪(\epsilon )\right].$$
(50)
Numerical values for the amplitude ratio (46) can be obtained by replacing the value of $`z_0^{}`$ form (50) and taking some specific values of the small parameter $`\epsilon `$. Note that the scaling variable $`z`$ is proportional to $`\sqrt{\epsilon }`$ as it was found previously (see Ref. for example) in the case of short-range forces. Furthermore it coincides with the result of Ref. for the scaling variable $`x`$ in the case of long-range interaction. Consequently all the thermodynamic function will be computed in powers of $`\sqrt{\epsilon }`$.
#### 3 Magnetic Susceptibility
As we mentioned earlier, there is no phase transition in the finite system under consideration. Consequently there will be no ‘true’ correlation length. An expression for it can be deduced from that of the susceptibility (21) trough the relation:
$$\xi ^{2\eta }=\chi .$$
(51)
The analyticity of the susceptibility is a consequence of that the coupling constants $`R`$ and $`U`$.
From (21) in the region $`tL^\sigma 1`$ (i.e. $`z1`$), we obtain for the susceptibility
$`\chi `$ $`=`$ $`{\displaystyle \frac{L^\sigma }{\sqrt{\epsilon }}}{\displaystyle \frac{2\sqrt{2}}{\sqrt{(4\pi )^\sigma \mathrm{\Gamma }(\sigma )}}}{\displaystyle \frac{\sqrt{n+8}}{n}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{4}(n+2)\right)}{\mathrm{\Gamma }\left(\frac{n}{4}\right)}}[1z({\displaystyle \frac{n}{4}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{n}{4}\right)}{\mathrm{\Gamma }\left(\frac{1}{4}(n+2)\right)}}{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{1}{4}(n+2)\right)}{\mathrm{\Gamma }\left(\frac{n}{4}\right)}})`$ (53)
$`+z^2({\displaystyle \frac{1n}{4}}+{\displaystyle \frac{\mathrm{\Gamma }^2\left(\frac{1}{4}(n+2)\right)}{\mathrm{\Gamma }^2\left(\frac{n}{4}\right)}})\widehat{g}{\displaystyle \frac{n+8}{\sigma }}(1+{\displaystyle \frac{\sigma }{2}}\mathrm{\Gamma }(\sigma )𝒞_\sigma )+𝒪(\widehat{g}z,z^3)]`$
at the bulk critical point $`T_c`$.
To the lowest order in $`\epsilon `$, after taking the limit $`n\mathrm{}`$, we find that the correlation length scales like
$$\xi \epsilon ^{1/2\sigma }L,$$
confirming the results obtained in the spherical model and showing that this behavior is not a characteristic of the spherical limit i.e. $`n\mathrm{}`$.
In the region $`tL^\sigma 1`$ (i.e. $`z1`$), from Eq. (21), we get
$`\chi `$ $`=`$ $`{\displaystyle \frac{1}{t}}[1{\displaystyle \frac{n+2}{\sigma }}\widehat{g}\mathrm{ln}t2^{\sigma 1}(n+2)\mathrm{\Gamma }(\sigma )\widehat{g}\left(tL^\sigma \right)^1F_{2\sigma ,\sigma }\left(tL^\sigma \right)`$ (55)
$`{\displaystyle \frac{1}{2}}(n+2)(4\pi )^\sigma \mathrm{\Gamma }(\sigma )\widehat{g}\left(tL^\sigma \right)^2+𝒪\left(\widehat{g}^2\right)].`$
The function $`F_{d,\sigma }(y)`$ has the following large $`y`$ asymptotic behavior (see Appendix B)
$$F_{d,\sigma }(y)\frac{(2\pi )^\sigma }{y}+\frac{4^\sigma \pi ^{\sigma d/2}\mathrm{\Gamma }\left(\frac{d+\sigma }{2}\right)}{y^2\mathrm{\Gamma }(\frac{\sigma }{2})}\underset{𝒍}{}^{}\frac{1}{|𝒍|^{d+\sigma }}$$
(57)
for the case $`0<\sigma <2`$, and
$$F_{d,2}(y)\frac{4\pi ^2}{y}+d(2\pi )^{(5d)/2}y^{(d3)/4}e^\sqrt{y}$$
(58)
for the particular case $`\sigma =2`$. These results show that the last term in Eq. (55) is just canceled by the first term in Eqs. (III C 3).
In the case of long-range interaction $`0<\sigma <2`$, we obtain for the susceptibility
$$\chi =\chi _{\mathrm{}}\left[1\sigma \widehat{g}(n+2)2^{3\sigma 2}\left(tL^\sigma \right)^3\frac{\mathrm{\Gamma }\left(3\sigma /2\right)\mathrm{\Gamma }(\sigma )}{\mathrm{\Gamma }(1\sigma /2)}\underset{𝒍}{}^{}𝒍^{3\sigma }+𝒪(\widehat{g}^2)\right]$$
(59)
in agreement with the finite-size scaling hypothesis (21). Eq. (59) shows that the finite-size scaling behavior of the system is dominated by the bulk critical behavior, with small correction in powers of $`L`$. It should be noted that the above result cannot be continued smoothly to the case of short-range interaction $`\sigma =2`$, since then $`F_{4,2}(y)`$ (see Eq. (58)) falls off exponentially fast and, correspondingly, the the finite-size corrections to $`\chi `$ are exponentially small:
$$\chi =\chi _{\mathrm{}}\left[18\widehat{g}\sqrt{2\pi }(n+2)\left(tL^2\right)^{3/4}e^{\sqrt{t}L}\right]$$
(60)
At this point we are in disagreement with the statement given in Ref. that the approach used in Refs. yields an incorrect non-exponential result. Note that all our calculations are up to the order $`\widehat{g}^1`$. It is interesting to see what happen in higher order, e.g. $`\widehat{g}^2`$, in this case, however, we need to have at our disposal the corresponding high order terms in Eqs. (35) and (39). Indeed it is beyond the scope of the present study. First the power law fall off of the finite-size corrections to the bulk critical behavior, due to long-range nature of the interaction, was found in the framework of the spherical model . Here, we extended this result to finite $`n`$ using a perturbative approach.
## IV conclusions
In this paper, we have investigated the finite-size scaling properties in the $`𝒪(n)`$-symmetric $`\phi ^4`$ model with long-range interaction potential decaying algebraically with the interparticle distance. We have found that the methods developed in Refs. can be successfully extended to systems with long-range interaction by combining them with other known techniques. These techniques allow the investigation to be simplified and express the results for various thermodynamic functions in terms of simple and known mathematical functions.
Here we restricted our calculations to the critical domain $`TT_c`$ and investigated the model in dimensions less than the upper critical one, which turns out to be $`2\sigma `$ ($`0<\sigma 2`$). We constructed an effective Hamiltonian, from the initial one, with new coupling constants $`R`$ and $`U`$. These constants obey the scaling hypothesis (16). We found that the even momenta of the field $`\phi `$, related to the thermodynamics of the finite system, are scaling functions of the characteristic variable
$$z=RU^{1/2}L^{2\eta \epsilon /2}.$$
This variable has the required scaling form predicted by the finite-size scaling theory. From the obtained forms of the constant $`R`$ and $`U`$ one concludes that $`z`$ is a universal quantity, which does not depend of the details of the model.
We evaluated the finite-size shift, the susceptibility and the amplitude ratio $`r=_4/_2^2`$ at the tree level (lowest order in $`\epsilon `$). We observed that the critical behavior of the system is dominated by its bulk critical behavior away from the critical domain and that the finite-size scaling is relevant in the vicinity of the critical point. The amplitude ratio $`r`$ is evaluated as an expansion in powers of $`z\sqrt{\epsilon }`$. Our result is in consistency with that of Ref. . But it is disagreement with the numerical results of the same paper. There, it has been found, using the Monte Carlo method, that $`r`$ has an expansion in $`\epsilon `$ instead of its square root. At this time, we do not have a reasonable explanation of this fact. It is also possible that higher order in $`\epsilon `$ could improve the result. An amelioration of the result could also come from accounting finite cutoff effects, which were to be relevant in the investigation of finite systems and the comparison of the results with numerical works . However this is the subject of another publication.
Notice that in the only work devoted to the exploration of finite-size scaling in $`𝒪(n)`$ systems with long-range interaction (Ref. ) the pertinent integrals have to be evaluated only numerically, due to the choice of a parametrization that does not reduce the $`d`$-dimesional problem to the effective one dimensional one. The approach we used here is more efficient in the sense that the corresponding final expressions can be handled by analytical means. Consequently, we cannot make a direct comparison between the results of this paper and those obtained there.
Let us note that it would be interesting and useful to extend the result obtained here in the static limit to models including dynamics, since we believe that this is closely related to the extensively investigated filed of quantum critical points i.e. phase transitions occurring at zero-temperature. In particular we find it useful to investigate the critical dynamic of the quantum model considered in Ref in the large $`n`$ limit.
###### Acknowledgements.
The authors thank Dr. E. Koroutcheva for helpful discussion. This work is supported by The Bulgarian Science Foundation under Project F608.
## A Some properties of the Mittag-leffler type functions
The Mittag-leffler type functions are defined by the power series :
$$E_{\alpha ,\beta }(z)=\underset{k=0}{\overset{\mathrm{}}{}}\frac{z^k}{\mathrm{\Gamma }(\alpha k+\beta )},\alpha ,\beta >0.$$
(A1)
They are entire functions of finite order of growth. Let us mention that the function corresponding the particular case $`\beta =1`$ was introduced by Mittag-Leffler. These function are very popular in the field of fractional calculus (for a recent review see Ref. ).
One of the most striking properties of these functions is that they obey the following useful identity )
$$\frac{1}{1+z}=_0^{\mathrm{}}𝑑xe^xx^{\beta 1}E_{\alpha ,\beta }\left(x^\alpha z\right),$$
(A2)
which is obtained by means of term-by-term integration of the series (A1). The integral in Eq. (A2) converges in the complex plane to the left of the line $`\text{Re}z=1^{1/\alpha }`$, $`|\mathrm{arg}z|\frac{1}{2}\alpha \pi `$. The identity (A2) lies in the basis of the mathematical investigation of finite-size scaling in the spherical model with algebraically decaying long-range interaction (see Ref. and references therein).
In some particular cases the functions $`E_{\alpha ,\beta }(z)`$ reduces to known functions. For example, in the case corresponding to the short range case we have:
$$E_{1,1}(z)=\mathrm{exp}(z).$$
(A3)
Setting $`z=y^\alpha `$, $`y>0`$, and $`x=ty`$, we obtain the Laplace transform
$$\frac{y^{\alpha \beta }}{1+z^\alpha }=_0^{\mathrm{}}𝑑te^{zt}t^{\beta 1}E_{\alpha ,\beta }\left(t^\alpha \right)$$
(A4)
from which we derive the identity (III B) by setting $`\beta =\alpha `$.
The asymptotic behavior of the Mittag-Leffler functions is given by the Lemma :
Let $`0<\alpha <2`$, $`\beta `$ be an arbitrary complex number, and $`\gamma `$ be a real number obeying the condition
$$\frac{1}{2}\alpha \pi <\gamma <\mathrm{min}\{\pi ,\alpha \pi \}.$$
Then for any integer $`p1`$ the following asymptotic expressions hold when $`|z|\mathrm{}`$:
* At $`|\mathrm{arg}z|\gamma `$,
$$E_{\alpha ,\beta }(z)=\frac{1}{\alpha }z^{(1\beta )/\alpha }e^{z^{1/\alpha }}\underset{k=1}{\overset{\mathrm{}}{}}\frac{z^k}{\mathrm{\Gamma }(\beta \alpha k)}+𝒪\left(|z|^{p1}\right).$$
(A5)
* At $`\gamma |\mathrm{arg}z|\pi `$,
$$E_{\alpha ,\beta }(z)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{z^k}{\mathrm{\Gamma }(\beta \alpha k)}+𝒪\left(|z|^{p1}\right).$$
(A6)
## B Asymptotic behavior of the function $`F_{d,\sigma }(y)`$
To obtain the small $`y`$ behavior (41) of the function $`F_{d,\sigma }(y)`$ we use the identity
$$\mathrm{ln}\varphi =\alpha _0^{\mathrm{}}\frac{dx}{x}\left[E_{\alpha ,1}\left(x^\alpha \right)E_{\alpha ,1}\left(yx^\alpha \right)\right]$$
(B1)
and the definition of the function $`F_{d,\sigma }(y)`$:
$$F_{d,\sigma }\left(y\right)=_0^{\mathrm{}}𝑑xx^{\frac{\sigma }{2}1}E_{\frac{\sigma }{2},\frac{\sigma }{2}}\left(\frac{yx^{\sigma /2}}{(2\pi )^\sigma }\right)\left[𝒜^d(x)1\left(\frac{\pi }{x}\right)^{d/2}\right].$$
(B2)
After some algebra one obtains:
$$F_{2\sigma ,\sigma }(y)=F_{2\sigma ,\sigma }(0)+2^\sigma y𝒞_\sigma \frac{2^{1\sigma }}{\sigma \mathrm{\Gamma }(\sigma )}y\mathrm{ln}y+𝒪(y^2),$$
(B4)
where
$$𝒞_\sigma =\frac{1}{\mathrm{\Gamma }(\sigma )}_0^{\mathrm{}}\frac{du}{u}\left[E_{\frac{\sigma }{2},1}\left(\frac{u^{\sigma /2}}{(2\pi )^\sigma }\right)\frac{u^\sigma }{\pi ^\sigma }𝒜^{2\sigma }(u)+\frac{u^\sigma }{\pi ^\sigma }\right].$$
(B5)
To obtain the large $`y`$ asymptotic behavior (III C 3) of the function $`F_{d,\sigma }(y)`$ we rewrite (B2) in the form
$$F_{d,\sigma }(y)=\pi ^{d/2}_0^{\mathrm{}}𝑑xx^{\frac{\sigma }{2}\frac{d}{2}1}E_{\frac{\sigma }{2},\frac{\sigma }{2}}\left(\frac{yx^{\sigma /2}}{(2\pi )^\sigma }\right)\underset{𝒍}{}^{}e^{\pi ^2𝒍^2/x}_0^{\mathrm{}}𝑑xx^{\frac{\sigma }{2}1}E_{\frac{\sigma }{2},\frac{\sigma }{2}}\left(\frac{yx^{\sigma /2}}{(2\pi )^\sigma }\right).$$
(B6)
Using the identity
$$_0^{\mathrm{}}𝑑xx^{\frac{\sigma }{2}1}E_{\frac{\sigma }{2},\frac{\sigma }{2}}\left(x^{\sigma /2}\right)=1,\sigma >0$$
(B7)
From the second term of Eq. (B6) we obtain the first terms of Eqs. (57) and (58) respectively.
Next taking into account Eq. (A6) or Eq. (A3) for the function $`E_{\alpha ,\beta }(z)`$ and after subsequent integration in the first term of Eq. (B6), we obtain finally the asymptotic behavior given by Eqs. (57) and (58).
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# Coulomb blockade in one-dimensional arrays of high conductance tunnel junctions
## 1 Introduction
The effect of Coulomb blockade in 1D arrays of normal tunnel junctions can be used for absolute thermometry -. The properties of such arrays have been extensively investigated both experimentally and theoretically -. In all these works arrays of high Ohmic junctions with the junction resistances, $`R_j`$, higher than or of the order of the quantum resistance $`R_\mathrm{q}=h/e^225.8`$ k$`\mathrm{\Omega }`$ have been studied. In this limit a theoretical description of the Coulomb blockade is well developed and it has been successfully applied - to explain experimental findings in the high temperature regime, $`k_BT>E_Ce^2/2C`$, where $`C`$ is the capacitance of a single junction.
From the practical point of view, large ($`N1`$) 1D arrays used as thermometers are advantageous over the smaller ones, because of higher accuracy of temperature measurements. At the same time, increasing the number of junctions in the array obviously yields an increase of its total resistance $`R_{\mathrm{tot}}N`$. Since in practice it is desirable to avoid very large values of $`R_{\mathrm{tot}}`$, it appears natural to relax the condition $`R_jR_\mathrm{q}`$ and consider 1D arrays of relatively highly conducting tunnel junctions with $`R_j`$ of the order of $`R_\mathrm{q}`$ or smaller. A natural way to avoid this is the parallel connection of several 1D arrays, and this is used extensively in Coulomb blockade thermometry (CBT). To avoid a large total number of junctions, a more straightforward solution would be to decrease the resistance of each individual junction. On the other hand, heating effects turn out to be much more pronounced for highly conducting junction arrays . Hence, from this point of view, it is better not to decrease $`R_j`$ down to very low values.
The above considerations motivated us to investigate the interplay between Coulomb blockade and strong tunneling effects in 1D arrays of normal metallic tunnel junctions. Do Coulomb blockade effects survive if the junction resistance becomes smaller than the resistance quantum? Both theory - and experiment - give a clear positive answer to this question. An adequate theoretical approach which enables one to study electron transport in the strong tunneling regime is well established . This so-called quasiclassical Langevin equation technique allows to proceed analytically and remains accurate at not very low temperatures and/or voltages :
$$\mathrm{max}[k_BT,eV]\stackrel{>}{}(\mathrm{}/R_jC)\mathrm{exp}(R_\mathrm{q}/2R_j)$$
(1)
for $`R_j\stackrel{<}{}R_\mathrm{q}`$, and $`\mathrm{max}[k_BT,eV]\stackrel{>}{}e^2/2C`$ otherwise. The condition (1) implies that in the strong tunneling regime $`R_jR_\mathrm{q}`$, which is of a primary interest for us here, the technique covers practically all experimentally accessible values of temperature and bias voltage. In several previous publications -, the Langevin equation technique was applied to analyze the Coulomb blockade and strong tunneling effects in single junctions and SET transistors, where the $`IV`$ curves were derived at arbitrary tunneling strength. The results of this theoretical analysis turned out to be in good agreement with available experimental data .
The structure of the paper is as follows: In Section 2 we develop a theoretical analysis of Coulomb blockade effects in 1D tunnel junction arrays in the strong tunneling regime. In Section 3 we present experimental results obtained for the arrays with resistances in the range $`123`$ k$`\mathrm{\Omega }`$ and compare these results with our theoretical predictions. Our main conclusions are summarized in Section 4.
## 2 Theory
In order to theoretically study the behavior of low resistance tunnel junction arrays we are going to use the technique of quasiclassical Langevin equation developed in Refs. . In this section we generalize the approach presented in to the case of 1D arrays of tunnel junctions.
We consider an array of $`N`$ normal metal tunnel junctions in series (Fig. 1). The system can be described by the following Langevin equations -:
$`C_j`$ $`{\displaystyle \frac{\mathrm{}\ddot{\phi }_j}{2e}}+{\displaystyle \frac{1}{R_j}}{\displaystyle \frac{\mathrm{}\dot{\phi }_j}{2e}}=\dot{q}+\stackrel{~}{\xi }_j,j=1,\mathrm{},N;`$ (2)
$`V_\mathrm{x}`$ $`=\dot{q}R_\mathrm{S}+{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{}\dot{\phi }}{2e}}\xi _SR_\mathrm{S}.`$
Here $`\phi _j(t)\frac{2e}{\mathrm{}}_0^t𝑑t^{}V_j(t^{})`$ is the effective phase and $`R_\mathrm{S}`$ is the resistance of the electromagnetic environment. The shot noise of the $`j`$-th junction depends on $`\phi _j`$ as :
$$\stackrel{~}{\xi }_j=\xi _{1j}\mathrm{cos}(\frac{\phi _j}{2})+\xi _{2j}\mathrm{sin}(\frac{\phi _j}{2}).$$
(3)
Here $`\xi _{jk}`$ are Gaussian stochastic variables with the following pair correlators:
$`\xi _{1j}(t_1)\xi _{1j}(t_2)`$ $`=\xi _{2j}(t_1)\xi _{2j}(t_2)={\displaystyle \frac{G(t_1t_2)}{R_j}},`$ (4)
$`\xi _{1j}(t_1)\xi _{2j}(t_2)`$ $`=0,\xi _S(t_1)\xi _S(t_2)={\displaystyle \frac{G(t_1t_2)}{R_S}}.`$
In Eqs. (4) we have defined
$$G(t)\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\frac{d\omega }{2\pi }\mathrm{}\omega \mathrm{coth}(\frac{\mathrm{}\omega }{2k_BT})\mathrm{e}^{i\omega t}=\frac{1}{\pi \mathrm{}}𝒫\frac{(\pi k_BT)^2}{\mathrm{sinh}^2(\frac{\pi k_BTt}{\mathrm{}})},$$
(5)
where $`𝒫`$ stands for the principal value. There exist no correlations between the noise terms from different junctions: $`\xi _{1(2)i}(t_1)\xi _{1(2)j}=0`$ for $`ij`$.
Averaging Eqs. (2) over the noise realizations (we will denote this average by angular brackets), we obtain the general expression for the current in the array:
$`I(V_x)`$ $`=`$ $`\dot{q}={\displaystyle \frac{V}{R_\mathrm{\Sigma }}}{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{R_j}{R_\mathrm{\Sigma }}}\stackrel{~}{\xi }_j`$
$`V(V_x)`$ $`=`$ $`{\displaystyle \underset{j=1}{\overset{N}{}}}{\displaystyle \frac{\mathrm{}\dot{\phi }_j}{2e}},`$ (6)
where $`R_\mathrm{\Sigma }_{j=1}^NR_j`$ is the total resistance of the array. These two equations define the $`IV`$ curve of the chain.
The problem, now, reduces to the evaluation of the noise averages $`\stackrel{~}{\xi }_j`$. In order to do this, we first exclude the current $`\dot{q}`$ from the equations (2) and get the following ones for the phases $`\phi _j`$:
$$C_j\frac{\mathrm{}\ddot{\phi }_j}{2e}+\frac{1}{R_j}\frac{\mathrm{}\dot{\phi }_j}{2e}+\frac{1}{R_S}\underset{k=1}{\overset{N}{}}\frac{\mathrm{}\dot{\phi }_k}{2e}=\frac{V_x}{R_S}+\xi _S+\stackrel{~}{\xi }_j.$$
(7)
Then we define the small deviations of the phase and the shot noise from their average values $`\delta \phi _j\phi _j\phi _j`$ and $`\delta \stackrel{~}{\xi }_j\stackrel{~}{\xi }_j\stackrel{~}{\xi }_j`$, respectively. They obey the following equations
$$C_j\frac{\mathrm{}\delta \ddot{\phi }_j}{2e}+\frac{1}{R_j}\frac{\mathrm{}\delta \dot{\phi }_j}{2e}+\frac{1}{R_S}\underset{k=1}{\overset{N}{}}\frac{\mathrm{}\delta \dot{\phi }_k}{2e}=\xi _S+\delta \stackrel{~}{\xi }_j.$$
(8)
Applying the Fourier transformation and solving the corresponding equations, we find
$$\frac{\mathrm{}\delta \dot{\phi }_j}{2e}|_\omega =Z_j(\omega )\delta \stackrel{~}{\xi }_{j,\omega }+\underset{kj}{}a_{jk}(\omega )\delta \stackrel{~}{\xi }_{k,\omega },$$
(9)
where $`Z_j(\omega )`$ is the total impedance seen by the $`j`$-th junction and the functions $`a_{jk}(\omega )`$ describe the mutual influence of the junctions on each other. Equation (9) applies to any array of tunnel junctions of any dimensionality. In the case of 1D arrays from Eq. (8) we find
$$Z_j(\omega )=\frac{R_j}{1i\omega R_jC_j}\frac{R_S+\underset{kj}{}\frac{R_k}{1i\omega R_kC_k}}{R_S+\underset{k}{}\frac{R_k}{1i\omega R_kC_k}}.$$
(10)
Here we do not present the expressions for the functions $`a_{jk}(\omega )`$ because, as we will see below, the contribution of corresponding terms turns out to vanish. We find
$`\delta \phi _j`$ $`=`$ $`{\displaystyle \frac{2e}{\mathrm{}}}{\displaystyle \underset{\mathrm{}}{\overset{t}{}}}𝑑t^{}K_j(tt^{})\delta \stackrel{~}{\xi }_j(t^{})`$ (11)
$`+{\displaystyle \underset{kj}{}}{\displaystyle \frac{2e}{\mathrm{}}}{\displaystyle \underset{\mathrm{}}{\overset{t}{}}}𝑑t^{}A_{jk}(tt^{})\delta \stackrel{~}{\xi }_k(t^{}),`$
where the response function $`K_j(t)`$ is defined as
$$K_j(t)=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}\frac{d\omega }{2\pi }\frac{Z_j(\omega )}{i\omega +0}\mathrm{e}^{i\omega t},$$
(12)
and $`A_{jk}(t)`$ are defined analogously.
It is important to emphasize that Eqs. (9) and (11) are not the explicit solutions for the phase, but the integral equations. The variables $`\delta \stackrel{~}{\xi }_j`$ on the right hand side of these equations depend on $`\delta \phi _j`$ through the $`\mathrm{sin}(\phi _j/2)`$ and $`\mathrm{cos}(\phi _j/2)`$ terms in the shot noise (3). These integral equations can be solved by iteration. Here we restrict ourselves to the first iteration and put $`\delta \phi _j=0`$ in the right hand side of Eq. (11).
The next step is to evaluate the average values $`\stackrel{~}{\xi }_j,`$ which enter the expression for the current (6). We make the following approximation:
$`\stackrel{~}{\xi }_j=\xi _{1j}\mathrm{cos}\left[{\displaystyle \frac{eV_jt}{\mathrm{}}}+{\displaystyle \frac{\delta \phi _j}{2}}\right]+\xi _{2j}\mathrm{sin}\left[{\displaystyle \frac{eV_jt}{\mathrm{}}}+{\displaystyle \frac{\delta \phi _j}{2}}\right]`$
$`\left[\xi _{2j}\mathrm{cos}{\displaystyle \frac{eV_jt}{\mathrm{}}}\xi _{1j}\mathrm{sin}{\displaystyle \frac{eV_jt}{\mathrm{}}}\right]{\displaystyle \frac{\delta \phi _j}{2}}.`$
Here $`V_j\frac{\mathrm{}\dot{\phi }_j}{2e}`$ is the average voltage on the $`j`$-th junction. Making use of Eq. (11) we get
$$\stackrel{~}{\xi }_j=\frac{2e}{\mathrm{}R_j}\underset{0}{\overset{\mathrm{}}{}}𝑑tG(t)K_j(t)\mathrm{sin}(\frac{eV_jt}{\mathrm{}}).$$
(13)
Here we note that the terms containing the kernels $`A_{jk}(t)`$ do not contribute to the result because there exists no correlation between the noise on different junctions. Now the current is expressed as follows:
$$I=\frac{V}{R_\mathrm{\Sigma }}+\frac{2e}{\mathrm{}R_\mathrm{\Sigma }}\underset{0}{\overset{\mathrm{}}{}}𝑑tG(t)\left[\underset{j=1}{\overset{N}{}}K_j(t)\mathrm{sin}(\frac{eV_jt}{\mathrm{}})\right].$$
(14)
Let us first assume that all the junctions in the chain are identical, i.e. they have the same resistance $`RR_j`$ and the same capacitance $`CC_j`$. Then we get
$$Z_j(\omega )=\frac{R_S+(N1)\frac{R}{1i\omega RC}}{R_S\left(\frac{1}{R}i\omega C\right)+N},$$
(15)
$`K_j(t)`$ $`=`$ $`{\displaystyle \frac{N1}{N}}R\left(1\mathrm{e}^{t/RC}\right)`$ (16)
$`+{\displaystyle \frac{R_SR}{N(R_S+NR)}}\left(1\mathrm{e}^{\frac{R_S+NR}{R_SRC}t}\right).`$
The current (14), then, can be found exactly:
$$I=\frac{Tv}{eR}\frac{eT}{\pi \mathrm{}}\left[\frac{N1}{N}F(v,u)+\frac{F(v,u_S)}{N(1+N\frac{R}{R_S})}\right].$$
(17)
Here $`veV/Nk_BT`$, $`u\mathrm{}/2\pi k_BTRC`$, $`u_Su(1+NR/R_S)`$, and
$`F(v,u)`$ $``$ $`v\left[\mathrm{Re}\mathrm{\Psi }\left(1+ui{\displaystyle \frac{v}{2\pi }}\right)\mathrm{Re}\mathrm{\Psi }\left(1i{\displaystyle \frac{v}{2\pi }}\right)\right]`$ (18)
$`2\pi u\mathrm{Im}\mathrm{\Psi }\left(1+ui{\displaystyle \frac{v}{2\pi }}\right).`$
In the limit $`T0`$ the final result of the above expressions reduces to that of our previous analysis .
The differential conductance is given by the following equation
$`R_\mathrm{\Sigma }{\displaystyle \frac{dI}{dV}}`$ $`=`$ $`1{\displaystyle \frac{e^2R}{\pi \mathrm{}}}[{\displaystyle \frac{N1}{N}}{\displaystyle \frac{F(v,u)}{v}}`$ (19)
$`+{\displaystyle \frac{1}{N(1+N\frac{R}{R_S})}}{\displaystyle \frac{F(v,u_S)}{v}}].`$
Now let us put $`R_\mathrm{S}=0`$ and consider the high temperature limit $`u1`$.
Then, in the first order in $`u`$, we find
$`I`$ $`=`$ $`{\displaystyle \frac{k_BTv}{eR}}{\displaystyle \frac{N1}{N}}{\displaystyle \frac{ek_BT}{\pi \mathrm{}}}u[v\mathrm{Re}\mathrm{\Psi }^{}(1i{\displaystyle \frac{v}{2\pi }})`$ (20)
$`2\pi \mathrm{Im}\mathrm{\Psi }(1i{\displaystyle \frac{v}{2\pi }})]`$
$`=`$ $`{\displaystyle \frac{k_BTv}{eR}}{\displaystyle \frac{N1}{2N}}{\displaystyle \frac{e}{RC}}\left[\mathrm{coth}{\displaystyle \frac{v}{2}}{\displaystyle \frac{v}{2\mathrm{sinh}^2\frac{v}{2}}}\right]`$
and
$$R_\mathrm{\Sigma }\frac{dI}{dV}=1\frac{N1}{N}\frac{e^2}{Ck_BT}\frac{v\mathrm{sinh}v4\mathrm{sinh}^2\frac{v}{2}}{8\mathrm{sinh}^4\frac{v}{2}}.$$
(21)
This result exactly coincides with that found for the high Ohmic junctions .
Equation (21) is basic for the Coulomb blockade thermometry. The half-width of the dip in the $`dI/dV`$ is related to the temperature as $`V_{1/2,0}=5.439Nk_BT/e`$ . This relation is proven to be very accurate for the high Ohmic arrays. To estimate its accuracy in the strong tunneling limit, we use the expression (19), put $`R_S=0`$ and numerically solve the following equation
$$\frac{F(v_{1/2}/2,u)}{v}=\frac{1}{2}\frac{F(0,u)}{v}.$$
(22)
Here $`v_{1/2}`$ is the normalized half-width, $`v_{1/2}V_{1/2}/V_{1/2,0}`$. The solution of this equation is plotted in Fig. 2.
In the limit of high temperatures (small $`u`$) we find
$`v_{1/2}=1+0.704u0.24u^2+\mathrm{}`$
$`=1+0.112{\displaystyle \frac{\mathrm{}}{k_BTRC}}0.006{\displaystyle \frac{\mathrm{}^2}{(k_BTRC)^2}}+\mathrm{}.`$ (23)
The zero bias conductance can be obtained from Eq. (19):
$`{\displaystyle \frac{G_0}{G_\mathrm{\Sigma }}}=1{\displaystyle \frac{N1}{N}}{\displaystyle \frac{e^2R}{\pi \mathrm{}}}\left\{\mathrm{\Psi }(1+u)+\gamma +u\mathrm{\Psi }^{}(1+u)\right\}.`$ (24)
Here $`G_\mathrm{\Sigma }1/R_\mathrm{\Sigma }`$, and $`\gamma `$ is the Euler’s constant. Note that the result (24) is equivalent to the one recently derived within the framework of a linear response approach based on the Kubo formula . At high temperatures (small $`u`$) we find from (24)
$`R_\mathrm{\Sigma }G_0=1{\displaystyle \frac{N1}{N}}{\displaystyle \frac{e^2R}{\pi \mathrm{}}}\left\{{\displaystyle \frac{\pi ^2}{3}}u3\zeta (3)u^2+\mathrm{}\right\}`$
$`=1{\displaystyle \frac{N1}{N}}{\displaystyle \frac{e^2}{6Ck_BT}}\left\{10.17446{\displaystyle \frac{\mathrm{}}{k_BTRC}}+\mathrm{}\right\}.`$ (25)
Now we will study the effect of the junction asymmetry on the properties of the thermometer. If the junctions are not identical, the equations become more complicated. Here, we will consider only the high temperature limit. In this limit the function $`G(t)`$ decays at short times of order $`\mathrm{}/k_BT`$, and one can replace the response function $`K_j(t)`$ in the integral (14) by its short time expansion, which starts from the linear-in-time term, $`K_j(t)\dot{K}_j(0)t`$. Then we get
$`I={\displaystyle \frac{V}{R_\mathrm{\Sigma }}}+{\displaystyle \frac{2e}{\mathrm{}R_\mathrm{\Sigma }}}{\displaystyle \underset{j=1}{\overset{N}{}}}\dot{K}_j(0){\displaystyle \underset{0}{\overset{\mathrm{}}{}}}𝑑tG(t)t\mathrm{sin}({\displaystyle \frac{eV_jt}{\mathrm{}}})`$
$`={\displaystyle \frac{Nk_BTv}{eR_\mathrm{\Sigma }}}{\displaystyle \frac{e}{2R_\mathrm{\Sigma }}}{\displaystyle \underset{j=1}{\overset{N}{}}}\dot{K}_j(0)\left\{\mathrm{coth}{\displaystyle \frac{v_j}{2}}{\displaystyle \frac{v_j}{2\mathrm{sinh}^2\frac{v_j}{2}}}\right\}.`$ (26)
$$R_\mathrm{\Sigma }\frac{dI}{dV}=1\underset{j=1}{\overset{N}{}}\frac{e^2R_j\dot{K}_j(0)}{k_BTR_\mathrm{\Sigma }}\frac{\frac{v_j}{2}\mathrm{sinh}v_j4\mathrm{sinh}^2\frac{v_j}{2}}{8\mathrm{sinh}^4\frac{v_j}{2}}.$$
(27)
Here $`V_jVR_j/R_\mathrm{\Sigma }`$ and $`v_jeV_j/k_BT`$. From Eqs. (12) and (15) we find
$$\dot{K}_j(0)=\underset{\omega \mathrm{}}{lim}(i\omega Z_j(\omega ))=\frac{1}{C_j}\frac{1}{C_j^2\underset{k=1}{\overset{N}{}}\frac{1}{C_k}}.$$
(28)
Assuming that deviations, $`\delta R_j`$ and $`\delta C_j`$, of the junction parameters from the reference values $`R`$ and $`C`$ are small, we expand Eq. (27) in powers of $`\delta R_j`$ and $`\delta C_j`$ up to the second order. Then we get
$`R_\mathrm{\Sigma }{\displaystyle \frac{dI}{dV}}`$ $`=`$ $`1{\displaystyle \frac{N1}{N}}{\displaystyle \frac{e^2}{Ck_BT}}f(v){\displaystyle \frac{e^2}{Nk_BT}}f(v){\displaystyle \underset{j}{}}\delta \dot{K}_j(0)`$ (29)
$`{\displaystyle \frac{N1}{N}}{\displaystyle \frac{e^2}{Ck_BT}}[f(v)+vf^{}(v)]{\displaystyle \underset{j}{}}\delta r_j`$
$`(N1){\displaystyle \frac{e^2}{Ck_BT}}\left[vf^{}(v)+{\displaystyle \frac{v^2}{2}}f^{\prime \prime }(v)\right]{\displaystyle \underset{j}{}}\delta r_j^2`$
$`{\displaystyle \frac{e^2}{k_BT}}[f(v)+vf^{}(v)]{\displaystyle \underset{j}{}}\delta r_j\delta \dot{K}_j(0).`$
Here
$`f(v)`$ $``$ $`{\displaystyle \frac{\frac{v}{2}\mathrm{sinh}v4\mathrm{sinh}^2\frac{v}{2}}{8\mathrm{sinh}^4\frac{v}{2}}},`$
$`\delta r_j`$ $``$ $`\delta \left[{\displaystyle \frac{R_j}{R_\mathrm{\Sigma }}}\right]={\displaystyle \frac{N\delta R_j\underset{k}{}\delta R_k}{N^2R}}\left[1{\displaystyle \frac{\underset{k}{}\delta R_k}{NR}}\right],`$
and
$`\delta \dot{K}_j(0)`$ $`=`$ $`{\displaystyle \frac{N2}{N}}{\displaystyle \frac{\delta C_j}{C^2}}{\displaystyle \frac{\underset{k}{}\delta C_k}{N^2C^2}}+{\displaystyle \frac{N2}{N}}{\displaystyle \frac{\delta C_j^2}{C^3}}`$ (30)
$`{\displaystyle \frac{\left[\underset{k}{}\delta C_kN\delta C_j\right]^2}{N^3C^3}}+{\displaystyle \frac{\underset{k}{}\delta C_k^2}{N^2C^3}}.`$
We observe that $`\delta r_j`$ is zero if the deviations of all the resistances are equal to each other. We also note that the terms linear in $`\delta R_j`$ vanish in the sum $`_j\delta r_j`$ and, as a consequence, in Eq. (29). Both these properties reflect the fact that the half-width depends only on temperature if the resistances of all the junctions are the same. Now the correction to the half-width of the conductance dip can be obtained peturbatively in $`\delta R_j`$ and $`\delta C_j`$. In the first non-vanishing order we get
$`{\displaystyle \frac{V_{1/2}}{V_{1/2,0}}}`$ $`=`$ $`1+{\displaystyle \frac{N2}{N(N1)}}\left\{{\displaystyle \underset{j}{}}{\displaystyle \frac{\delta R_j\delta C_j}{RC}}{\displaystyle \underset{i,j}{}}{\displaystyle \frac{\delta R_i\delta C_j}{NRC}}\right\}`$ (31)
$`{\displaystyle \frac{\alpha }{N}}\left\{{\displaystyle \underset{j}{}}{\displaystyle \frac{\delta R_j^2}{R^2}}{\displaystyle \frac{1}{N}}\left[{\displaystyle \underset{j}{}}{\displaystyle \frac{\delta R_k}{R}}\right]^2\right\},`$
where $`\alpha 1+\frac{v_0}{4}\frac{f^{\prime \prime }(v_0/2)}{f^{}(v_0/2)}0.734`$.
## 3 Experiment
To obtain suitable data for comparison between the theoretical predictions (presented in the previous section) and the experiment, we fabricated high conductance Al/AlO<sub>x</sub>/Al tunnel junction arrays by electron beam lithography and two-angle shadow evaporation techniques. As a substrate, we used nitridized silicon wafers. The number of junctions in the array, $`N`$, was twenty, and each junction had an area of about 0.025 $`\mu `$m<sup>2</sup>. Different high conductance samples with per-junction asymptotic resistances of $`12`$ k$`\mathrm{\Omega }`$, together with two samples with lower conductances (in the intermediate regime) with per-junction resistances equal to 20 k$`\mathrm{\Omega }`$ and 23 k$`\mathrm{\Omega }`$, were made and measured. In order to decrease heating effects in the high conductance arrays at higher bias voltages, the islands between the junctions in most of the samples (albeit not for those shown in Fig. 3) were made sufficiently large with cooling bars (see, e.g., ) attached to them. The measurements were carried out in the temperature range 1.5 K $`\stackrel{<}{}T\stackrel{<}{}`$ 4.5 K, i.e. at temperatures of liquid helium. To measure the temperatures as accurately as possible, we fabricated CBT sensors on the same sample stage in the vicinity of the samples to be measured.
As it was already discussed above, Coulomb blockade – although weakened – is not smeared out completely even in the strong tunneling limit. The zero bias conductance of the array is always lower than its asymptotic value at high voltages. In general, the zero bias conductance at high temperatures can be written as
$$\frac{G_0}{G_\mathrm{\Sigma }}=1\frac{N1}{N}\frac{E_C}{3k_BT}+A\frac{E_C^2}{(k_BT)^2}+\mathrm{}.$$
(32)
In the limit of strong tunneling, $`R_jR_\mathrm{q}`$, from Eq. (25) we find $`A=A_{strong}\frac{N1}{N}\frac{3\zeta (3)}{2\pi ^4}\frac{R_\mathrm{q}}{R}=0.0185\frac{N1}{N}\frac{R_\mathrm{q}}{R},`$ while in the opposite weak tunneling limit the approach based on the Master equations, , yields $`A=A_{weak}\frac{1}{15}\left(\frac{N1}{N}\right)^2.`$ Here, as before, $`RR_j`$ is the (per-junction) resistance of the homogeneous array at large bias voltages.
In the intermediate regime, which is appropriate for most of the arrays measured in the experiment, one can conjecture that $`A=A_{strong}+A_{weak}`$. Previously, this conjecture was verified for the specific case of SET transistors ($`N=2`$) .
Within the first order in $`u=\mathrm{}/2\pi k_BTRC`$, the inverse resistance enhancement at zero bias voltage, $`R/\mathrm{\Delta }R`$, where $`\mathrm{\Delta }RR(V=0)R`$, can be easily derived from Eq. (32) as:
$$\frac{R}{\mathrm{\Delta }R}=3\frac{N}{N1}\frac{k_B}{E_C}T+a\frac{R_\mathrm{q}}{R}+b,$$
(33)
where, according to the theory,
$$a=\frac{N}{N1}\frac{27\zeta (3)}{2\pi ^4}=0.175$$
(34)
and
$`b=\{\begin{array}{cc}1,\hfill & A=A_{strong};\hfill \\ 2/5,\hfill & A=A_{strong}+A_{weak}.\hfill \end{array}`$ (35)
The expression above has two characteristic features: linearity in $`T`$ and dependence of its slope on the capacitance of the junctions in the array. In addition, this equation predicts an offset which depends only on the number of junctions in the array, and on the ratio of the quantum and per-junction resistances. Below, while making comparison between the measured data and the predictions of Eq. (33), we will take the the weak tunneling correction into account, i.e, $`A=A_{strong}+A_{weak}`$.
The value $`R/\mathrm{\Delta }R`$, measured for two tunnel junction arrays with $`N=20`$ and at different temperatures, is displayed in Fig. 3. The asymptotic resistances of the samples were 23 k$`\mathrm{\Omega }`$ (open circles) and 1.2 k$`\mathrm{\Omega }`$ (solid triangles). One observes an almost perfect linear dependence of the value $`R/\mathrm{\Delta }R`$ on temperature (cf. Eq. (33). This dependence can be used to obtain a quantitative estimate for the junction capacitance. By fitting the slope of the experimental curves of Fig. 3 to Eq. (33) we find $`C=2.4`$ fF and 2.1 fF respectively for the arrays with $`R=23`$ and 1.2 K$`\mathrm{\Omega }`$. This way appears to be the most reliable to evaluate the per-junction capacitance in arrays of normal metal tunnel junctions .
The second and third terms in the right hand side of Eq. (33) yield an offset which – for a given number of junctions $`N`$ – should depend solely on the ratio $`R_\mathrm{q}/R`$. Eq. (33) gives the values 3.4 and -0.2 for this offset respectively for the samples represented by the solid triangles and the open circles. The corresponding numbers obtained from the lines fitted to the experimental data in Fig. 3 are 4.1 and 0.2, being in a reasonable agreement with the above theoretical values.
The offset values for two additional samples with (per-junction) resistances $`R=2.1`$ k$`\mathrm{\Omega }`$ and $`R=2.2`$ k$`\mathrm{\Omega }`$ (the samples for which the systematic $`T`$-dependence was measured) were equal to 3.6 and 2.9, respectively. The corresponding theoretical predictions are 1.8 and 1.2. The data for the offsets are presented in Fig. 4. We observe that arrays with lower resistances show larger offsets, as predicted by our theory, Eq. (33). However, the offset values presented as a function of $`R_\mathrm{q}/R`$ do not exactly fall on the straight line: they are scattered within an interval of $`\pm 0.5`$. This effect could probably be attributed to the inhomogeneity of the arrays. For instance, the measured offsets for two arrays with nearly identical values of $`R=2.1`$ k$`\mathrm{\Omega }`$ and 2.2 k$`\mathrm{\Omega }`$ differ from each other. This discrepancy can be explained, if one assumes a certain degree of junction asymmetry, e.g., a 20 $`\%`$ fluctuation of the per-junction resistance, $`R_j`$, around the mean value, $`R`$. This, in turn, can arise from a 40 % fluctuation in the area of the ”odd” and the ”even” tunnel junctions as a consequence of the two-angle evaporation technique employed in the sample fabrication.
The best linear fit for the offset – inverse resistance dependence (Fig. 4) yields $`a=0.19\pm 0.05`$, $`b=0.4\pm 0.7`$. The value $`a`$ agrees well with our theoretical prediction (34). At the same time the experimental value of $`b`$ turns out to be different from the predictions of the theory (35). It is also too uncertain due to the scattering of the data points. At present, possible reasons for this discrepancy remain unclear. We have checked asymmetry effects as well as those of the external environment within the framework of the Eq. (19). These effects can hardly help to improve the agreement between theoretical and experimental values of $`b`$.
In the weak tunneling regime $`RR_\mathrm{q}`$ the offset is given by $`b`$ because $`aR_\mathrm{q}/R0`$. In order to check the effect of finite external resistance on the offset in this limit we have performed Monte-Carlo simulations (for details see ) based on the phase-correlation theory . The only fitting parameter used in these simulations was the junction capacitance. By comparing the measured conductance curves of the array represented by the solid triangles in Fig. 3 ($`R=2.1`$ k$`\mathrm{\Omega }`$) to those obtained from the numerical simulations, we have obtained the value $`C=2.1`$ fF for this sample. Notice that this value is in excellent agreement with that derived from Eq. (33) as explained above. We have also found that, depending on the magnitude of the environmental resistance, $`b`$ may increase by $`0.1`$ at most. Again this is not sufficient to explain the observed values of $`b`$.
Figure 5 shows the variation of the normalized half-width $`\mathrm{\Delta }V_{1/2}/V_{1/2,0}`$ as a function of temperature ($`\mathrm{\Delta }V_{1/2}V_{1/2}V_{1/2,0})`$. The solid and open squares represent the data measured for two 20-junction arrays both with $`R=2.1`$ k$`\mathrm{\Omega }`$. The solid curve corresponds to our theoretical prediction, Eq. (23). We observe, that the difference between the experimental points and our theoretical curve typically does not exceed one percent, i.e. the ageement is fairly good.
Within the framework of our analysis, one can easily establish yet one more useful relation between the normalized half-width, $`\mathrm{\Delta }V_{1/2}/V_{1/2,0}`$, and the value of the conductance dip $`\mathrm{\Delta }GG_\mathrm{\Sigma }G(V=0)`$ \[cf, e.g., Eq. (19)\] evaluated in the linear regime. Combining Eqs. (23) and (25) in the high temperature limit one gets
$$\frac{\mathrm{\Delta }V_{1/2}}{V_{1/2,0}}=\chi (R_\mathrm{q}/R)\frac{\mathrm{\Delta }G}{G_\mathrm{\Sigma }},$$
(36)
where the function $`\chi (R_\mathrm{q}/R)`$ in the strong tunneling limit $`R_\mathrm{q}R`$ reads
$$\chi (R_\mathrm{q}/R)0.108\frac{N}{N1}\frac{R_\mathrm{q}}{R}.$$
(37)
Note, that Eq. (36) was previously derived in Ref. in the opposite weak tunneling limit $`RR_\mathrm{q}`$, in which case the $`\chi `$-function tends to the constant
$$\chi (R_\mathrm{q}/R)0.392.$$
(38)
We assume that in the intermediate regime the function $`\chi (R_\mathrm{q}/R)`$ is the sum of the expressions (37) and (38). This assumption, although not proven rigorously, provides a reasonable interpolation between the two limiting cases.
In addition to the arrays discussed above, three other samples with per-junction resistances of 20 k$`\mathrm{\Omega }`$, 1.4 k$`\mathrm{\Omega }`$, and 1.0 k$`\mathrm{\Omega }`$ were measured at $`T4.2`$ K. The exact temperature was obtained from the vapor pressure of liquid helium. The corresponding data are collected into Table 1. The second row in the table represents the measured half-widths obtained directly from the measurement, $`V_{1/2,\mathrm{meas}}`$. The corresponding conductance dips, $`\mathrm{\Delta }G/G_\mathrm{\Sigma }`$, are shown in the third row. The forth row demonstrates the ”corrected” half-widths (i.e. the weak tunneling correction \[Eqs. (36) and (38) subtracted\],
$$V_{1/2,\mathrm{corr}.}=V_{1/2,\mathrm{meas}.}[10.392\frac{\mathrm{\Delta }G}{G_\mathrm{\Sigma }}].$$
(39)
In the weak tunneling regime this value would coincide with $`V_{1/2,0}=5.439Nk_BT/e.`$ However in our experiments this is not the case. The fifth row shows the relative deviations of the corrected half-widths from that of the basic linear result, $`\mathrm{\Delta }V_{1/2,0}^{\mathrm{corr}}=(V_{1/2,\mathrm{corr}}V_{1/2,0})/V_{1/2,0}`$. In fact, the numbers in this row can be interpreted as the ”residual inaccuracies” of the measured half-widths as compared to the weak tunneling approximation. They can be explained by the strong tunneling effects. The corresponding theoretical values, obtained from the strong tunneling correction to the half-width of the conductance dip around zero bias voltage \[Eq. (36, 37)\], are shown in the last row of the table. These predictions are in a good agreement with the corrected measured data discussed above.
## 4 Conclusions
We have studied one-dimensional arrays of tunnel junctions in the strong tunneling regime, both theoretically and experimentally. Within the framework of the quasiclassical Langevin equation formalism, analytical expressions for the current-voltage characteristics of such arrays, together with the expressions for the half-width of the conductance dip around zero bias voltage, have been derived. Furthermore, the effect of external elecromagnetic environment has been studied theoretically. We have fabricated and measured several arrays of tunnel junctions with per-junction resistances ranging from 1 k$`\mathrm{\Omega }`$ to 23 k$`\mathrm{\Omega }`$. Our measured data are in rather good agreement with the theoretical predictions. The experiments demonstrate that Coulomb blockade effects survive even in the strong tunneling regime, and are clearly visible in arrays with per-junction tunnel resistances as low as 1 k$`\mathrm{\Omega }`$. These observations are in agreement with other recent experimental results on SET transistors and on single junctions -.
It has been verified, both in the theory and in the experiment, that high conductance arrays of tunnel junction are less favorable for Coulomb Blockade Thermometry (CBT) applications. This is a combined consequence of the more pronounced heating effects in such arrays, and the larger departure from the simple linear relation $`V_{1/2,0}=5.439Nk_BT/e`$, which is of central role in the CBT applications. The latter, however, is not such a severe limitation, since corrections to the linear relation are now well known for all values of the junction resistance. It has also been shown that in the high temperature limit (in the leading approximation in $`1/T`$) the results for the conductance dip, now derived for arbitrary tunneling strength, coincide with those previously obtained in the weak tunneling regime.
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# The afterglow of the short/intermediate-duration gamma-ray burst GRB~000301C: A jet at z=2.04Based on observations made with the Nordic Optical Telescope, operated on the island of La Palma jointly by Denmark, Finland, Iceland, Norway, and Sweden, in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofisica de Canarias.,Based on observations collected at the European Southern Observatory, La Silla and Paranal, Chile (ESO project No. 64.H-0573) ,Based on observations at the German-Spanish Astronomical Centre, Calar Alto, operated by the Max-Planck-Institute for Astronomy, Heidelberg, jointly with the Spanish National Commission for Astronomy.
## 1 Introduction
The discovery of the first $`X`$-ray afterglow (Costa et al. C1997 (1997)) and optical counterpart (van Paradijs et al. PGG1997 (1997)) to a long-duration gamma-ray burst (GRB) have led to a revolution in GRB research. The determination of a redshift of 0.835 for GRB~970508 (Metzger et al. MDK1997 (1997)), and the subsequent determination of redshifts of 13 bursts with a median redshift of $``$1.0, have firmly established their cosmological origin (Kulkarni et al. 2000a ; This work; Bloom et al. BDD2000 (2000)).
The intriguing case of an association of the peculiar supernova SN1998bw with GRB~980425 (Galama et al. GVP1998 (1998)) was the first indication of a possible connection with supernovae. Evidence for supernova signatures in the late light curves of GRB~970228 (Reichart Rei1999 (1999); Galama et al. GTV1999 (1999)) and GRB~980326 (Castro-Tirado & Gorosabel CTG1999 (1999); Bloom et al. BKD1999 (1999)) suggests that at least some long-duration GRBs may be related to the collapse of massive ($`>25`$ M) stars. Breaks in the power-law declines of GRB~990123 (Kulkarni et al. KDO1999 (1999)) and GRB~990510 (Harrison et al. HBF1999 (1999)) are interpreted as evidence for collimated outflows (‘jets’) (see also Holland et al. HBH2000 (2000)). Further evidence for this collapsar + jet model (e.g., MacFadyen & Woosley MW1999 (1999)) comes from the light curve of GRB~980519 which is best interpreted as a jet expanding into a preexisting circumburst stellar wind (Jaunsen et al. J2001 (2001)).
The high-energy properties of GRBs show a bi-modal distribution of burst durations (Kouveliotou et al. K1995 (1995)) which, in the simplest scenario, may indicate the existence of binary compact mergers as the progenitors of the short-duration bursts (T$`{}_{90}{}^{}<2`$ s). From an analysis of the Third BATSE Catalog, Mukherjee et al. (MFB1998 (1998)) have shown that, in addition to the short (T$`{}_{90}{}^{}<2`$ s) and long (T$`{}_{90}{}^{}>5`$ s) classes, there may exist a third, intermediate soft-spectrum class of GRBs with duration 2 $`\mathrm{s}<\mathrm{T}_{90}<5`$ s.
In this paper we report the discovery and subsequent observations and analysis of the afterglow of the short-to-intermediate duration GRB~000301C (Fynbo et al. 2000a ).
Sect. 2 reports the detection, IPN localisation and the high-energy data of the GRB obtained from Ulysses and NEAR. Sect. 3 describes the discovery of the optical counterpart and our subsequent optical and infrared observations. Sect. 4 details the optical and infrared photometry and Sect. 5 describes the VLT spectroscopy. Sect. 6 describes the results obtained on the spectroscopy and spectral energy distribution and Sect. 7 is devoted to the discussion and interpretation, with Sect. 8 presenting our conclusions.
Throughout this paper, we adopt a Hubble constant of H<sub>0</sub> = 65 km s<sup>-1</sup> Mpc<sup>-1</sup> and assume $`\mathrm{\Omega }_\mathrm{m}=0.3`$ and $`\mathrm{\Omega }_\mathrm{\Lambda }=0.7`$.
## 2 Detection and localisation of the gamma-ray burst
GRB~000301C was recorded by the Ulysses GRB experiment and by the NEAR $`X`$-Ray/Gamma-Ray Spectrometer. Because this burst was relatively weak, it did not trigger the Ulysses Burst Mode, and the only data available from Ulysses is the Observation Mode 1 0.25 s resolution 25–150 keV light curve (Hurley et al. 1992b ). NEAR records the light curves of bursts in the 150–1000 keV energy range with 1 second resolution, but takes high-energy spectra only with 40 min resolution.
Analysis of the Ulysses and NEAR relative timing data yields an annulus centred at $`(\alpha ,\delta )_{2000}=(20^\mathrm{h}34^\mathrm{m}7.56^\mathrm{s},+20^{}32\mathrm{}19.62\mathrm{})`$, with a radius of 57.520$`\pm 0.083`$ degrees (at 3$`\sigma `$ full-width). This annulus intersected the error-box of the All-Sky Monitor (ASM) on the RXTE spacecraft, at near-right angles to create a composite localisation of a parallelogram of area $`50`$ arcmin<sup>2</sup> (see Fig. 2).
Since no high-energy spectra are available, we have estimated the peak fluxes and fluences for trial power-law spectra with indices between 1 and 4 using the Ulysses data. For a typical power-law index of 2, we find a 25–100 keV fluence of $`2.1\times 10^6\mathrm{erg}\mathrm{cm}^2`$, and a peak flux over the same energy range, and over 0.25 s, of $`6.3\times 10^7\mathrm{erg}\mathrm{cm}^2s^1`$. The uncertainties in these numbers are partly due to the lack of a high-energy spectrum. For example, the fluence estimates range from $`1.45\times 10^6`$ to $`2.24\times 10^6\mathrm{erg}\mathrm{cm}^2`$ as the spectral index is varied from 4 to 1. The statistical uncertainty is approximately 30%. From the NEAR data we estimate the 150–1000 keV fluence to be approximately $`2\times 10^6\mathrm{erg}\mathrm{cm}^2`$.
To date, the only GRBs with identified long-wavelength counterparts have been long-duration bursts. As measured by both Ulysses and NEAR, in the $`>`$25 keV energy range, the duration of this burst was approximately 2 s. (Note that the earlier estimate of a 10 s duration of GRB~000301C by Smith et al. (SHC2000 (2000)) was based on the $`<`$10 keV energy range). Thus it falls in the short class of bursts, though it is consistent with belonging to the proposed intermediate class or the extreme short end of the distribution of long-duration GRBs (Hurley et al. 1992a ; Mukherjee et al. MFB1998 (1998)). Although we do not have any measurements of the high-energy spectra above 25 keV, it is possible to derive a crude estimate of the spectral index, and therefore the hardness ratio (the 100–300 keV fluence divided by the 50–100 keV fluence), from the Ulysses and NEAR count rates. We obtain a hardness ratio of 2.7$`\pm `$0.6(cutoff)$`\pm `$30%(statistical error) from fitting a power-law, with the index as a free parameter, to the count rates from NEAR and Ulysses, assuming a range of cut-off energies.
Fig. 1 shows the location of GRB~000301C in a hardness vs. duration plot. The contour plot contains 1959 GRBs for which data on fluence and duration were available in the Fourth BATSE GRB Catalog (revised) (Paciesas et al. PMP1999 (1999)) and the BATSE Current GRB Catalog<sup>1</sup><sup>1</sup>1Data on current GRBs are available through the BATSE homepage http://www.batse.msfc.nasa.gov/batse/. The symbols represent the 10 GRBs included in the BATSE catalogs for which an afterglow has been identified, with GRB~000301C located near the center of the plot. Triangles are bursts where a break has been found in the optical light curve. From this sparse set of data there does not appear to be any marked difference in the distributions of bursts with, or without, an identified break.
Of the 1959 BATSE bursts in Fig. 1, the ratio between bursts with a duration of T$`{}_{90}{}^{}2.0`$ s and with T$`{}_{90}{}^{}<2.0`$ s is 3:1. To date, at least 23 GRB optical afterglows have been discovered (Kulkarni et al. 2000a ; Andersen et al. AH2000 (2000); this work; Klose et al. KSF2000 (2000); Fynbo et al. 2001a ; Fynbo et al. 2001b ; Henden H2001 (2001); grb-webpage of J. Greiner<sup>2</sup><sup>2</sup>2http://www.aip.de/~jcg/grbgen.html). If the distribution of the 23 GRBs with identified counterparts follows the general BATSE distribution, one would expect that $`17\pm 4`$ bursts were in the long class, and $`6\pm 2`$ bursts were in the short class. However, GRB~000301C is the only GRB with a duration consistent with the short-duration class. The expected number of identified short burst counterparts is moderated by the strong selection bias caused by the technical difficulties of obtaining precise localisations for the short GRBs.
## 3 Discovery and observations of the afterglow
The IPN/RXTE error-box of GRB~000301C (Smith et al. SHC2000 (2000)) was observed with the 2.56-m Nordic Optical Telescope (NOT) on 2000 March 3.14–3.28 UT ($``$1.8 days after the burst) using the Andalucía Faint Object Spectrograph (ALFOSC). Comparing with red and blue Palomar Optical Sky Survey II exposures, a candidate Optical Transient (OT) was found at the position $`(\alpha ,\delta )_{2000}=(16^\mathrm{h}20^\mathrm{m}18.56^\mathrm{s},+29^{}26\mathrm{}36.1\mathrm{})`$ (Fynbo et al. 2000a ). A finding-chart of the IPN-errorbox and the two ALFOSC pointings used to cover the field can be seen in Fig. 2.
The transient nature of the candidate was subsequently confirmed at optical and infrared wavelengths (Bernabei et al. BMB2000 (2000); Stecklum et al. SKF2000 (2000); Garnavich et al. GBJ2000 (2000); Veillet et al. V2000 (2000); Fynbo et al. 2000b ; Kobayashi et al. KGT2000 (2000)). At the time of discovery the magnitude of the OT was R=20.09$`\pm `$0.04 (see Sect. 4 for a detailed discussion of the photometry and Fig. 3 for a finding chart and UBRI-images of the OT on 2000 March 3 UT).
We obtained subsequent optical observations using the NOT, the Antu telescope (UT1) of ESO’s Very Large Telescope (VLT), the USNOFS 1.0-m telescope, the 2.2-m telescope at Calar Alto (CAHA) and the Wide Field Imager (WFI) on the ESO 2.2-m telescope. In addition we obtained near-infrared (NIR) data from the United Kingdom Infra-Red Telescope (UKIRT).
The journal of our optical and NIR observations, including the derived magnitudes, is given in Table LABEL:TABLE:photometry.
The OT was also observed from several other optical and infrared telescopes, and the counterpart was subsequently detected in the mm-band at 250 GHz (Bertoldi B2000 (2000)), and at radio (8.46 GHz) wavelengths (Berger & Frail BF2000 (2000))<sup>3</sup><sup>3</sup>3See the GCN Circulars Archive for further details (http://gcn.gsfc.nasa.gov/gcn/gcn3\_archive.html). Papers detailing the properties of the counterpart of GRB~000301C have been presented at radio and mm wavelengths (Berger et al. BSF2000 (2000)) and in the infrared (Rhoads & Fruchter RF2001 (2001)), and optical bands (Masetti et al. MBB2000 (2000); Sagar et al. SMP2000 (2000)).
## 4 Photometry
### 4.1 Optical data
To avoid contamination from the nearby star A (located at a separation of $`6\mathrm{}`$ west and $`1\mathrm{}`$ south of the OT in Fig. 3), we measured the magnitude of the OT relative to stars in the field by performing Point Spread Function (PSF) photometry, using DAOPHOT II (Stetson S1987 (1987)S1997 (1997)). There are several bright and unsaturated stars in the field from which a good PSF could be determined.
For the data presented here, there is no indication of a contribution from a host galaxy to the emission at the position of the OT (see Sect. 7.2 for a discussion of the host galaxy). Hence, extended emission from a faint galaxy at the position of the OT will not affect the PSF photometry appreciably (much less than observational errors). The quality of the PSF photometry was checked by subtracting the PSFs from the images of star A and the OT. In all frames the residuals are consistent with being shot-noise.
To avoid errors due to colour terms or colour differences in our photometry (the conditions during most of the observations at the NOT were possibly non-photometric due to increasing amounts of Saharan dust in the atmosphere over the telescope at the time of observations), the magnitudes of the OT for all our optical photometry were calibrated relative to stars of similar colours in the field.
The photometric standard UBVR<sub>C</sub>I<sub>C</sub> calibration of the field was performed at the USNOFS 1.0-m telescope and is available in Henden (H2000 (2000)). This calibration has an estimated zero-point uncertainty of 2 percent, which is well below the errors in the relative magnitudes. The results of the PSF photometry are presented in Table LABEL:TABLE:photometry. The 2000 March 6.39 VLT R-data point has been derived from the March 6 combined VLT-spectrum (Table 2).
Based on this photometric calibration we conclude that star A showed no sign of variability within observational errors throughout our observations, and that it had the following magnitudes: U = 20.427$`\pm `$0.133, B = 19.837$`\pm `$0.030, V = 18.767$`\pm `$0.018, R = 18.084$`\pm `$0.043, and I = 17.526$`\pm `$0.044.
### 4.2 Near-infrared data
The UKIRT images were processed using the ORAC imaging data reduction routines developed for UKIRT (Bridger et al. BWE2000 (2000)). The J, H and K magnitudes of the OT were then measured from the UKIRT data as follows. First we measured the magnitude of the OT relative to star A using DAOPHOT II PSF photometry as described above. Then, in order to transform this magnitude to the standard UKIRT system, we performed aperture photometry in an aperture with a diameter of $`2\stackrel{}{.}7`$ on calibration images obtained of the standard stars S868-G and p389-d from the list of UKIRT faint standards<sup>4</sup><sup>4</sup>4http://www.jach.hawaii.edu/JACpublic/UKIRT/ astronomy/calib/faint\_stds.html and on star A. The estimated error in the zero-point is about 0.05 in each of J, H and K. We have assumed negligible extinction difference between standard and program field. The results of these measurements are presented in Table LABEL:TABLE:photometry.
## 5 Spectroscopy
Spectroscopic observations were carried out on 2000 March 5 and 6 UT with VLT-Antu equipped with FORS1 (for details see the observing log in Table 2). We used the GRIS\_300V+10 grism and the GG375 order separation filter, which provide a spectral coverage from 3600 Å to 8220 Å and a dispersion of 2.64 Å/pixel. The effective exposure time was 800 s on March 5.39 UT and 1200 s on March 6.38 UT. Standard procedures were used for bias and flat field correction, and the optimal extraction procedure for faint spectra (described in Møller (M2000 (2000))) was used to extract one dimensional spectra.
The position angle of the long slit was chosen such that both the OT and star A were centred onto the slit. From the magnitude of star A we calibrated the flux of the trace of the optical counterpart. The spectral flux calibration derived for the OT on March 5 is consistent with the optical photometry displayed in Table LABEL:TABLE:photometry. We derive a value for the spectral index of $`\beta =1.15\pm 0.26`$ on March 5 and $`\beta =1.43\pm 0.28`$ on March 6, corrected for interstellar extinction of $`\mathrm{E}(\mathrm{B}\mathrm{V})=0.053\pm 0.020`$, using the dust maps of Schlegel et al. (SFD1998 (1998)).
## 6 Results
### 6.1 GRB absorption lines and redshift determination
The combined spectrum, reproduced in Fig 3, has a resolution of 14 Å FWHM, and signal–to–noise (S/N) in the range 15–30 per resolution element redwards of 4000 Å. From 4000 Å to 3600 Å the S/N drops rapidly.
Due to the poor resolution, only very strong absorption lines can be detected individually. In Table 3 we list the only four absorption features which were detected at a S/N in excess of 4.5. Two of the features were found bluewards of 4000 Å, and were initially ignored. The line at 4712 Å is broader than the resolution profile, and we tentatively identified it as a possible C iv absorption complex with redshifts in the range 2.038 to 2.042. With this identification, the other three features would fit the proposed identifications given in Table 3. Note, however, that the Si ii 1260 line is far too strong and wide to be a single line, and we hence assume that it is a blended feature. The Fe ii 2600 line is strong and narrow, but also this line seems excessively strong given the lack of other strong Fe ii lines.
In order to provide a more strict test of our proposed identification, and to obtain an accurate value for the redshift, we proceeded as follows. First we shifted and stacked pieces of the spectra where we would expect common low ionization lines. We selected the singly ionized species of Si, C and Fe, all of which are known as strong absorbers in quasar absorption line systems. In total our spectrum covers positions of the lines Fe ii 1608, Fe ii 2344, Fe ii 2374, Fe ii 2382, Fe ii 2586, Fe ii 2600, Si ii 1260, Si ii 1304, Si ii 1526, and C ii 1334. Si ii 1260 (at 3832 Å) was in the very low S/N part of spectrum, and almost certainly blended, so it was not included. Treating each ion separately, a weighted mean absorption feature was calculated using the oscillator strength of each line as statistical weight.
The regions of the co–added spectra for each of the ions Fe ii, Si ii and C ii, transformed into redshift space, is shown in the left panel of Fig. 5. A combined “Low Ionization” absorption feature (bottom of left panel of Fig. 5) was obtained by co–addition of the three sets of features, but using the number of lines as statistical weights. The redshift range searched for low ionization absorption systems by this method was 1.95–2.14 and no other candidate systems were found in this range. For comparison of the redshifts, we plot in the right panel of Fig. 5 again the combined “Low Ionization” absorption feature together with the C iv absorption trough. The bottom panel here shows the combination of all the lines. Given the significance of this combined “Metal Absorption Feature”, we conclude that the tentative identification of this system is confirmed.
It is commonly seen in quasar absorption line systems that the low ionization species, tracing the cold dense gas, have a more well-defined redshift than the high ionization species. Hence, we shall adopt the redshift $`z_{\mathrm{abs}}=2.0404\pm 0.0008`$, measured from the combined Low Ionization feature, as the systemic redshift. It is seen from the bottom right panel of Fig. 5 that inclusion of C iv would result in a slightly higher redshift. Note that $`z_{\mathrm{abs}}=2.0404\pm 0.0008`$ is consistent with the redshift z=1.95$`\pm 0.1`$ based on the Lyman break (Smette et al. SFG2000 (2000); Feng et al. FWW2000 (2000)), but is significantly higher than the value z=2.0335$`\pm `$0.0003 reported by Castro et al. (CDD2000 (2000)).
The oscillator strength weighted mean observed equivalent width of the Fe ii lines is 2.56 Å, which is strong enough that by comparison to known quasar absorbers one would expect this to likely have a column density of neutral Hydrogen in excess of $`2\times 10^{20}\mathrm{cm}^2`$. Such absorbers are known as Damped Ly$`\alpha `$ Absorbers (DLAs), and hold a special interest because of the large amounts of cold gas locked up in those objects (Wolfe et al. WLF1995 (1995); Storrie-Lombardi et al. SLI1997 (1997)). It is commonly assumed that the DLAs are the progenitors of present day disk galaxies, but they have proven extremely difficult to identify (see e.g. Møller & Warren 1993; Kulkarni et al. 2000b ). Observational evidence has been accumulating (Møller & Warren 1998; Fynbo et al. FMW1999 (1999)) which suggests that a likely reason why DLA galaxies are so hard to identify is their small gas cross–sections and faint magnitudes, causing them to stay hidden under the point spread functions of the bright quasars. A GRB selected DLA galaxy sample would not be hampered by this problem once the OT has faded sufficiently, and could as such help greatly in understanding the nature of DLA galaxies. We shall therefore now briefly consider the low S/N part of the GRB spectrum below 4000 Å, to investigate if any information concerning the H i column density can be extracted.
In Fig. 6 (lower panel) we have plotted the spectral region around Ly$`\alpha `$, and for comparison of redshifts, the “Combined Metals” feature (upper panel). Also plotted on the lower panel is the noise per bin (for redshift binsize = 0.003). It is clearly seen that the spectrum drops steeply before the expected central position of the Ly$`\alpha `$ line, and well before the S/N drops below detection. One likely explanation for this is the presence of a very broad Ly$`\alpha `$ absorption line. To quantify this we have modelled several Ly$`\alpha `$ absorption lines, all at redshift 2.0404, and calculated the $`\chi ^2`$ of their fit to the data in the range 3700 Å to 3750 Å. For N$`(H\text{i})=0`$ the $`\chi ^2`$ per degree-of-freedom DOF is 6.46 which confirms that an absorption feature is indeed present. The formal $`\chi ^2`$ minimum is found at N$`(H\text{i})=1.5\times 10^{21}\mathrm{cm}^2`$ ($`\chi ^2`$ per DOF = 0.86), but any value within a factor 3 of this is acceptable.
It should be recalled that the above estimate, in a strict sense, only applies in the case the OT lies “behind” the absorbing cloud. In case the OT is in fact embedded “inside” a DLA cloud, resonant scattering of Ly$`\alpha `$ photons may alter the profile of the absorption line somewhat. In the present case this is a detail which the quality of our data will not allow us to discern.
### 6.2 The multi-wavelength spectrum around March 4.5 UT
The wide wavelength coverage obtained around 2000 March 4.5 UT allows us to construct the multi-wavelength spectrum of the afterglow, by using our USNO and UKIRT data. Most fireball models (e.g. Sari et al. SPN1998 (1998); Piran P1999 (1999); Mészáros M1999 (1999) and references therein) and observations of previous afterglows suggest a power-law Spectral Energy Distribution (SED). However, a global fit to the broadband SED of GRB~000301C by Berger et al. (BSF2000 (2000)) demonstrates that the mm–optical range cannot be described by a single power-law. Thus, we have only considered wavelengths shorter than IR when fitting the SEDs presented in this section. This gives eight measurements in the period March 4.39–4.55 UT, namely UBVRIJHK. The eight measurements, plotted against $`\mathrm{log}\nu `$ together with the VLT spectrum in Fig. 7, allow us to constrain the SED of the OT at this epoch.
As extinction is highly wavelength dependent, a progressive deviation from a pure power-law fit can be explained as being due to dust extinction in the GRB host galaxy. In order to test this possibility, we first fit a power-law to the data in the NIR–optical range with the addition of an intrinsic extinction law, i.e., using the expression $`F_\nu \nu ^\beta \times 10^{(0.4A_\nu )}`$, where A<sub>ν</sub> is the rest frame extinction in magnitudes at the rest frame frequency $`\nu `$. Four extinction laws have been applied in order to establish a relationship between A<sub>ν</sub> and $`\nu `$. Following Rhoads & Fruchter (RF2001 (2001)), we have fitted the extinction curves for the Galaxy and Magellanic Clouds published by Pei (P1992 (1992)). Pei (P1992 (1992)) provides extinction laws for the Milky Way (MW), the Large Magellanic Cloud (LMC), and the Small Magellanic Cloud (SMC) based on their different proportions in the dust-to-gas ratio (1:1/5:1/10) and in the abundance of heavy elements (1:1/3:1/8). The most significant difference is the sequential change in the strength of the 2175 Å extinction feature, prominent for the MW, moderate in the LMC, and nonexistent for the SMC extinction curve. It is important to note that for the redshift of GRB~000301C this extinction feature falls in the observed R-band. So, the presence of this feature would result in a clear decrease in the R-band flux compared to the I and V-bands. The MW extinction curve requires an equal amount of graphite and silicate grains, while the SMC extinction curve can be explained by silicate grains only, with the LMC extinction law as an intermediate stage.
We have applied these four extinction laws to our data in order to infer qualitative information about the dust-to-gas ratio, the abundance of heavy elements and the composition of the dust in the host galaxy of GRB~000301C. We leave A<sub>ν</sub> as a free parameter, so fitting a function like $`F_\nu \nu ^\beta \times 10^{(0.4A_\nu )}`$ allows us to determine $`\beta `$ and A<sub>ν</sub> simultaneously. The values of $`\chi ^2`$ are displayed in Table 4 for each case.
A pure power-law fit ($`F_\nu \nu ^\beta `$) to the eight data points (after correction for Galactic extinction) leads to a reduced $`\chi ^2`$ of 1.69 (see Table 4), making an acceptable description of the NIR–optical range of the SED. However, the fit can be improved if a modest amount of extinction is introduced. This is because the SED is slightly bending down towards higher frequencies (see Fig. 7). The eight data points show that there is not any presence of a redshifted 2175 Å absorption bump in the R-band at all. In short, the near-IR SED of GRB~000301C can be described as a curved power-law but with no broad absorption features.
As expected by the lack of the absorption bump in the R-band, the MW and the LMC extinction laws are completely inconsistent with our data. In fact, both fits imply an unphysical negative extinction (see Table 4). This is because the R-band flux is slightly over the linear interpolation between the I-band and V-band fluxes (in a Log-Log space), and both of these two extinction laws fit the 2175 Å bump as an emission feature instead of an absorption bump. To illustrate the problem with the MW and LMC extinction curves, we have, in the lower panel of Fig. 7, plotted the effect of having the MW and LMC extinction laws with the parameters derived for the SMC extinction law ($`\beta `$=-0.7, A<sub>V</sub>=0.09). As seen, the shapes of both SEDs are incompatible with our UBVRIJHK measurements.
The quality of the MW, LMC and SMC and unextincted SEDs can also be compared checking the flux predicted at 250 $`\times `$ (1+z) GHz (rest-frame), where Berger et al. (BSF2000 (2000)) reports a flux of 2.1 $`\pm `$ 0.3 mJy at 250 GHz on March 4.29 UT. Extrapolating the four extinction curves, we obtain the following fluxes in increasing order; 5.0 mJy (SMC), 18.1 mJy (No extinction), 33.9 mJy (MW) and 51.6 mJy (LMC). As with the NIR–optical range, the SMC extinction provides the most reasonable results. The actual measured flux in mm is below the value predicted by the mm–optical extrapolation, because the pure power-law assumption is not correct in the mm-NIR spectral range and an additional curvature effect is present in the SED, as demonstrated by Berger et al. (BSF2000 (2000)) (see their Fig. 1). Thus, the value of A$`{}_{V}{}^{}=0.09\pm 0.04`$ for the SMC-fit given in Table 4 should be taken as a good indication of the real extinction, although strictly speaking it is just an upper limit.
In conclusion, the featureless SMC extinction law provides the best fit to our data, improving the quality of the fit obtained for an unextincted afterglow (see Table 4). It is interesting to note the dramatic dependence of the quality of fit on the existence of the 2175 Å absorption bump. Extinction laws with high moderate dust-to-gas ratios that produce such an absorption feature do not provide good fits to our data points. Therefore, the spectral energy distribution of GRB~000301C supports a scenario where the host is in the early stages of chemical enrichment.
The power-law + extinction fits to the SED in the NIR–optical range allow us to predict the UV flux at the spectral range of MAMA-HST when the UV spectrum was obtained at 2000 March 6.375 UT (Smette et al. SFG2000 (2000)) assuming that the shape of the SED has not changed between the two epochs. This assumption is supported by the imaging data (Sect. 6.3). First, we consider the best fit to the SED at 2000 March 4.39–4.55 UT and calculate the flux at March 4.39-4.55 UT at 3000 Å. Then, making use of the light curve models (presented in Table 6), we estimate the value of the flux at 3000 Å for March 6.375 UT. The predicted flux ranges from 5.9 $`\times 10^{18}`$ erg cm$`^2`$s<sup>-1</sup> Å<sup>-1</sup> to 7.9 $`\times 10^{18}`$ erg cm$`^2`$s<sup>-1</sup> Å<sup>-1</sup>, depending on the light curve model. A final analysis (Smette et al. SFG2001 (2001)) of the MAMA-HST data revealed a flux of $`7.3_{1.8}^{+0.8}10^{18}`$ erg cm<sup>-1</sup> consistent with our extrapolation.
### 6.3 Evolution of the spectral energy distribution
From our UBRI photometric data (presented in Table LABEL:TABLE:photometry) we have multi-band optical coverage from 2 to 10 days after the burst-trigger (on March 1.41 UT). When analysing the colours, we find that the simplest reliable fit is for constant colours. Thus we find no evidence for optical chromatic evolution for the afterglow during the period of observations (see Fig. 8). For these constant fits, we obtain the values presented in Table 5. These values are not corrected for Galactic or intrinsic extinction.
### 6.4 The light curve
According to the simple fireball model the optical afterglow should follow a power-law decay, $`F_\nu \nu ^\beta t^\alpha `$ (Sari et al. SPN1998 (1998)). However, a single power-law is excluded at more than the $`99.9`$% confidence level. The parameters for this power-law are given in Table 6. The photometry suggests that the optical afterglow follows a shallow power-law decay for the first few days and then steepens. This behaviour has been seen previously in GRB~980519 (Jaunsen et al. J2001 (2001)), GRB~990123 (Kulkarni et al. KDO1999 (1999)), GRB~990510 (Harrison et al. HBF1999 (1999)), GRB~991208 (Castro-Tirado et al. CT2001 (2001)) and GRB~000926 (Fynbo et al. 2001b and is predicted by many models for gamma-ray bursts (see below). Sagar et al. (SMP2000 (2000)) report that there are seven components to the R-band light curve. Here we are primarily interested in the overall structure of the light curve, not the structure at small time scales. Therefore, we fit a broken power-law of the form,
$$f_\nu (t)=\{\begin{array}{ccc}f_\nu (t_b)\left(\frac{t}{t_b}\right)^{\alpha _1},\hfill & \mathrm{if}\hfill & tt_b\hfill \\ f_\nu (t_b)\left(\frac{t}{t_b}\right)^{\alpha _2},\hfill & \mathrm{if}\hfill & tt_b,\hfill \end{array}$$
(1)
to the UBRI data presented in Table LABEL:TABLE:photometry. The flux, in $`\mu `$Jy, at time $`t`$ days after the burst is denoted by $`f_\nu (t)`$. The time of the break in the decay is denoted $`t_b`$. The slope before the break is $`\alpha _1`$, and the slope after the break is $`\alpha _2`$. The flux at the time of the break is $`f_\nu (t_b)`$. We used CERN’s Minuit function minimization package, and a chi-square minimization scheme, to simultaneously solve for the four free parameters ($`\alpha _1`$, $`\alpha _2`$, $`t_b`$, and $`f_\nu (t_b)`$) and their formal 1-$`\sigma `$ errors in the fit for each parameter.
The data was corrected for Galactic reddening and extinction before the fits were made. No corrections were made for reddening or extinction in the host galaxy. The photometry was transformed to the R band using the colours given in Table 5 and then converted to units of flux using a photometric zero point of $`f_{\nu ,0}=3.02\times 10^{20}`$ erg cm<sup>-1</sup> s<sup>-1</sup> Hz<sup>-1</sup> (Fukugita et al. FSI1995 (1995)).
The best fit to the combined UBRI photometry is listed in Table 6, and shown in Fig. 9. To test the sensitivity of the results to the fitting function, we also fit our data with the smooth function used by Stanek et al. (SGK1999 (1999), their Eq. 1) on GRB~990510. The results are given in Table 6. The broken power-law gives the smallest chi-square value, and the errors in the individual parameters are smaller for the broken power-law fit than they are for the smooth function. The correlation coefficient between $`t_b`$ and $`\alpha _1`$ is $`0.39`$ and the coefficient between $`t_b`$ and $`\alpha _2`$ is $`0.84`$. The broken power-law fit is consistent with the data at the 43% confidence level. Even though we, from the theory of fireballs, would expect that the light-curve evolution is a smooth function, we find in the case of GRB~000301C, in agreement with Berger et al. (BSF2000 (2000)), that the broken power-law provides the most reliable fit. Additionally, a broken power-law provides the most reliable metod of determining the time when the decay of the light has steepened, and thus is a useful way of parameterising the data.
We combined all our UBRI-data (Table LABEL:TABLE:photometry) with those from the literature. Fig. 10 shows the best-fitting broken power-law for all of the UBRI data in the literature (Sagar et al. SMP2000 (2000) and references therein) and Table LABEL:TABLE:photometry. This data was shifted to the R band in the manner described above. The parameters of the fit are shown in Fig. 10 and are not significantly different from the parameters of the broken power-law that was fit to our data (see Table 6). The conspicuous short-term behaviour of the light-curve has been detailed by Masetti et al. (MBB2000 (2000)), Sagar et al. (SMP2000 (2000)) and Berger et al. (BSF2000 (2000)). Garnavich, Loeb & Stanek (GLS2000 (2000)) find that the variation of the lightcurve can be interpreted as a microlensing event, peaking about 3.5 days after the burst, superposed on a power-law broken at $`t_b=7.6`$ days. This superposed event peaks at a more sparsely sampled period in our data, coinciding partly with where we identify the break. Thus it is not possible, from our data, to further constrain the existence of such an event. We choose here to work only with the data presented in Table LABEL:TABLE:photometry as it represents a consistently derived set.
## 7 Discussion
### 7.1 Interpretation of the light curve
The fit to the multi-colour light curves shows a break at $`t_b=4.39\pm 0.26`$ days, with the light curve steepening from $`\alpha _1=0.72\pm 0.06`$ to $`\alpha _2=2.29\pm 0.17`$, i.e., by $`\mathrm{\Delta }\alpha =\alpha _1\alpha _2=1.57\pm 0.18`$. A broken light curve can arise in a number of circumstances: i) If the frequency separating fast cooling electrons from slow cooling ones moves through the optical at $`t_b`$, the resulting light curve would steepen by $`\mathrm{\Delta }\alpha 0.25`$ (Sari et al. SPN1998 (1998)). ii) The light curve may also steepen if a spherical fireball slows down to a non-relativistic expansion (Dai & Lu DL1999 (1999)), resulting in $`\mathrm{\Delta }\alpha =(\alpha _1+3/5)=0.12`$ for our value of $`\alpha _1`$. iii) If the outflow is collimated with a fixed opening angle, the break in the light curve occurs when the relativistic beaming of the synchrotron radiation becomes wider than the jet opening angle (Mészáros & Rees MR1999 (1999)). In this case the break is a geometrical effect and the steepening is $`\mathrm{\Delta }\alpha =3/4`$. iv) If the afterglow arises in a sideways expanding jet, the steepening will be $`\mathrm{\Delta }\alpha =(1\alpha _1/3)=1.24`$ (Rhoads R1999 (1999)) for our value of $`\alpha _1`$. The above estimates all assume a constant mean density distribution of the ambient medium. We note that collimated outflows in general result in faster decaying light curves than the spherically symmetric ones. If the mean density distribution is not constant, e.g. it has a stellar wind density profile, the light curves also decay faster, but the break will be less pronounced (Panaitescu et al. PMR1998 (1998)).
Based on the light-curve properties alone, the model that best fits the observations is that of a sideways expanding jet in an ambient medium with a constant mean density distribution. In that interpretation, the observed light curve indices imply an electron energy distribution index of $`p=2.13\pm 0.09`$ that results in a theoretical spectral index of $`\beta =(p1)/2=0.56\pm `$0.05. This is in agreement with the spectral index of $`\beta =0.70\pm `$0.09 inferred from our spectroscopic observations when correcting for extinction in the host galaxy (Sect. 6.2), independently strengthening the described model for the afterglow.
With a combined fluence in the 25 – 1000 KeV range of $`4\times 10^6`$ erg cm<sup>-2</sup>, and a redshift of $`z=2.0404`$, the isotropic energy release of GRB 000301C is $`E=4.6\times 10^{52}`$ erg. Following Rhoads (R1999 (1999)), the energy estimate and the light curve break time, $`t_b=4.39\pm 0.26`$ days, implies a jet opening angle, at that time, of $`\theta 15^{}n^{1/8}`$, where $`n`$ is the number density of the ambient medium (in units of cm<sup>-3</sup>), and the break is assumed to occur when the opening angle equals the bulk Lorentz factor, $`\theta =\mathrm{\Gamma }`$.
This interpretation is similar to that of GRB 990510 which was almost 5 times more energetic, but had a jet opening angle of approximately $`5^{}`$, leading to an earlier break in the light curve. The best fit to GRB 990510 was a smooth function (e.g. Stanek et al. SGK1999 (1999); Harrison et al. HBF1999 (1999); Holland et al. HBH2000 (2000)), as compared to a broken power-law in the case of GRB 000301C.
### 7.2 The host galaxy
The host galaxy appears to be very faint. There is no evidence for any extended emission from a host galaxy in any of the data presented in this paper. Deep images obtained with HST+STIS about 1 month after the GRB indicate that any host galaxy must be R$``$27.8$`\pm `$0.25 (Fruchter et al. FSG2000 (2000)). Hence, we can safely conclude that the host galaxy of GRB~000301C is very faint compared to other known populations of galaxies at high redshifts. From fitting extinction curves to our photometry, we have found evidence for some extinction ($`\mathrm{A}_V`$ $``$ 0.1) in the host galaxy. The A<sub>V</sub> derived from the best fit corresponds to an absorption at rest-frame 1500 Å of about 0.4 magnitudes.
For comparison, Lyman Break Galaxies (LBGs) at slightly higher redshifts (z $``$ 2.5–3.5) on average have values of extinction at rest frame wavelength 1500 Å of approximately 1.7 magnitudes, and, in rare cases up to 5 magnitudes (Steidel et al. SADP1999 (1999)). Hence, the faint optical appearance of the host galaxy relative to the star-forming LBGs at high redshift is most likely not imposed by massive extinction, but is rather due to a lower overall star formation rate of the host galaxy.
As the host galaxy furthermore has a very high H i column density, log(NH i)=21.2$`\pm `$0.5 as derived from the Ly$`\alpha `$ absorption feature and supported by the strong Lyman break (Smette et al. SFG2000 (2000)), it is interesting to compare with the population of galaxies identified as Damped Ly$`\alpha `$ Absorbers (Wolfe et al. WTS1986 (1986)) in the spectra of background QSOs. These galaxies have H i column densities higher than log(NH i)=20.3. Based on the luminosity function of LBGs and the typical impact parameters of DLAs, Fynbo et al. (FMW1999 (1999)) show that the majority of DLAs at $`z=3`$ must be fainter than the current flux limit for LBGs of R=25.5 and that there hence is a very abundant population of galaxies fainter than the LBG flux limit. A similar conclusion has been reached by Haehnelt et al. (HSR2000 (2000)). The dust-to-gas ratio towards the line-of-sight of GRB 000301C gives a value of A<sub>V</sub>/N($`H\text{i})0.1/1.6\times 10^{21}`$ cm$`{}_{}{}^{2}=0.6\times 10^{22}`$ cm<sup>-2</sup>. This upper limit is, within errors, consistent with the expected A<sub>V</sub>/N($`H\text{i})`$ for DLAs. The corresponding value for the Milky Way is $`2\times 10^{22}`$ cm<sup>-2</sup> (Allen ALL2000 (2000)).
It is still uncertain what fraction of the integrated star formation rate at high redshift is accounted for by the LBGs and what fraction has to be accounted for by galaxies further down the luminosity function. The relative occurrence of GRBs in a given population of galaxies is expected to be proportional to its relative contribution to the total star formation rate (Totani et al. T1997 (1997); Wijers et al. WBB1998 (1998); Mao & Mo MM1999 (1999); Blain & Natarajan BN2000 (2000)). However, so far only one GRB (GRB~971214 at $`z=3.418`$) is confirmed to have occurred in a galaxy similar to the faint members of the LBGs selected in ground based surveys (Odewahn et al. ODK1998 (1998)). The fact that GRB~000301C occurred in an intrinsically very faint galaxy and that most GRBs with identified OTs have occurred in L or sub-L galaxies, suggest that a large fraction of total star formation at high redshift occurs in a population of galaxies that is further down the luminosity function than the bright LBGs found in ground based surveys and that is likely to have a large overlap with the DLAs.
## 8 Conclusion
GRB~000301C is so far the GRB of shortest duration, for which a counterpart has been detected. The high-energy properties of the burst are consistent with membership of the short-duration class of GRBs, though GRB~000301C could belong to the proposed intermediate class of GRBs or the extreme short end of the distribution of long-duration GRBs. Our VLT-spectra show that GRB~000301C occurred at a redshift of 2.0404$`\pm `$0.0008. The light curve of the optical transient is well-fitted by a broken power-law and it is consistent with being achromatic. From the light-curve properties we find that the best model for GRB~000301C is that of a sideways expanding jet in an ambient medium of constant density. This interpretation is further supported by the achromatic light-curve evolution, and by the agreement between the theoretically predicted and observationally derived spectral indices. The spectral energy distribution at March 4.5 reveals SMC-like extinction in the host galaxy at a level of $`\mathrm{A}_V<0.10`$, which is significantly lower than for the strongly star-forming LBGs. Hence, the extreme faintness of the host galaxy indicates a low overall star-formation rate in the host galaxy, raising the possibility that the host may be a chemically less evolved, relatively low-luminosity galaxy containing SMC-type dust. We argue that there may be a connection between the host galaxy of GRB~000301C and DLAs, suggesting that substantial star-forming activity at high redshift takes place in relatively faint galaxies. Future studies of high redshift GRBs will further help explore this connection.
## Acknowledgments
We wish to thank our anonymous referee for providing many helpful comments and suggestions. Support for this programme by the director of the Nordic Optical Telescope, professor Piirola, is much appreciated. We also acknowledge the assistance given by the ESO service observing team. J. Gorosabel acknowledges the receipt of a Marie Curie Research Grant from the European Commission. J. Hjorth acknowledges support from the Danish Natural Science Research Council (SNF). B. Thomsen acknowledges support from the Danish Natural Science Research Council for funding the Danish Centre for Astrophysics with the HST. G. Björnsson acknowledges support from the Icelandic Research Council and the University of Iceland Research Fund. I. Burud is supported by Pôle d’Attraction Interuniversitaire, P4/05 (SSTC, Belgium). K. Hurley acknowledges support for Ulysses operations under JPL Contract 958056, for IPN operations under NASA LTSA grant NAG5-3500, and for NEAR operations under the NEAR Participating Scientist program. We are grateful to R. Gold and R. McNutt for their assistance with the NEAR spacecraft. R. Starr is supported by NASA grant NCC5-380. We are indebted to T. Sheets for her excellent work on NEAR data reduction. Special thanks also go to the NEAR project office for its support of post-launch changes to XGRS software that made these measurements possible. In particular, we are grateful to John R. Hayes and Susan E. Schneider for writing the GRB software for the XGRS instrument and to Stanley B. Cooper and David S. Tillman for making it possible to get accurate universal time for the NEAR GRB detections. The data presented here have been taken using ALFOSC, which is owned by the Instituto de Astrofisica de Andalucia (IAA) and operated at the Nordic Optical Telescope under agreement between IAA and the NBIfAFG of the Astronomical Observatory of Copenhagen. Additionally, the availability of the GRB Coordinates Network (GCN) and BACODINE services, maintained by Scott Barthelmy, is greatly acknowledged. We acknowledge the availability of POSS-II exposures, used in this work; The Second Palomar Observatory Sky Survey (POSS-II) was made by the California Institute of Technology with funds from the National Science Foundation, the National Aeronautics and Space Administration, the National Geographic Society, the Sloan Foundation, the Samuel Oschin Foundation, and the Eastman Kodak Corporation.
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# Isospin Symmetry Breaking in Nuclei — ONS Anomaly —
## INTRODUCTION
The discrepancy between the calculated binding energy differences of mirror nuclei and those measured is a long-standing problem in nuclear physics. It is known as the Okamoto-Nolen-Schiffer (ONS) anomaly oka ; nol . Although it was first thought that electromagnetic effects could almost account for the observed binding energy differences, it is now believed that the ONS anomaly has its origin in charge symmetry breaking (CSB) in the strong interaction mil . In addition to calculations based on charge symmetry violating meson exchange potentials mil ; blu ; sha , a number of quark-based calculations have been performed hen ; sai in an attempt to resolve this anomaly. In such calculations, CSB enters through the up (u) and down (d) current quark mass difference in QCD. Despite these efforts, the difficulty of producing a realistic description of nuclear structure on the basis of explicit quark degrees of freedom has hindered the direct calculation of the binding energy differences.
In this study we report the results for the binding energy differences of the valence (excess) proton and neutron of the mirror nuclei, <sup>15</sup>O – <sup>15</sup>N, <sup>17</sup>F – <sup>17</sup>O, <sup>39</sup>Ca – <sup>39</sup>K and <sup>41</sup>Sc – <sup>41</sup>Ca, calculated using a quark-based model involving explicit nuclear structure and shell effects, namely the quark-meson coupling (QMC) model gui . This model has been successfully applied not only to traditional nuclear problems gui but also to other new areas as well tsu . Although some exploratory QMC results on the ONS anomaly have already been reported sai , an early version of the model was used there, and it was applied to finite nuclei only through local density approximation, rather than a consistent shell model calculation.
## THE QUARK-MESON COUPLING MODEL
In this section, we introduce the QMC model, and then report the medium modification of the nucleon properties in finite nuclei gui .
### Effect of Nucleon Structure
Let us suppose that a free nucleon (at the origin) consists of three light (u and d) quarks under a (Lorentz scalar) confinement potential, $`V_c`$. Then, the Dirac equation for the quark field $`\psi _q`$ is given by
$$[i\gamma m_qV_c(r)]\psi _q(r)=0,$$
(1)
where $`m_q`$ is the bare quark mass.
Next we consider how Equation (1) is modified when the nucleon is bound in static, uniformly distributed (iso-symmetric) nuclear matter. In the QMC model gui it is assumed that each quark feels scalar, $`V_s^q`$, and vector, $`V_v^q`$, potentials, which are generated by the surrounding nucleons, as well as the confinement potential. This assumption seems appropriate when the nuclear density $`\rho _B`$ is near around normal nuclear matter density ($`\rho _0=0.15`$ fm<sup>-3</sup>). If we use the mean-field approximation (MFA) for the meson fields, Equation (1) may be rewritten as
$$[i\gamma (m_qV_s^q)V_c(r)\gamma _0V_v^q]\psi _q(r)=0.$$
(2)
The potentials generated by the medium are constants because the matter distributes uniformly. As the nucleon is static, the time-derivative operator in the Dirac equation can be replaced by the quark energy, $`iϵ_q`$. By analogy with the procedure applied to the nucleon in QHD ser , if we introduce the effective quark mass by $`m_q^{}=m_qV_s^q`$, the Dirac equation (2) can be rewritten in the same form as that in free space with the mass $`m_q^{}`$ and the energy $`ϵ_qV_v^q`$, instead of $`m_q`$ and $`ϵ_q`$. In other words, the vector interaction has no effect on the nucleon structure except for an overall phase in the quark wave function, which gives a shift in the nucleon energy. This fact does not depend on how to choose the confinement potential, $`V_c`$. Then, the nucleon energy at rest in the medium is given by $`E_N=M_N^{}(V_s^q)+3V_v^q`$, where the effective nucleon mass $`M_N^{}`$ depends on only the scalar potential.
We can extend this idea to finite nuclei gui . Let us suppose that the scalar and vector potentials in Equation (2) are mediated by the $`\sigma `$ and $`\omega `$ mesons, and introduce their mean-field values, which now depend on position $`\stackrel{}{r}`$, by $`V_s^q(\stackrel{}{r})=g_\sigma ^q\sigma (\stackrel{}{r})`$ and $`V_v^q(\stackrel{}{r})=g_\omega ^q\omega (\stackrel{}{r})`$, respectively, where $`g_\sigma ^q`$ ($`g_\omega ^q`$) is the coupling constant of the quark-$`\sigma `$ ($`\omega `$) meson. Furthermore, we shall add the isovector vector meson, $`\rho `$, and the Coulomb field, $`A`$, to describe finite nuclei realistically. Then, the effective Lagrangian density for finite nuclei would be given by gui
$`_{QMC}`$ $`=`$ $`\overline{\psi }[i\gamma M_N^{}g_\omega \omega \gamma _0g_\rho {\displaystyle \frac{\tau _3^N}{2}}b\gamma _0{\displaystyle \frac{e}{2}}(1+\tau _3^N)A\gamma _0]\psi `$
$``$ $`{\displaystyle \frac{1}{2}}[(\sigma )^2+m_\sigma ^2\sigma ^2]+{\displaystyle \frac{1}{2}}[(\omega )^2+m_\omega ^2\omega ^2]+{\displaystyle \frac{1}{2}}[(b)^2+m_\rho ^2b^2]+{\displaystyle \frac{1}{2}}(A)^2,`$
where $`\psi `$ and $`b`$ are respectively the nucleon and the $`\rho `$ fields. $`m_\sigma `$, $`m_\omega `$ and $`m_\rho `$ are respectively the masses of the $`\sigma `$, $`\omega `$ and $`\rho `$ mesons. $`g_\omega `$ and $`g_\rho `$ are respectively the $`\omega `$-N and $`\rho `$-N coupling constants, which are given by $`g_\omega =3g_\omega ^q`$ and $`g_\rho =g_\rho ^q`$ (where $`g_\rho ^q`$ is the quark-$`\rho `$ coupling constant).
If we define the field-dependent $`\sigma `$-N coupling constant, $`g_\sigma (\sigma )`$, by gui
$$M_N^{}(\sigma (\stackrel{}{r}))M_Ng_\sigma (\sigma (\stackrel{}{r}))\sigma (\stackrel{}{r}),$$
(4)
where $`M_N`$ is the free nucleon mass, it is easy to compare with QHD ser . The difference between QMC and QHD lies only in the coupling constant $`g_\sigma `$, which depends on the scalar field in QMC while it is constant in QHD. However, this difference leads to a lot of favorable results gui .
Now we need a model for the structure of the nucleon involved to perform actual calculations. We here use the MIT bag model. In the present model, the bag constant, $`B`$, and the $`z`$ parameter for the nucleon are fixed to reproduce the free nucleon mass ($`M_N`$ = 939 MeV) and the free bag radius $`R_N`$ = 0.8 fm. In the following we choose $`m_q`$ = 5 MeV and set $`m_\sigma `$ = 550 MeV, $`m_\omega `$ = 783 MeV and $`m_\rho `$ = 770 MeV. (Variations of the quark mass and $`R_N`$ only lead to numerically small changes in the calculated results gui .) We then find that $`B^{1/4}`$ = 170.0 MeV and $`z`$ = 3.295.
For infinite nuclear matter, from the Lagrangian density (Effect of Nucleon Structure), we can easily find the total energy per nucleon, $`E_{tot}/A`$, and the mean-field values of $`\omega `$ and $`\rho `$ (which are respectively given by baryon number conservation and the difference in proton and neutron densities). The scalar mean-field is given by a self-consistency condition, $`E_{tot}/\sigma =0`$ gui . The coupling constants, $`g_\sigma ^2`$ and $`g_\omega ^2`$, are fixed to fit the average binding energy ($`15.7`$ MeV) at $`\rho _0`$ for nuclear matter. Furthermore, the $`\rho `$-N coupling constant is used to reproduce the bulk symmetry energy, 35 MeV. We then find gui : $`g_\sigma ^2/4\pi `$ = 5.40, $`g_\omega ^2/4\pi `$ = 5.31, $`g_\rho ^2/4\pi `$ = 6.93, and the nuclear incompressibility, $`K280`$ MeV. Note that the model gives the variation of the nucleon bag radius, $`\delta R_N^{}/R_N=0.02`$, the lowest quark eigenvalue, $`\delta x_q^{}/x_q=0.16`$ and the root-mean-square radius of the quark wave function, $`\delta r_q^{}/r_q=+0.02`$, at saturation density.
Using these parameters, we can solve a finite nuclear system. As an example, we show charge density distribution of <sup>40</sup>Ca in Figure 1. The QMC model can reproduce the properties of not only nuclear matter but also finite nuclei (for more details, see gui ).
### Nucleon Mass in Nuclear Matter
Here we consider the nucleon mass in matter furthermore. The nucleon mass is a function of the scalar field. Because the scalar field is small at low density the mass may be expanded in terms of $`\sigma `$ as
$$M_N^{}=M_N+\left(\frac{M_N^{}}{\sigma }\right)_{\sigma =0}\sigma +\frac{1}{2}\left(\frac{^2M_N^{}}{\sigma ^2}\right)_{\sigma =0}\sigma ^2+𝒪(\sigma ^3).$$
(5)
Since the interaction Hamiltonian between the nucleon and the $`\sigma `$ field at the quark level is given by $`H_{int}=3g_\sigma ^q𝑑\stackrel{}{r}\overline{\psi }_q\sigma \psi _q`$, the derivative of $`M_N^{}`$ with respect to $`\sigma `$ is $`3g_\sigma ^q𝑑\stackrel{}{r}\overline{\psi }_q\psi _q3g_\sigma ^qS_N(\sigma )`$, where we have defined the quark scalar charge in the nucleon, $`S_N(\sigma )`$, which is itself a function of $`\sigma `$. Because of a negative value of the derivative, the nucleon mass decreases in matter at low density.
Furthermore, we define the scalar-charge ratio, $`S_N(\sigma )/S_N(0)`$, to be $`C_N(\sigma )`$ and the $`\sigma `$-N coupling constant in free space to be $`g_\sigma `$ (i.e., $`g_\sigma =g_\sigma (\sigma =0)=3g_\sigma ^qS_N(0)`$). Using these quantities, we find
$$M_N^{}=M_Ng_\sigma \sigma \frac{1}{2}g_\sigma C_N^{}(0)\sigma ^2+𝒪(\sigma ^3).$$
(6)
In general, $`C_N`$ is a decreasing function because the quark in matter becomes more relativistic than in free space. Thus, $`C_N^{}(0)`$ takes a negative value. If the nucleon were structureless $`C_N`$ would not depend on $`\sigma `$. Therefore, only the first two terms in the RHS of Equation (6) remain, which is exactly the same as the effective nucleon mass in QHD ser .
## CHARGE SYMMETRY BREAKING IN QMC
Now we introduce the charge symmetry breaking in the QMC model sai ; tsu2 . The charge symmetry is explicitly broken at the quark level through their masses. We use different values for the u and d current quark masses, and the effective proton, $`M_p^{}`$, and neutron, $`M_n^{}`$, masses. At position $`\stackrel{}{r}`$ in a nucleus (the coordinate origin is taken at the center of the nucleus), the Dirac equations for the quarks in the proton or neutron bag are given by
$$\left[i\gamma _x\left(\left(\begin{array}{c}m_u\\ m_d\end{array}\right)V_\sigma ^q(\stackrel{}{r})\right)\gamma ^0\left(V_\omega ^q(\stackrel{}{r})\pm \frac{1}{2}V_\rho ^q(\stackrel{}{r})\right)\right]\left(\begin{array}{c}\psi _u(x)\\ \psi _d(x)\end{array}\right)=0,$$
(7)
where $`|\stackrel{}{x}\stackrel{}{r}|R_j^{}`$ ($`j`$ specifies proton or neutron). Note that we have assumed that the scalar potential is common to both the u and d quarks. The nucleon and meson fields are calculated self-consistently by solving a set of coupled non-linear differential equations, derived from the effective Lagrangian density (Effect of Nucleon Structure) with the proper modifications caused by the different proton and neutron (or u and d quark) masses in MFA. Thus, the present calculation is free from the sort of double counting questioned by Auerbach aue , and includes the shell effects, which were discussed by Cohen et al. coh .
Before discussing the results obtained, we again need to specify the parameters and inputs used in the calculation tsu2 . They are summarized in Table 1. The bag constant, $`B`$, and the $`z`$ parameter are determined by the bare proton mass, after allowing for the electromagnetic self-energy correction $`+`$0.63 MeV, with the bag radius, $`R_p=0.8`$ fm, in free space. For the neutron, the procedure is the same as that for the proton, allowing for the electromagnetic self-energy correction, $``$0.13 MeV, but using the values of $`B`$ and $`z`$ determined above and calculating the d current quark mass and the bag radius for the neutron. Thus, the u current quark mass ($`m_u=5`$ MeV) is the basic input parameter used to fix the model parameters so as to reproduce the bare proton and neutron masses in free space. We found $`m_d=9.2424`$ MeV in the present calculation.
The coupling constants, $`g_\sigma ^q`$ and $`g_\omega ^q`$, are determined so as to fit the saturation properties of symmetric nuclear matter tsu2 . In Table 1, SU(2) stands for the parameters and inputs obtained and used for the calculation when SU(2) symmetry is assumed, namely $`m_u=m_d=5`$ MeV. We then found: ($`g_\sigma ^q`$, $`g_\omega ^q)=`$ (5.698, 2.744) for CSB, and (5.685, 2.721) for SU(2). For the quark-$`\rho `$ meson coupling constant, to make a realistic estimate, we here use the phenomenological value, $`g_\rho ^q=4.595`$, the value at zero three-momentum transfer corresponding to Hartree approximation, from Table 4.1 of Ref. mac . (Note that because the QMC model does not contain the $`\rho `$-nucleon tensor coupling gui , this gives an unrealistically large value for the coupling constant tsu2 .)
### Proton and Neutron Masses in Nuclear Matter
As in Equation (5), the proton and neutron masses are again given by functions of $`\sigma `$ in matter, and may be expanded in terms of $`\sigma `$ at low $`\rho _B`$
$`M_p^{}`$ $`=`$ $`M_p+(3g_\sigma ^q){\displaystyle \frac{1}{3}}[2S_{u/p}(0)+S_{d/p}(0)]\sigma +𝒪(\sigma ^2),`$ (8)
$`M_n^{}`$ $`=`$ $`M_n+(3g_\sigma ^q){\displaystyle \frac{1}{3}}[S_{u/n}(0)+2S_{d/n}(0)]\sigma +𝒪(\sigma ^2).`$ (9)
Because $`m_um_d`$, the u-quark scalar charge is no longer the same as the d-quark scalar charge. We have therefore introduced four kinds of quark scalar charges in the expansion: $`S_{i/j}(\sigma )=_{V_j}𝑑\stackrel{}{r}\overline{\psi }_{i/j}\psi _{i/j}`$, where $`i`$ denotes u or d quark, $`V_j`$ is the volume of $`j`$ (= p or n) and $`\psi _{i/j}`$ is the $`i`$ quark wave function in $`j`$. Since the proton consists of two u quarks and one d quark, the derivative of $`M_p^{}`$ with respect to $`\sigma `$ is given by $`2S_{u/p}+S_{d/p}`$. Similarly, the derivative for the neutron is given by $`S_{u/n}+2S_{d/n}`$.
Taking the difference between the in-medium neutron and proton masses, we find
$$\mathrm{\Delta }_{np}^{}M_n^{}M_p^{}=\mathrm{\Delta }_{np}(3g_\sigma ^q)[S_n(0)S_p(0)]\sigma +𝒪(\sigma ^2),$$
(10)
where $`\mathrm{\Delta }_{np}=M_nM_p`$, $`S_n(0)=\frac{1}{3}[S_{u/n}(0)+2S_{d/n}(0)]`$ and $`S_p(0)=\frac{1}{3}[2S_{u/p}(0)+S_{d/p}(0)]`$. Here we may expect that $`S_{u/j}<S_{d/j}`$ because the u quark is more relativistic than the d quark in nuclear matter ($`m_u<m_d`$) — note that the quark scalar charge is given in terms of $`\overline{\psi }_q\psi _q`$ in matter. Thus, we find that $`S_n(0)>S_p(0)`$ and then $`\mathrm{\Delta }_{np}^{}<\mathrm{\Delta }_{np}`$ in nuclear medium.
In Figure 2 we show the neutron-proton effective mass difference calculated in symmetric nuclear matter, including the electromagnetic self-energy corrections. One notices that the mass difference becomes smaller as the density increases. This behavior works in the right direction to resolve the ONS anomaly.
### The ONS Anomaly in Mirror Nuclei
Now we are in a position to show our results of the ONS anomaly in mirror nuclei tsu2 . We first present the calculated single-particle energies for <sup>17</sup>F and <sup>17</sup>O in Table 2. These mirror nuclei have a common core nucleus, <sup>16</sup>O, and have an extra valence proton for <sup>17</sup>F and neutron for <sup>17</sup>O. In order to focus on the strong interaction effect for the valence proton and neutron, the Dirac equations are solved without the Coulomb and $`\rho `$-meson potentials, or the electromagnetic self-energy corrections, and keeping only the charge symmetric $`\sigma `$ and $`\omega `$ mean field potentials. Consistently, the valence nucleon contributions are not included in the Coulomb and $`\rho `$-mean field source densities in the core nucleus. However, for the nucleons in the core nucleus, electromagnetic self-energy corrections and the Coulomb potential as well as the $`\rho `$ mean field potential are included in addition to the $`\sigma `$ and $`\omega `$ mean field potentials in solving the Dirac equations. Results are shown for two cases in Table 2: calculation with charge symmetry breaking (denoted by CSB) and calculation performed assuming SU(2) symmetry (denoted by SU(2)).
The SU(2) results for <sup>17</sup>F and <sup>17</sup>O agree perfectly with each other as they should. Single-particle energies in the cores of <sup>17</sup>F and <sup>17</sup>O are slightly different for CSB. This difference is induced by the different (effective) masses for the valence proton and neutron, arising from the charge and density dependence of their coupling to the self-consistent scalar mean field. This also causes a second order effect on the Coulomb and $`\rho `$-meson potentials through the self-consistency procedure.
It is interesting to compare the binding energy differences between the valence proton in <sup>17</sup>F and neutron in <sup>17</sup>O. In CSB, the result gives, $`E(p)(1d_{5/2})E(n)(1d_{5/2})0.18`$ MeV, while the SU(2) case is zero as it should be. This amount already shows a magnitude similar to that of the observed binding energy differences.
In Table 3, we summarize the calculated single-particle energies for the valence proton and neutron of several mirror nuclei (in CSB) tsu2 . Comparing the $`\rho `$-potential contributions for the hole states with core plus valence states, one notices the shell effects due to the $`\rho `$-potentials. These results reflect the difference in the shell structure, namely the hole states tend to have larger $`\rho `$-potential contributions than the core plus valence nucleon states.
The binding energy differences obtained indicate that the prime CSB effects originate in the u-d current quark mass difference. The calculated binding energy differences give of the order of about a few hundred keV. This is precisely the order of magnitude which is observed as the ONS anomaly mil ; sha .
## SUMMARY
Using the QMC model, we have discussed CSB in nuclear medium and calculated the ONS anomaly in mirror nuclei, including the quark degrees of freedom explicitly. We stress that the present contribution to the ONS anomaly is based on a very simple but novel idea, namely the slight difference between the quark scalar densities of the u and d quarks in a bound nucleon, which stems from the u and d quark mass difference sai ; tsu2 . This implies that the in-medium proton-$`\sigma `$ and neutron-$`\sigma `$ coupling constants differ from their values in free space and that the neutron-proton effective mass difference is reduced in matter.
Our results were obtained within an explicit shell model calculation, based on quark degrees of freedom. They show that once CSB is set through the u and d current quark mass difference so as to reproduce the proton and neutron masses in free space, it can produces binding energy differences for the valence (excess) proton and neutron of mirror nuclei of a few hundred keV. The origin of this effect is so simple that it is natural to conclude that a sizable fraction of CSB in mirror nuclei arises from the density dependence of the u and d quark scalar densities in a bound nucleon.
It is a fascinating challenge for the future to compare this result with the traditional mechanism involving $`\rho \omega `$ mixing blu . This will involve the issue of the possible momentum dependence of the $`\rho \omega `$ mixing amplitude mil ; gol . In addition, one would need to examine whether there is any deeper connection between these apparently quite different sources of charge symmetry violation.
## ACKNOWLEDGMENTS
The author would like to thank K. Tsushima, A.W. Thomas and A.G. Williams for valuable discussions. This work was supported by the Australian Research Council and the Japan Society for the Promotion of Science. In the present paper, Fugure 1 was reprinted from Nucl. Phys. A 609 (1996), Saito et al., “Self-consistent description of finite nuclei based on a relativistic quark model”, p.352 (Figure 4), and Figure 2 and Table 1-3 were reprinted from Phys. Lett. B 465 (1999), Tsushima et al., “Charge symmetry breaking in mirror nuclei from quarks”, p.38 (Figure 1 and Table 1), p.39 (Table 2) and p.41 (Table3), with permission from Elsevier Science.
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# 1 The colored Hubbard model. The lattice sites of the coarse lattice are symbolized by a ◆. The numbering of the sites of our original lattice is also shown.
The Hubbard model is one of the most studied models for electron systems. In particular, the two dimensional model appears to be a good candidate for an explanation of high $`T_c`$ superconductivity. Despite its simplicity, several obstacles render even its approximate solution a difficult theoretical task. As a fermionic system it is not easily accessible to numerical simulations. Furthermore, there seems to be a competition between antiferromagnetic order and d–wave superconductivity. Both operators do not correspond to fermion bilinears on the same lattice site.
Recently, promising investigations using approximations to exact renormalization group equations have been started , . The difficulty of these approaches, however, consists in the high complexity of the equations if the full momentum dependence of correlation functions for several fermions is kept. In particular, the low temperature phases can only be realistically described if effective interactions involving more than four fermions are included. In our opinion a prerequisite for a successful use of these methods in the ordered phases is a simplification of the momentum dependence of the interactions. This can be done if the most prominent physical degrees of freedom are identified. We propose here a version of the Hubbard model where the relevant order parameters correspond to translationally invariant vacuum expectation values of scalar fields. We will see that this formulation can describe the low temperature phases in a very simple way. We therefore hope that it constitutes a good starting point for a detailed renormalization group analysis.
We consider the partition function of the Hubbard model
$$Z=D\widehat{\psi }D\widehat{\psi }^{}\mathrm{exp}\left\{S[\widehat{\psi },\widehat{\psi }^{}]+_0^\beta 𝑑\tau \underset{i}{}(\eta _i^{}\widehat{\psi }_i+\eta _i\widehat{\psi }_i^{})\stackrel{~}{S}_j\right\}$$
(1)
$`S[\widehat{\psi },\widehat{\psi }^{}]`$ $`=`$ $`{\displaystyle _0^\beta }d\tau \{{\displaystyle \underset{i}{}}(\widehat{\psi }_i^{}_\tau \widehat{\psi }_i\mu \widehat{\psi }_i^{}\widehat{\psi }_i`$ (2)
$`{\displaystyle \frac{1}{6}}U(\widehat{\psi }_i^{}\stackrel{}{\tau }\widehat{\psi }_i)(\widehat{\psi }_i^{}\stackrel{}{\tau }\widehat{\psi }_i))+{\displaystyle \underset{ij}{}}\widehat{\psi }_i^{}𝒯_{ij}\widehat{\psi }_j\}`$
as a functional of sources $`\eta ,\eta ^{}`$ of the fermions as well as sources of fermion bilinears ($`\stackrel{~}{S}_j`$) that will be specified below (cf. eq. (S0.Ex27)). The spinors $`\widehat{\psi }_i=(\widehat{\psi }_i,\widehat{\psi }_i)^T`$, $`\widehat{\psi }_i^{}=(\widehat{\psi }_i^{},\widehat{\psi }_i^{})`$ (as well as the fermionic sources $`\eta _i`$, $`\eta _i^{}`$) are two-component Grassmann variables depending on the Euclidean “time” $`\tau `$ with antiperiodicity $`\widehat{\psi }_i(\beta )=\widehat{\psi }_i(0),\widehat{\psi }_i^{}(\beta )=\widehat{\psi }_i^{}(0)`$. Here $`\beta =1/T`$ is the inverse temperature. We treat $`\psi `$ and $`\psi ^{}`$ as independent Grassmann variables, even though the notation is reminiscent of a type of complex conjugation which also inverts $`\tau `$ and reorders all Grassmann variables. (In quantum field theory the invariance of the action under this discrete transformation is related to Osterwalder-Schrader positivity.)
The index $`i`$ labels the sites of the lattice. We concentrate on a quadratic lattice in two dimensions, $`i=(m,n)`$, $`m,n`$, with next neighbor interactions where $`𝒯_{ij}=t`$ for $`i,j`$ next neighbors and $`𝒯_{ij}=0`$ otherwise. They describe the probability of fermion tunneling between different lattice sites. After a Fourier transform the Fermi surface for $`U0`$ is given by
$`2t\left(\mathrm{cos}(aq_1)+\mathrm{cos}(aq_2)\right)=\mu `$
$`|q_\mu |2\mathrm{\Lambda },\mathrm{\Lambda }=\pi /2a`$ (3)
with the lattice spacing $`a`$. The local<sup>4</sup><sup>4</sup>4Here the term local refers to our neglect of the interaction of fermions located at different lattice sites. Coulomb interaction of the fermions involves the Pauli matrices $`\stackrel{}{\tau }`$ and can be rearranged, e.g. $`(\widehat{\psi }_i^{}\stackrel{}{\tau }\widehat{\psi }_i)(\widehat{\psi }_i^{}\stackrel{}{\tau }\widehat{\psi }_i)=3(\widehat{\psi }_i^{}\widehat{\psi }_i)(\widehat{\psi }_i^{}\widehat{\psi }_i)=6n_in_i`$ with $`n_{i()}=\widehat{\psi }_{i()}^{}\widehat{\psi }_{i()}`$. As usual, the expectation values of operators are related to appropriate derivatives of $`Z`$ with respect to the sources. In addition to the discrete lattice symmetries, the model has two obvious continuous symmetries: the $`SU(2)`$-spin rotations and the $`U(1)`$-phase rotations corresponding to charge conservation.
In units where $`\mathrm{}=c=k_B=1`$, the parameters $`T`$, $`\mu `$, $`t`$ and $`U`$ have dimension of mass whereas $`\widehat{\psi },\widehat{\psi }^{}`$ are dimensionless. The partition function is invariant under the rescaling $`(\alpha _+)`$ $`\tau \tau /\alpha ,\mu \alpha \mu ,t\alpha t,U\alpha U,T\alpha T`$. It can therefore only depend<sup>5</sup><sup>5</sup>5This holds up to a possible temperature-dependent factor from the functional measure. on the dimensionless parameter ratios $`\mu /U,T/U,t/U`$ and is independent of the dimensionless combination $`aU`$. Furthermore, the invariance of $`Z`$ under the discrete transformation $`\widehat{\psi }(\tau )\widehat{\psi }(\beta \tau ),\widehat{\psi }^{}(\tau )\widehat{\psi }^{}(\beta \tau ),\mu \mu ,tt`$ (with appropriate transformations of the sources) permits a restriction to $`\mu 0`$. Finally, we may divide the lattice sites $`i`$ into two classes, $`iI_1`$, if $`m`$ and $`n`$ are both even or both odd, $`iI_2`$ otherwise. The transformation $`\widehat{\psi }_{iI_2}\widehat{\psi }_{iI_2}`$, $`\widehat{\psi }_{iI_2}^{}\widehat{\psi }_{iI_2}^{}`$ while leaving $`\widehat{\psi }_{iI_1},\widehat{\psi }_{iI_1}^{}`$ invariant maps $`Z(t)Z(t)`$ (again with an appropriate mapping for the sources). We restrict our discussion to positive $`t`$ and $`\mu `$. Predictions for models with negative $`t`$ or negative $`\mu `$ can easily be obtained from our results by an appropriate mapping.
In order to represent the fermion bilinears of interest in a simple form, we choose to label the variables at four neighboring lattice sites by different colors red, yellow, green and blue, $`\widehat{\psi }_a,a=1,\mathrm{},4`$. (This also allows us to easily extend the formalism to lattices with different atoms.) We take $`\stackrel{}{x}=(x,y)=(ma,na)`$ with $`m,n`$ even as the sites of a coarse lattice with lattice distance $`2a`$ and define
$`\widehat{\psi }_1(\stackrel{}{x})=\widehat{\psi }_{m,n}`$ , $`\widehat{\psi }_2(\stackrel{}{x})=\widehat{\psi }_{m+1,n}`$
$`\widehat{\psi }_4(\stackrel{}{x})=\widehat{\psi }_{m,n1}`$ , $`\widehat{\psi }_3(\stackrel{}{x})=\widehat{\psi }_{m+1,n1}`$ (4)
and similar for $`\widehat{\psi }_a^{}`$ (see fig.1).
The lattice symmetries $`T_x`$ (translation in $`x`$ by $`a`$), $`T_y`$ (translation in $`y`$ by $`a`$), $`R`$ (clockwise rotation by $`90^{}`$ around $`\stackrel{}{x}=0`$), $`I_x`$ (reflection at $`x`$-axis) and $`\stackrel{~}{I}_x`$ (reflection at the axis $`y=\frac{1}{2}a`$) act as
$`T_x`$ $`:`$ $`\widehat{\psi }_1(x,y)\widehat{\psi }_2(x,y),\widehat{\psi }_2(x,y)\widehat{\psi }_1(x+2a,y)`$
$`\widehat{\psi }_4(x,y)\widehat{\psi }_3(x,y),\widehat{\psi }_3(x,y)\widehat{\psi }_4(x+2a,y)`$
$`T_y`$ $`:`$ $`\widehat{\psi }_1(x,y)\widehat{\psi }_4(x,y+2a),\widehat{\psi }_2(x,y)\widehat{\psi }_3(x,y+2a)`$
$`\widehat{\psi }_3(x,y)\widehat{\psi }_2(x,y),\widehat{\psi }_4(x,y)\widehat{\psi }_1(x,y)`$
$`R`$ $`:`$ $`\widehat{\psi }_1(x,y)\widehat{\psi }_2(y,x),\widehat{\psi }_2(x,y)\widehat{\psi }_3(y,x)`$
$`\widehat{\psi }_3(x,y)\widehat{\psi }_4(y,x),\widehat{\psi }_4(x,y)\widehat{\psi }_1(y,x)`$
$`I_x`$ $`:`$ $`\widehat{\psi }_1(x,y)\widehat{\psi }_4(x,y),\widehat{\psi }_2(x,y)\widehat{\psi }_3(x,y)`$
$`\stackrel{~}{I}_x`$ $`:`$ $`\widehat{\psi }_{1,2}(x,y)\widehat{\psi }_{1,2}(x,y),\widehat{\psi }_{3,4}(x,y)\widehat{\psi }_{3,4}(x,y+2a).`$ (5)
The lattice symmetries of the coarse lattice can be composed from the generators $`T_x,R,I_x`$. We note that they also act in color space<sup>6</sup><sup>6</sup>6The local interaction is also invariant under relabelings at fixed $`\stackrel{}{x}`$ $`r`$ $`:`$ $`\widehat{\psi }_1\widehat{\psi }_2\widehat{\psi }_3\widehat{\psi }_4`$ $`i`$ $`:`$ $`\widehat{\psi }_1\widehat{\psi }_4,\widehat{\psi }_2\widehat{\psi }_3`$ $`d`$ $`:`$ $`\widehat{\psi }_1\widehat{\psi }_3,\widehat{\psi }_2\widehat{\psi }_4.`$ This, however, is not a symmetry of the next neighbor interaction..
Useful fermion bilinear operators are $`(c=i\tau _2)`$
$`\stackrel{~}{\sigma }_{ab}(\stackrel{}{x})`$ $`=`$ $`\widehat{\psi }_b^{}(\stackrel{}{x})\widehat{\psi }_a(\stackrel{}{x})`$
$`\stackrel{}{\stackrel{~}{\phi }}_{ab}(\stackrel{}{x})`$ $`=`$ $`\widehat{\psi }_b^{}(\stackrel{}{x})\stackrel{}{\tau }\widehat{\psi }_a(\stackrel{}{x})`$
$`\stackrel{~}{\chi }_{ab}^{(1)}(\stackrel{}{x})`$ $`=`$ $`\widehat{\psi }_b^T(\stackrel{}{x})c\widehat{\psi }_a(\stackrel{}{x})=\stackrel{~}{\chi }_{ba}^{(1)}(\stackrel{}{x})`$
$`\stackrel{~}{\chi }_{ab}^{(2)}(\stackrel{}{x})`$ $`=`$ $`\widehat{\psi }_b^{}(\stackrel{}{x})c\widehat{\psi }_a^T(\stackrel{}{x})=\stackrel{~}{\chi }_{ba}^{(2)}(\stackrel{}{x}).`$ (6)
(We have omitted for simplicity fermion-fermion pairs in the triplet of the spin group, similar to the fermion-antifermion pairs $`\stackrel{}{\stackrel{~}{\phi }}`$. They are antisymmetric in the color indices.) Among the electron-electron (or hole-hole) pairs we concentrate on
$`\stackrel{~}{s}^{(\alpha )}`$ $`=`$ $`\stackrel{~}{\chi }_{11}^{(\alpha )}+\stackrel{~}{\chi }_{22}^{(\alpha )}+\stackrel{~}{\chi }_{33}^{(\alpha )}+\stackrel{~}{\chi }_{44}^{(\alpha )}`$
$`\stackrel{~}{c}^{(\alpha )}`$ $`=`$ $`\stackrel{~}{\chi }_{11}^{(\alpha )}\stackrel{~}{\chi }_{22}^{(\alpha )}+\stackrel{~}{\chi }_{33}^{(\alpha )}\stackrel{~}{\chi }_{44}^{(\alpha )}`$
$`\stackrel{~}{d}^{(\alpha )}`$ $`=`$ $`\stackrel{~}{\chi }_{12}^{(\alpha )}\stackrel{~}{\chi }_{23}^{(\alpha )}+\stackrel{~}{\chi }_{34}^{(\alpha )}\stackrel{~}{\chi }_{41}^{(\alpha )}`$
$`\stackrel{~}{e}^{(\alpha )}`$ $`=`$ $`\stackrel{~}{\chi }_{12}^{(\alpha )}+\stackrel{~}{\chi }_{23}^{(\alpha )}+\stackrel{~}{\chi }_{34}^{(\alpha )}+\stackrel{~}{\chi }_{41}^{(\alpha )}`$
$`\stackrel{~}{v}_x^{(\alpha )}`$ $`=`$ $`\stackrel{~}{\chi }_{23}^{(\alpha )}\stackrel{~}{\chi }_{41}^{(\alpha )},\stackrel{~}{v}_y^{(\alpha )}=\stackrel{~}{\chi }_{12}^{(\alpha )}\stackrel{~}{\chi }_{34}^{(\alpha )}`$
$`\stackrel{~}{t}_1^{(\alpha )}`$ $`=`$ $`\stackrel{~}{\chi }_{11}^{(\alpha )}\stackrel{~}{\chi }_{33}^{(\alpha )},\stackrel{~}{t}_2^{(\alpha )}=\stackrel{~}{\chi }_{22}^{(\alpha )}\stackrel{~}{\chi }_{44}^{(\alpha )}.`$ (7)
which transform as $`s`$-wave $`(\stackrel{~}{s})`$, $`d_{xy}`$-wave $`(\stackrel{~}{c}),d_{x^2y^2}`$-wave $`(\stackrel{~}{d})`$, extended $`s`$-wave $`(\stackrel{~}{e})`$, $`p`$-wave $`(\stackrel{~}{v}_x,\stackrel{~}{v}_y)`$ in the spin singlet state. Similarly, for the electron-hole pairs we select
$`\stackrel{~}{\rho }`$ $`=`$ $`\stackrel{~}{\sigma }_{11}+\stackrel{~}{\sigma }_{22}+\stackrel{~}{\sigma }_{33}+\stackrel{~}{\sigma }_{44}`$
$`\stackrel{~}{p}`$ $`=`$ $`\stackrel{~}{\sigma }_{11}\stackrel{~}{\sigma }_{22}+\stackrel{~}{\sigma }_{33}\stackrel{~}{\sigma }_{44}`$
$`\stackrel{~}{q}_1`$ $`=`$ $`\stackrel{~}{\sigma }_{11}\stackrel{~}{\sigma }_{33},\stackrel{~}{q}_2=\stackrel{~}{\sigma }_{22}\stackrel{~}{\sigma }_{44}`$
$`\stackrel{}{\stackrel{~}{m}}`$ $`=`$ $`\stackrel{}{\stackrel{~}{\phi }}_{11}+\stackrel{}{\stackrel{~}{\phi }}_{22}+\stackrel{}{\stackrel{~}{\phi }}_{33}+\stackrel{}{\stackrel{~}{\phi }}_{44}`$
$`\stackrel{}{\stackrel{~}{a}}`$ $`=`$ $`\stackrel{}{\stackrel{~}{\phi }}_{11}\stackrel{}{\stackrel{~}{\phi }}_{22}+\stackrel{}{\stackrel{~}{\phi }}_{33}\stackrel{}{\stackrel{~}{\phi }}_{44}`$
$`\stackrel{}{\stackrel{~}{g}}_+`$ $`=`$ $`\stackrel{}{\stackrel{~}{\phi }}_{11}\stackrel{}{\stackrel{~}{\phi }}_{33},\stackrel{}{\stackrel{~}{g}}_{}=\stackrel{}{\stackrel{~}{\phi }}_{22}\stackrel{}{\stackrel{~}{\phi }}_{44}`$ (8)
They correspond to the charge density $`\stackrel{~}{\rho }`$, the charge modulation or charge density wave $`\stackrel{~}{p}`$, the ferromagnetic and antiferromagnetic spin densities $`\stackrel{}{\stackrel{~}{m}},\stackrel{}{\stackrel{~}{a}}`$ and the diagonal spin density $`\stackrel{}{\stackrel{~}{g}}_\pm `$. Correspondingly, we specify the source term for the bilinears as
$`\stackrel{~}{S}_j`$ $`=`$ $`{\displaystyle }d\tau {\displaystyle \underset{\stackrel{}{x}}{}}({\displaystyle \underset{\beta }{}}\{j_\beta ^{}(\stackrel{}{x})\stackrel{~}{u}_\beta ^{(1)}(\stackrel{}{x})+j_\beta (\stackrel{}{x})\stackrel{~}{u}_\beta ^{(2)}(\stackrel{}{x})+\stackrel{~}{r}_\beta j_\beta ^{}(\stackrel{}{x})j_\beta (\stackrel{}{x})\}`$
$`+{\displaystyle \underset{\gamma }{}}\{l_\gamma (\stackrel{}{x})\stackrel{~}{w}_\gamma (\stackrel{}{x})+{\displaystyle \frac{1}{2}}r_\gamma l_\gamma (\stackrel{}{x})l_\gamma (\stackrel{}{x})\}\mu \stackrel{~}{w}_\rho (\stackrel{}{x}))+{\displaystyle \frac{1}{8}}r_\rho \beta V_2\mu ^2/a^2`$
with
$`\stackrel{~}{u}^{(\alpha )}`$ $`=`$ $`(\stackrel{~}{s},\stackrel{~}{c},\stackrel{~}{d},\stackrel{~}{e},\stackrel{~}{v}_x,\stackrel{~}{v}_y,\stackrel{~}{t}_1,\stackrel{~}{t}_2)^{(\alpha )},`$
$`\stackrel{~}{w}`$ $`=`$ $`(\stackrel{~}{\rho },\stackrel{~}{p},\stackrel{~}{q}_1,\stackrel{~}{q}_2,\stackrel{}{\stackrel{~}{m}},\stackrel{}{\stackrel{~}{a}},\stackrel{}{\stackrel{~}{g}}_+,\stackrel{}{\stackrel{~}{g}}_{}).`$
The complex sources $`j=(j_s,j_c,j_d,j_e,j_{v_x},j_{v_y},j_{t_1},j_{t_2})`$ and the real sources $`l=(\mu +l_\rho ^{},l_p,l_{q_1},l_{q_2},\stackrel{}{l}_m,\stackrel{}{l}_a,\stackrel{}{l}_{g_+},\stackrel{}{l}_g_{})`$ also depend on $`\tau `$ and from now on we include the chemical potential $`\mu `$ in the source for $`\stackrel{~}{\rho }`$. Here $`V_2`$ is the two-dimensional volume and we will specify the constants $`\stackrel{~}{r}_\beta `$ and $`r_\gamma `$ below.
Using the identities
$`\stackrel{~}{\chi }_{ab}^{(2)}\stackrel{~}{\chi }_{cd}^{(1)}`$ $`=`$ $`\stackrel{~}{\sigma }_{ca}\stackrel{~}{\sigma }_{db}+\stackrel{~}{\sigma }_{cb}\stackrel{~}{\sigma }_{da}`$
$`\stackrel{}{\stackrel{~}{\phi }}_{ab}\stackrel{}{\stackrel{~}{\phi }}_{cd}`$ $`=`$ $`\stackrel{~}{\sigma }_{ab}\stackrel{~}{\sigma }_{cd}2\stackrel{~}{\sigma }_{ad}\stackrel{~}{\sigma }_{cb}`$ (10)
we establish the relations
$`{\displaystyle \frac{1}{2}}\left(\stackrel{~}{s}^{(2)}\stackrel{~}{s}^{(1)}+\stackrel{~}{c}^{(2)}\stackrel{~}{c}^{(1)}\right)+\stackrel{~}{t}_1^{(2)}\stackrel{~}{t}_1^{(1)}+\stackrel{~}{t}_2^{(2)}\stackrel{~}{t}_2^{(1)}`$ $`=`$ $`4{\displaystyle \underset{a}{}}\stackrel{~}{\sigma }_{aa}^2`$
$`{\displaystyle \frac{1}{2}}\left(\stackrel{~}{\rho }^2+\stackrel{~}{p}^2\right)+\stackrel{~}{q}_1^2+\stackrel{~}{q}_2^2`$ $`=`$ $`2{\displaystyle \underset{a}{}}\stackrel{~}{\sigma }_{aa}^2`$
$`{\displaystyle \frac{1}{2}}\left(\stackrel{}{\stackrel{~}{m}}^2\stackrel{}{\stackrel{~}{a}}^2\stackrel{~}{\rho }^2+\stackrel{~}{p}^2\right)+\stackrel{~}{d}^{(2)}\stackrel{~}{d}^{(1)}+\stackrel{~}{e}^{(2)}\stackrel{~}{e}^{(1)}`$
$`+2\left(\stackrel{~}{v}_x^{(2)}\stackrel{~}{v}_x^{(1)}+\stackrel{~}{v}_y^{(2)}\stackrel{~}{v}_y^{(1)}\right)`$ $`=`$ $`0`$
$`{\displaystyle \frac{1}{2}}\left(\stackrel{}{\stackrel{~}{a}}^2+\stackrel{}{\stackrel{~}{m}}^2\right)+\stackrel{}{\stackrel{~}{g}}_+^2+\stackrel{}{\stackrel{~}{g}}_{}^2`$ $`=`$ $`6{\displaystyle \underset{a}{}}\stackrel{~}{\sigma }_{aa}^2.`$ (11)
The interaction term in (2) reads $`\frac{U}{2}_a\stackrel{~}{\sigma }_{aa}^2`$ and it is obvious that the decomposition into products of the above fermion bilinears is not unique.
It is now straightforward to derive a partially bosonized version of the Hubbard model by introducing a suitable identity in the functional integral (1). Inversely, the bosonized partition function
$`Z`$ $`=`$ $`\mathrm{exp}\left\{2{\displaystyle \frac{\mathrm{\Lambda }^2}{h_\rho ^2}}\beta V_2\mu ^2\right\}{\displaystyle D\widehat{\psi }D\widehat{\psi }^{}D\widehat{u}D\widehat{u}^{}D\widehat{w}}`$
$`\mathrm{exp}\left\{{\displaystyle 𝑑\tau \underset{\stackrel{}{x}}{}(_{kin}+_Y+_B+_j)}\right\}`$
$`_{kin}`$ $`=`$ $`{\displaystyle \underset{a}{}}\widehat{\psi }_a^{}(\stackrel{}{x})_\tau \widehat{\psi }_a(\stackrel{}{x})`$
$`t\{\widehat{\psi }_1^{}(x,y)[\widehat{\psi }_2(x,y)+\widehat{\psi }_2(x2a,y)+\widehat{\psi }_4(x,y+2a)+\widehat{\psi }_4(x,y)]`$
$`+\widehat{\psi }_2^{}(x,y)[\widehat{\psi }_1(x+2a,y)+\widehat{\psi }_1(x,y)+\widehat{\psi }_3(x,y+2a)+\widehat{\psi }_3(x,y)]`$
$`+\widehat{\psi }_3^{}(x,y)[\widehat{\psi }_4(x+2a,y)+\widehat{\psi }_4(x,y)+\widehat{\psi }_2(x,y)+\widehat{\psi }_2(x,y2a)]`$
$`+\widehat{\psi }_4^{}(x,y)[\widehat{\psi }_3(x,y)+\widehat{\psi }_3(x2a,y)+\widehat{\psi }_1(x,y)+\widehat{\psi }_1(x,y2a)]\}`$
$`_B`$ $`=`$ $`4\pi ^2{\displaystyle \underset{\beta }{}}\widehat{u}_\beta ^{}(\stackrel{}{x})\widehat{u}_\beta (\stackrel{}{x})+2\pi ^2{\displaystyle \underset{\gamma }{}}\widehat{w}_\gamma (\stackrel{}{x})\widehat{w}_\gamma (\stackrel{}{x})`$
$`_Y`$ $`=`$ $`{\displaystyle \underset{\beta }{}}\stackrel{~}{h}_\beta (\widehat{u}_\beta ^{}(\stackrel{}{x})\stackrel{~}{u}_\beta ^{(1)}(\stackrel{}{x})+\widehat{u}_\beta (\stackrel{}{x})\stackrel{~}{u}_\beta ^{(2)}(\stackrel{}{x})){\displaystyle \underset{\gamma }{}}h_\gamma \widehat{w}_\gamma (\stackrel{}{x})\stackrel{~}{w}_\gamma (\stackrel{}{x})`$
$`_j`$ $`=`$ $`{\displaystyle \underset{\beta }{}}{\displaystyle \frac{4\pi ^2}{\stackrel{~}{h}_\beta }}(j_\beta ^{}(\stackrel{}{x})\widehat{u}_\beta (\stackrel{}{x})+j_\beta (\stackrel{}{x})\widehat{u}_\beta ^{}(\stackrel{}{x})){\displaystyle \underset{\gamma }{}}{\displaystyle \frac{4\pi ^2}{h_\gamma }}l_\gamma (\stackrel{}{x})\widehat{w}_\gamma (\stackrel{}{x})`$ (12)
$`{\displaystyle \underset{a}{}}\left(\eta _a^{}(\stackrel{}{x})\widehat{\psi }_a(\stackrel{}{x})+\eta _a(\stackrel{}{x})\widehat{\psi }_a^{}(\stackrel{}{x})\right)`$
can easily be transformed into a purely fermionic functional integral by performing the Gaussian integration over the complex scalar fields $`\widehat{u}`$ and real scalar fields $`\widehat{w}`$. We choose $`\stackrel{~}{r}_\beta =4\pi ^2/\stackrel{~}{h}_\beta ^2,r_\gamma =4\pi ^2/h_\gamma ^2`$ such that the partition functions (1) and (S0.Ex36) coincide except for the quartic interactions. Indeed, the four fermion interaction resulting from the bosonic functional integration can be more complex than in the original Hubbard model, i.e.
$$_{int}=\underset{\beta }{}\frac{\stackrel{~}{h}_\beta ^2}{4\pi ^2}\stackrel{~}{u}_\beta ^{(2)}\stackrel{~}{u}_\beta ^{(1)}\underset{\gamma }{}\frac{h_\gamma ^2}{8\pi ^2}\stackrel{~}{w}_\gamma \stackrel{~}{w}_\gamma .$$
(13)
Only for particular values of the real positive Yukawa couplings $`\stackrel{~}{h}_\beta ,h_\gamma `$ the partition function (S0.Ex36) is equal to the partition function (1) of the Hubbard model<sup>7</sup><sup>7</sup>7This holds up to an irrelevant source-independent overall normalization factor., namely for (cf. eq.(S0.Ex31))
$`\stackrel{~}{h}_\beta ^2={\displaystyle \frac{\pi ^2}{3}}\stackrel{~}{H}_\beta U,h_\gamma ^2={\displaystyle \frac{\pi ^2}{3}}H_\gamma U`$
$`2\stackrel{~}{H}_s=2\stackrel{~}{H}_c=\stackrel{~}{H}_{t_1}=\stackrel{~}{H}_{t_2}=3\lambda _1`$
$`2\stackrel{~}{H}_d=2\stackrel{~}{H}_e=\stackrel{~}{H}_{v_x}=\stackrel{~}{H}_{v_y}=6\lambda _3`$
$`H_{q_1}=H_{q_2}=6\lambda _2`$
$`H_\rho =3\left(\lambda _2\lambda _3\right),H_p=3\left(\lambda _2+\lambda _3\right)`$
$`H_\stackrel{}{a}=2\lambda _1+\lambda _23\lambda _3+1`$
$`H_\stackrel{}{m}=2\lambda _1+\lambda _2+3\lambda _3+1`$
$`H_{\stackrel{}{g}_+}=H_\stackrel{}{g}_{}=4\lambda _1+2\lambda _2+2,`$ (14)
where the parameters $`\lambda _i`$ obey
$`\lambda _i`$ $`>`$ $`0i=1\mathrm{}3,`$
$`\lambda _2`$ $`>`$ $`\lambda _3,`$
$`2\lambda _1+\lambda _2+1`$ $`>`$ $`3\lambda _3.`$
We emphasize that the choice (S0.Ex45) of the Yukawa couplings is not unique since it depends on the three parameters $`\lambda _i`$. Arbitrary values of $`\lambda _i`$ (within the allowed range) all describe the same Hubbard model. The independence of physical results on the values of $`\lambda _i`$ can be used as a check for the validity of approximations. Furthermore, a large variety of different four-fermion interactions can be described by varying the Yukawa couplings away from the “Hubbard values” (S0.Ex45).
The symmetries $`R`$ and $`I_x`$ as well as the translations by $`2a`$ are easily implemented on the space of bilinears $`\stackrel{~}{u}_\beta ,\stackrel{~}{w}_\gamma `$ and correspondingly for the scalar fields $`\widehat{u}_\beta ,\widehat{w}_\gamma `$. This is not the case for the translations by $`a`$. The above formulation of the bosonization is therefore not optimal yet if – beyond the symmetries of the coarse-grained lattice – the symmetries like $`T_x`$ play an important role (as for the Hubbard model). It is easy to remedy this shortcoming by an extension of the space of bilinears and the corresponding scalar fields. We introduce an additional color index for the fermion bilinears and the scalars by
$`\stackrel{~}{w}_{1\gamma }(\stackrel{}{x})=T_yT_x^1\stackrel{~}{w}_\gamma (\stackrel{}{x}),\stackrel{~}{w}_{2\gamma }(\stackrel{}{x})=T_y\stackrel{~}{w}_\gamma (\stackrel{}{x}),`$
$`\stackrel{~}{w}_{3\gamma }(\stackrel{}{x})=\stackrel{~}{w}_\gamma (\stackrel{}{x}),\stackrel{~}{w}_{4\gamma }(\stackrel{}{x})=T_x^1\stackrel{~}{w}_\gamma (\stackrel{}{x})`$ (15)
and similar for $`\stackrel{~}{u}^{(\alpha )},\widehat{w},\widehat{u}^{(\alpha )},l,j`$. Products like $`w_\gamma w_\gamma `$ are now understood as scalar products
$$w_\gamma w_\gamma \frac{1}{4}\underset{a}{}w_{a\gamma }w_{a\gamma }.$$
(16)
With these replacements<sup>8</sup><sup>8</sup>8The term $`\mu \stackrel{~}{w}_\rho (\stackrel{}{x})`$ in eq. (S0.Ex27) becomes $`(\mu /4)_a\stackrel{~}{w}_{a\rho }(\stackrel{}{x})`$. it is straightforward to check that the partition function (S0.Ex36) with the choice of Yukawa couplings (S0.Ex45) is again exactly equal to the one of the Hubbard model if all sources except $`\mu `$ are zero. The translations $`T_x,T_y`$ are now directly implemented on the scalar fields, e.g. $`T_x(\widehat{w}_{1\gamma }(\stackrel{}{x}))=\widehat{w}_{2\gamma }(\stackrel{}{x})`$. The same holds true for the rotation $`\stackrel{~}{R}`$ or the reflection $`\stackrel{~}{I}_x`$.
In conclusion, we have developed an extended version of the Hubbard model – the colored Hubbard model – which coincides with the Hubbard model for special values of the Yukawa couplings (S0.Ex45) and the sources. Particularly simple modifications arise for $`\stackrel{}{x}`$-independent and $`\tau `$-independent sources. As an example, the source
$$l_{3\rho }^{}=l_{3p}=2l_{3q_1}=\frac{1}{4}\nu $$
(17)
induces an additional energy for the occupation of the sites $`(m,n)`$ with both $`m`$ and $`n`$ odd
$$S_\nu =\nu 𝑑\tau \underset{\genfrac{}{}{0pt}{}{m\mathrm{odd}}{n\mathrm{odd}}}{}\widehat{\psi }_{mn}^{}\widehat{\psi }_{mn}.$$
(18)
For $`\nu \mathrm{}`$ these sites are effectively removed from the lattice and we therefore deal with the Hubbard model on a non-quadratic lattice structure.
Analytic computations for the partition function (S0.Ex36) are most easily done in momentum space. It is straightforward to perform a Fourier transform using
$`\widehat{\psi }_a(\stackrel{}{x},\tau )`$ $`=`$ $`\sqrt{2a}T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2q}{(2\pi )^2}\mathrm{exp}\left(i\{(\stackrel{}{x}+\stackrel{}{z}_a)\stackrel{}{q}+2\pi nT\tau \}\right)\widehat{\psi }_{an}(\stackrel{}{q})}`$
$`\widehat{\psi }_a^{}(\stackrel{}{x},\tau )`$ $`=`$ $`i\sqrt{2a}T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2q}{(2\pi )^2}\mathrm{exp}\left(i\{(\stackrel{}{x}+\stackrel{}{z}_a)\stackrel{}{q}+2\pi nT\tau \}\right)\widehat{\overline{\psi }}_{an}(\stackrel{}{q})\gamma ^0}`$
$`\widehat{\psi }^{}`$ $`=`$ $`i\widehat{\overline{\psi }}\gamma ^0,\gamma ^0=\left(\begin{array}{cc}\tau _3& 0\\ 0& \tau _3\end{array}\right).`$ (21)
Here the Matsubara frequencies are labeled by half integer $`n=\pm 1/2,\pm 3/2,\mathrm{})`$ and the momentum integration is in the range $`\mathrm{\Lambda }<q_x<\mathrm{\Lambda },\mathrm{\Lambda }<q_y<\mathrm{\Lambda }`$ as appropriate for the coarse lattice with lattice distance $`2a`$. We choose
$`\stackrel{}{z}_1=({\displaystyle \frac{a}{2}},{\displaystyle \frac{a}{2}}),\stackrel{}{z}_2=({\displaystyle \frac{a}{2}},{\displaystyle \frac{a}{2}}),`$
$`\stackrel{}{z}_3=({\displaystyle \frac{a}{2}},{\displaystyle \frac{a}{2}}),\stackrel{}{z}_4=({\displaystyle \frac{a}{2}},{\displaystyle \frac{a}{2}})`$ (22)
corresponding to an expansion in the coordinates of the $`(m,n)`$ lattice. This yields for the kinetic term
$`S_{kin}`$ $`=`$ $`{\displaystyle 𝑑\tau \underset{\stackrel{}{x}}{}_{kin}}`$
$`=`$ $`T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2q}{(2\pi )^2}\widehat{\overline{\psi }}_{an}(\stackrel{}{q})P_{ab}^{(0)}(n,\stackrel{}{q})\widehat{\psi }_{bn}(\stackrel{}{q})}`$
$`P^{(0)}`$ $`=`$ $`2\pi nT\gamma ^02it\gamma ^0\{\mathrm{cos}(aq_x)A_1+\mathrm{cos}(aq_y)B_1\}`$ (23)
where we use matrices ($`\tau _0=\mathrm{𝟙}_2`$; $`\mu =0\mathrm{}3`$; $`i,j=1\mathrm{}3`$)
$`A_\mu =\left(\begin{array}{cc}\tau _\mu & 0\\ 0& \tau _\mu \end{array}\right),B_\mu =\left(\begin{array}{cc}0& \tau _\mu \\ \tau _\mu & 0\end{array}\right)`$ (28)
$`\{A_i,B_j\}=2\delta _{ij}B_0,\{A_i,A_j\}=\{B_i,B_j\}=2\delta _{ij}`$
$`[A_i,B_j]=2iϵ_{ijk}B_k`$
$`B_0A_i=A_iB_0=B_i,B_0B_\mu =B_\mu B_0=A_\mu .`$ (29)
Spontaneous symmetry breaking with “extended order parameters” like the antiferromagnetic spin density $`(\stackrel{}{\stackrel{~}{a}}_1\stackrel{}{\stackrel{~}{a}}_2+\stackrel{}{\stackrel{~}{a}}_3\stackrel{}{\stackrel{~}{a}}_4)`$ or $`d`$-wave superconductivity with order parameter $`(\stackrel{~}{d}_1+\stackrel{~}{d}_2+\stackrel{~}{d}_3+\stackrel{~}{d}_4)`$ can be directly investigated in our formalism by looking for the minima of the effective scalar potential. The notion of the effective potential is a very powerful concept since it describes simultaneously situations with vanishing and nonvanishing sources, i.e. in addition to the Hubbard model for arbitrary $`\mu `$ it also comprises many extended models. The effective potential corresponds to the effective action for homogeneous colored scalar fields. We define the scalar expectation values in the presence of sources<sup>9</sup><sup>9</sup>9The variation with respect to $`l_{a\rho }`$ is performed at fixed $`\mu `$.
$`u_{a\beta }(\stackrel{}{x})={\displaystyle \frac{\mathrm{\Lambda }^2}{\pi ^2}}{\displaystyle \frac{}{J_{a\beta }^{}(\stackrel{}{x})}}\mathrm{ln}Z=<\widehat{u}_{a\beta }(\stackrel{}{x})>`$
$`u_{a\beta }^{}(\stackrel{}{x})={\displaystyle \frac{\mathrm{\Lambda }^2}{\pi ^2}}{\displaystyle \frac{}{J_{a\beta }(\stackrel{}{x})}}\mathrm{ln}Z=<\widehat{u}_{a\beta }^{}(\stackrel{}{x})>`$
$`w_{a\gamma }(\stackrel{}{x})={\displaystyle \frac{\mathrm{\Lambda }^2}{\pi ^2}}{\displaystyle \frac{}{L_{a\gamma }(\stackrel{}{x})}}\mathrm{ln}Z=<\widehat{w}_{a\gamma }(\stackrel{}{x})>`$ (30)
with
$$J_{a\beta }=\frac{\mathrm{\Lambda }^2}{\stackrel{~}{h}_\beta }j_{a\beta },L_{a\gamma }=\frac{\mathrm{\Lambda }^2}{h_\gamma }l_{a\gamma }.$$
(31)
With the usual Legendre transform one obtains<sup>10</sup><sup>10</sup>10We concentrate in the following on $`\psi _a=\widehat{\psi }_a=0`$, $`\psi _a^{}=\widehat{\psi }_a^{}=0`$ the effective action $`\mathrm{\Gamma }`$
$`\mathrm{\Gamma }[u,w,\psi ,\psi ^{}]`$ $`=`$ $`\mathrm{ln}Z+{\displaystyle }d\tau {\displaystyle \underset{\stackrel{}{x}}{}}{\displaystyle \underset{a}{}}\{{\displaystyle \frac{\pi ^2}{\mathrm{\Lambda }^2}}{\displaystyle \underset{\beta }{}}(J_{a\beta }^{}u_{a\beta }+J_{a\beta }u_{a\beta }^{})`$ (32)
$`+{\displaystyle \frac{\pi ^2}{\mathrm{\Lambda }^2}}{\displaystyle \underset{\gamma }{}}L_{a\gamma }w_{a\gamma }+\eta _a^{}\psi _a\psi _a^{}\eta _a\}`$
which obeys
$$\frac{\mathrm{\Gamma }}{w_{a\gamma }}=\frac{\pi ^2}{\mathrm{\Lambda }^2}L_{a\gamma }\mathrm{etc}.$$
(33)
Performing the derivatives (S0.Ex63) in the fermionic functional integral, we can directly relate the scalar expectation values to the expectation values of fermionic bilinears
$`u_{a\beta }`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{h}_\beta }{4\pi ^2}}<\stackrel{~}{u}_{a\beta }^{(1)}>+{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}J_{a\beta }`$
$`u_{a\beta }^{}`$ $`=`$ $`{\displaystyle \frac{\stackrel{~}{h}_\beta }{4\pi ^2}}<\stackrel{~}{u}_{a\beta }^{(2)}>+{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}J_{a\beta }^{}`$
$`w_{a\gamma }`$ $`=`$ $`{\displaystyle \frac{h_\gamma }{4\pi ^2}}<\stackrel{~}{w}_{a\gamma }>+{\displaystyle \frac{1}{\mathrm{\Lambda }^2}}L_{a\gamma }.`$ (34)
In particular, if all sources except $`L_{a\rho }(\stackrel{}{x})=\mathrm{\Lambda }^2\mu /h_\rho `$ vanish, the scalar $`w_{a\rho }`$ has contributions linear in the electron density $`n=<\stackrel{~}{\rho }>/4a^2`$ and the chemical potential
$$w_{a\rho }=\frac{h_\rho }{4\mathrm{\Lambda }^2}n+\frac{\mu }{h_\rho }.$$
(35)
We will mainly be interested in homogenous expectation values and therefore in the effective scalar potential which can be obtained from $`\mathrm{\Gamma }`$ for $`\stackrel{}{x}`$\- and $`\tau `$-independent scalar fields and vanishing fermion fields by
$$U_0=T\mathrm{\Gamma }/V_2.$$
(36)
The ground state of the Hubbard model corresponds to the minimum of $`U_0`$ with respect to all fields except $`\rho =\frac{1}{4}_aw_{a\rho }`$ given by eq. (35), where $`\mu `$ obeys
$$\mu =\frac{h_\rho }{4\mathrm{\Lambda }^2}\frac{U_0}{\rho }.$$
(37)
We are interested in possible expectation values of scalars different from $`\rho `$. Such a spontaneous symmetry breaking arises if for some range of $`\rho `$ the minimum of $`U_0`$ (at fixed $`\rho `$) occurs for a nonvanishing scalar field.
In this paper we compute the effective potential $`U_0`$ in the “mean field” approximation. This means that only the fermionic part of the functional integral (S0.Ex36) is performed in a homogenous “background” $`\widehat{u}=u,\widehat{w}=w`$. This integral is Gaussian, and we can write the mean field expression for $`U_0`$ as
$`U_0`$ $`=`$ $`U_{cl}+\mathrm{\Delta }U`$
$`U_{cl}`$ $`=`$ $`\mathrm{\Lambda }^2{\displaystyle \underset{\beta }{}}{\displaystyle \underset{a}{}}u_{a\beta }^{}u_{a\beta }+{\displaystyle \frac{\mathrm{\Lambda }^2}{2}}{\displaystyle \underset{\gamma }{}}{\displaystyle \underset{a}{}}w_{a\gamma }w_{a\gamma }+{\displaystyle \frac{2\mathrm{\Lambda }^2}{h_\rho ^2}}\mu ^2`$
$`\mathrm{\Delta }U`$ $`=`$ $`{\displaystyle \frac{1}{2}}T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2q}{(2\pi )^2}\mathrm{ln}detP(n,\stackrel{}{q})}.`$ (38)
Here $`P(q)`$ is a $`16\times 16`$ matrix (including spinor indices) for the inverse fermion propagation in the presence of scalar background fields. It is defined by the part of the action quadratic in the fermion fields
$`S_2`$ $`=`$ $`{\displaystyle \frac{1}{2}}T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2q}{(2\pi )^2}\stackrel{~}{\psi }_n^T(\stackrel{}{q})P(n,\stackrel{}{q})\stackrel{~}{\psi }_n(\stackrel{}{q})}`$ (39)
$`\stackrel{~}{\psi }_n(\stackrel{}{q})`$ $`=`$ $`\left(\begin{array}{c}\widehat{\psi }_{a,n}(\stackrel{}{q})\\ \widehat{\overline{\psi }}_{a,n}(\stackrel{}{q})\end{array}\right)`$ (42)
and we find from eq. (S0.Ex60)
$$P=\left(\begin{array}{cc}0& P_0^T(n)\\ P_0(n)& 0\end{array}\right)+\stackrel{~}{P}(\rho ,d,\stackrel{}{a},\mathrm{})$$
(43)
where the second term reflects the influence of the background through the Yukawa couplings. We explore here the dependence of the effective potential on the charge density $`\rho `$, $`d`$-wave pair condensation $`d`$ and antiferromagnetic order parameter $`\stackrel{}{a}`$. We therefore take
$`w_{1\rho }=w_{2\rho }=w_{3\rho }=w_{4\rho }=\rho `$
$`u_{1d}=u_{2d}=u_{3d}=u_{4d}=d`$
$`\stackrel{}{w}_{1a}=\stackrel{}{w}_{2a}=\stackrel{}{w}_{3a}=\stackrel{}{w}_{4a}=\stackrel{}{a}`$ (44)
and find<sup>11</sup><sup>11</sup>11The choice of $`\gamma ^0`$ in (S0.Ex56) is not crucial for this calculation – any orthogonal $`\gamma ^0`$ will do.
$`\stackrel{~}{P}`$ $`=`$ $`ih_\rho \rho \left(\begin{array}{cc}0& \gamma _{}^{0}{}_{}{}^{T}\\ \gamma ^0& 0\end{array}\right)ih_a\stackrel{}{a}\left(\begin{array}{cc}0& A_3\gamma _{}^{0}{}_{}{}^{T}\stackrel{}{\tau }^T\\ \gamma ^0A_3\stackrel{}{\tau }& 0\end{array}\right)`$ (52)
$`h_d\left(\begin{array}{c}d^{}[\mathrm{cos}(aq_x)A_1\mathrm{cos}(aq_y)B_1]0\\ 0d\gamma ^0[\mathrm{cos}(aq_x)A_1\mathrm{cos}(aq_y)B_1]\gamma _{}^{0}{}_{}{}^{T}\end{array}\right)c`$
$`=`$ $`\stackrel{~}{P}^T.`$ (53)
With $`G=\mathrm{diag}(\mathrm{𝟙},i\gamma _{}^{0}{}_{}{}^{T})`$ and $`\underset{¯}{P}=GPG^T`$ one obtains
$`\mathrm{ln}detP`$ $`=`$ $`\mathrm{ln}det\underset{¯}{P}={\displaystyle \frac{1}{2}}\mathrm{ln}det\left(\underset{¯}{P}B_0\underset{¯}{P}^TB_0\right)`$ (54)
$`=`$ $`\mathrm{ln}\mathrm{det}_8\{(2\pi nT)^2+[2t\mathrm{cos}(aq_x)A_1+2t\mathrm{cos}(aq_y)B_1+h_\rho \rho ]^2`$
$`+h_a^2\stackrel{}{a}\stackrel{}{a}+2h_\rho h_a\rho \stackrel{}{a}A_3\stackrel{}{\tau }`$
$`+h_d^2d^{}d[\mathrm{cos}(aq_x)A_1\mathrm{cos}(aq_y)B_1]^2\}.`$
It is easy to see that $`\mathrm{\Delta }U`$ depends only on the invariants $`\delta =d^{}d,\alpha =\stackrel{}{a}\stackrel{}{a}`$ and $`\rho `$. Up to an additive ($`T`$-dependent) constant one finds
$`\mathrm{\Delta }U`$ $`=`$ $`{\displaystyle \frac{1}{2}}T{\displaystyle \underset{n}{}}{\displaystyle }{\displaystyle \frac{d^2q}{(2\pi )^2}}\mathrm{tr}_8\mathrm{ln}\{\mathrm{𝟙}`$ (55)
$`+([2t\mathrm{cos}(aq_x)A_1+2t\mathrm{cos}(aq_y)B_1+h_\rho \rho +h_a\sqrt{\alpha }A_3\tau _3]^2`$
$`+h_d^2\delta [\mathrm{cos}(aq_x)A_1\mathrm{cos}(aq_y)B_1]^2)/(2\pi nT)^2\}`$
$`U_{cl}`$ $`=`$ $`2\mathrm{\Lambda }^2\alpha +4\mathrm{\Lambda }^2\delta +2\mathrm{\Lambda }^2\rho ^2+{\displaystyle \frac{2\mathrm{\Lambda }^2\mu ^2}{h_\rho ^2}},`$
and, evaluating the Matsubara sum and the trace, finally
$`U_0`$ $`=`$ $`2\mathrm{\Lambda }^2\alpha +4\mathrm{\Lambda }^2\delta +{\displaystyle \frac{2\mathrm{\Lambda }^2\mu ^2}{h_\rho ^2}}2T{\displaystyle \frac{d^2q}{(2\pi )^2}\underset{\epsilon _i,\epsilon _j\{1,1\}}{}}`$
$`\mathrm{ln}\mathrm{cosh}\left({\displaystyle \frac{1}{2T}}\sqrt{\left(h_\rho \rho +\epsilon _i\sqrt{4t^2(c_x+\epsilon _jc_y)^2+h_a^2\alpha }\right)^2+h_d^2\delta (c_x\epsilon _jc_y)^2}\right)`$
with $`c_x=\mathrm{cos}(aq_x)`$, $`c_y=\mathrm{cos}(aq_y)`$.
For large temperature the fluctuation contribution $`\mathrm{\Delta }U`$ is suppressed $`T^1`$. The minimum of $`U_0`$ therefore occurs for all $`\rho `$ at $`\alpha =0,\delta =0`$. As $`T`$ is lowered, the fluctuations tend to destabilize the “symmetric minimum”. In particular, the fluctuation contribution to the mass term for $`\stackrel{}{a}`$ is negative for not too large $`\rho ^2`$ and the one for $`d`$ is negative for all $`\rho `$
$`\mathrm{\Delta }M_a^2`$ $`=`$ $`2{\displaystyle \frac{}{\alpha }}\mathrm{\Delta }U_{|\alpha =\delta =0}`$
$`=`$ $`2h_a^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2q}{(2\pi )^2}\mathrm{tr}_4\{P_\rho ^22h_\rho ^2\rho ^2A_3P_\rho ^2A_3P_\rho ^2\}}`$
$`\mathrm{\Delta }M_d^2`$ $`=`$ $`{\displaystyle \frac{}{\delta }}\mathrm{\Delta }U_{|\alpha =\delta =0}`$ (56)
$`=`$ $`h_d^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{d^2q}{(2\pi )^2}\mathrm{tr}_4\{(\mathrm{cos}(aq_x)A_1\mathrm{cos}(aq_y)B_1)^2P_\rho ^2\}}`$
with
$$P_\rho ^2=(2\pi nT)^2+(2t\mathrm{cos}(aq_x)A_1+2t\mathrm{cos}(aq_y)B_1+h_\rho \rho )^2$$
(57)
and $`\mathrm{tr}_4`$ the trace in color space only. These contributions should be compared with $`(M_a^{(0)})^2=(M_d^{(0)})^2=4\mathrm{\Lambda }^2`$. The zeroes of $`P_\rho ^2`$ for $`T=0`$ correspond to the Fermi surface (S0.Ex2) with shifted chemical potential $`\mu _{\mathrm{eff}}=h_\rho \rho `$. (Neglecting contributions from $`\mathrm{\Delta }U`$ eqs. (37) and (S0.Ex68) imply $`\mu _{\mathrm{eff}}=\mu `$.) We recall that the momenta are restricted to the range corresponding to the coarse lattice $`|q_{x,y}|\pi /(2a)=\mathrm{\Lambda }`$. On the other hand we have now possible zeroes for different linear color combinations. Noting that the eigenvalues of $`A_1`$ and $`B_1`$ are $`\pm 1`$, they precisely correspond to the original Fermi surface – the original zeros in the four ranges $`|q_{x,y}|\mathrm{\Lambda },\mathrm{\Lambda }|q_{x,y}|2\mathrm{\Lambda },|q_x|\mathrm{\Lambda },\mathrm{\Lambda }|q_y|2\mathrm{\Lambda }`$ and $`\mathrm{\Lambda }|q_x|2\mathrm{\Lambda },|q_y|\mathrm{\Lambda }`$ appear now for different color combinations in the range $`|q_{x,y}|\mathrm{\Lambda }`$. Due to these zeros one finds $`lim_{T0}\mathrm{\Delta }M_d^2\mathrm{}`$ for not too large $`\rho `$ and similar for $`\mathrm{\Delta }M_a^2`$ in the appropriate range of $`\rho `$. This clearly indicates spontaneous symmetry breaking with $`d`$–wave superconductivity or/and antiferromagnetic order parameter. Note that in contrast to its derivatives the potential is not singular for $`T0`$. Since for large $`\alpha `$ and $`\delta `$ $`U_0`$ grows $`(\alpha ,\delta )`$, the minimum occurs necessarily for finite $`\alpha 0`$ or $`\delta 0`$ if $`M_a^2=4\mathrm{\Lambda }^2+\mathrm{\Delta }M_a^2`$ or $`M_d^2=4\mathrm{\Lambda }^2+\mathrm{\Delta }M_d^2`$ become negative.
The spontaneous symmetry breaking cuts off the singularity near the Fermi surface or reduces its strength. An antiferromagnetic expectation value typically produces a gap for the fermionic fluctuations. For $`\alpha >0,\delta =0,\rho =0`$ this can be seen from a search for possible zeroes of $`det_8`$ in eq. (54) for $`T=0`$. On the other hand, for $`\alpha =0,\delta >0`$ the condition $`det_8=0`$ requires $`\mathrm{cos}(aq_x)=\pm \mathrm{cos}(aq_y)=\pm h_\rho \rho /(4t)`$. In the superconducting phase the singularity therefore only occurs for special points in momentum space instead of a whole Fermi surface. As a consequence, the momentum integrations for the bosonic mass terms (similar to eq. (S0.Ex81)) remain finite even for $`T0`$.
For vanishing sources $`j_a,j_d`$ the minimum of $`U_0`$ obeys the “field equations”
$`{\displaystyle \frac{U_0}{\stackrel{}{a}}}`$ $`=`$ $`2\stackrel{}{a}{\displaystyle \frac{U_0}{\alpha }}=2\stackrel{}{a}(2\mathrm{\Lambda }^2{\displaystyle \frac{1}{2}}h_aT{\displaystyle \underset{n}{}}{\displaystyle }{\displaystyle \frac{d^2q}{(2\pi )^2}}`$
$`\mathrm{tr}_8\left\{(h_\rho \rho \alpha ^{1/2}A_3\tau _3+h_a)\overline{P}_\rho ^2(\alpha ,\delta )\right\})=0`$
$`{\displaystyle \frac{U_0}{d^{}}}`$ $`=`$ $`d{\displaystyle \frac{U_0}{\delta }}=d(4\mathrm{\Lambda }^2{\displaystyle \frac{1}{2}}h_d^2T{\displaystyle \underset{n}{}}{\displaystyle }{\displaystyle \frac{d^2q}{(2\pi )^2}}`$
$`\mathrm{tr}_8\{[\mathrm{cos}(aq_x)\mathrm{cos}(aq_y)B_0]^2\overline{P}_\rho ^2(\alpha ,\delta )\})=0`$
with
$`\overline{P}_\rho ^2(\alpha ,\delta )`$ $`=`$ $`P_\rho ^2+2h_\rho h_a\rho \sqrt{\alpha }A_3\tau _3`$ (60)
$`+h_a^2\alpha +h_d^2\delta \left[\mathrm{cos}(aq_x)\mathrm{cos}(aq_y)B_0\right]^2.`$
One always has the symmetric solution $`\stackrel{}{a}=0,\delta =0`$ which corresponds to a local minimum if $`\overline{M}_a^2>0,\overline{M}_d^2>0`$ and to a maximum or saddlepoint otherwise. Consider next solutions with $`\stackrel{}{a}=0,\delta 0`$ which require
$`h_d^4\delta T{\displaystyle \underset{n}{}}{\displaystyle }{\displaystyle \frac{d^2q}{(2\pi )^2}}\mathrm{tr}_4\{[\mathrm{cos}(aq_x)\mathrm{cos}(aq_y)B_0]^4P_\rho ^2`$
$`(P_\rho ^2+h_d^2\delta [\mathrm{cos}(aq_x)\mathrm{cos}(aq_y)B_0]^2)^1\}`$
$`=4\mathrm{\Lambda }^2+\mathrm{\Delta }M_d^2`$ $`=`$ $`\overline{M}_d^2.`$ (61)
Solutions with $`\delta >0`$ exist only for $`\overline{M}_d^2<0`$ and $`\delta `$ vanishes as the mass term $`\overline{M}_d^2`$ approaches zero from below. One concludes that the transition from the symmetric phase $`(\stackrel{}{a}=0,\delta =0)`$ to a possible superconducting phase without antiferromagnetism $`(\stackrel{}{a}=0,\delta >0)`$ is of second order.
We have analyzed the phase diagram for different Yukawa couplings numerically. Due to the free parameters $`\lambda _i`$ in eq. (S0.Ex45) the Yukawa couplings are largely undetermined. They only have to obey the inequalities $`h_\rho ^2>0,h_d^2>0,h_a^2>0,h_a^2>\pi ^2U/3+h_\rho ^2/32h_d^2/3.`$ For example, $`\lambda _1=\lambda _3=1/2,\lambda _2=1`$ leads to $`h_d^2=h_\rho ^2=h_a^2=\pi ^2U/2`$. Because of our meanfield approximation, the partition function becomes dependent on the particular choice of the parameters $`\lambda _i`$. We investigate the cases $`h_\rho =h_\alpha =h_\delta =\sqrt{10U}`$ (fig. 2), $`h_\rho =h_\alpha =h_\delta =2\sqrt{10U}`$ (fig. 3), $`h_\alpha =\sqrt{10U}`$, $`h_\rho =h_\delta =2\sqrt{10U}`$ (fig. 4), and $`h_\delta =\sqrt{10U}`$, $`h_\rho =h_\alpha =2\sqrt{10U}`$ (fig. 5). We choose $`t/U=1`$ and investigate the phase diagram in the $`\rho /\sqrt{U}`$-$`T/U`$-plane. Expressed in the variables $`t/U`$, $`T/U`$, $`\rho /\sqrt{U}`$, our results do not depend on $`U`$ and the lattice distance $`a`$, as discussed in the beginning. As we increase all three Yukawa couplings simultaneously, the antiferromagnetic phase dominates over the superconducting phase (compare figs. 2 and 3). An interesting result of the mean field analysis is the appearance of a phase transition of first order into the antiferromagnetic phase for small $`T/U`$ and for high values of $`\rho /\sqrt{U}`$. The phase transition between the symmetric and the superconducting phase remains of second order. Both results were anticipated when examining the above formulae analytically. If we increase $`h_d/\sqrt{U}`$ compared to $`h_a/\sqrt{U}`$ the superconductivity phase dominates for low $`T/U`$, whereas in the opposite case it is the antiferromagnetic phase. This is illustrated in figures 4 and 5.
We note that for negative $`t`$ our results apply if the antiferromagnetic condensate $`\stackrel{}{a}`$ is replaced by the ferromagnetic condensate $`\stackrel{}{m}`$. Furthermore, small disturbances can easily be taken into account by source terms. For example, an interaction between spin and angular momentum will explicitely break the continuous $`SU(2)`$ invariance and typically amount to a source term $`l_\stackrel{}{a}`$ or $`l_\stackrel{}{m}`$.
In conclusion, the mean field approximation for the colored Hubbard model can give a qualitatively reasonable picture of the phases in high $`T_c`$ superconductors. On the other hand, the shortcomings of this approximation are also apparent from the figures. All phase diagrams in figs. 2, 3, 4 and 5 correspond to different mean field approximations for the same model. It is impossible to resolve this ambiguity within the mean field approximation without additional input on the selection of the Yukawa couplings. The reason is the neglect of fluctuations of the bosonic fields. Only if these are included, the different equivalent choices of the Yukawa couplings should lead to the same physical results. The differences between the figures reveal the importance of the neglected bosonic fluctuations, at least for some choices of the Yukawa couplings<sup>12</sup><sup>12</sup>12It is conceivable that an ”optimal choice” of the Yukawa couplings minimizes the impact of the bosonic fluctuations..
The inclusion of the bosonic fluctuations is a complex problem which can be attacked by means of nonperturbative renormalization group equations . Studies for similar QCD-motivated models of fermions with Yukawa coupling to scalars have already been carried out successfully . One of the dominant effects will be the scale dependence of the Yukawa couplings. It is conceivable that this running is dominated by partial infrared fixed points for ratios of Yukawa couplings. For large couplings, as relevant here, such partial fixed points would be approached fast. In this case the “memory” of the initial choice of Yukawa couplings could be erased rapidly and unambiguous physical predictions become possible.
A second important ingredient is the appearance of Goldstone bosons for $`\stackrel{}{a}0`$ or $`d0`$, corresponding to flat directions in the effective potential (S0.Ex68). For a superconducting condensate $`d`$ the $`U(1)`$-symmetry would be spontaneously broken and the question arises if this is self consistent. For a large correlation length $`\xi `$, i.e. $`\xi T1`$, one expects that the dominant fluctuations near a second order phase transition are well described by an effective dimensional reduction to two dimensional classical statistics. The Mermin–Wagner theorem then suggests that the Goldstone boson fluctuations prevent a continuous symmetry from being spontaneously broken. In the case of a $`U(1)`$-symmetry the natural solution to this puzzle is a second order phase transition of the Kosterlitz–Touless type: only a renormalized expectation value differs from zero, whereas the wave function renormalization will lead to a vanishing expectation value for the unrenormalized scalar field . This reconciles the Mermin–Wagner theorem with the existence of Goldstone bosons and superconductivity in presence of electromagnetic fields.
For a possible “antiferromagnetic phase” the nonabelian interactions between the Goldstone bosons of the effective two dimensional model have a tendency to push the minimum of $`U_0`$ towards $`\alpha =0`$ and to make $`U_0/\alpha `$ positive . If only the nonabelian Goldstone bosons are present in the effective long distance model their fluctuations would destroy the nontrivial minimum of the potential. One may therefore speculate about a new type of low temperature phase, which is characterized by the presence of massless Goldstone bosons as well as massless fermions. Alternatively no true antiferromagnetic phase with Goldstone bosons may occur. For all practical purposes the physics nevertheless will look qualitatively similar to the phase transition in the mean field approximation: the effects from Goldstone fluctuations are only logarithmic in ratios of mass scales and would be cut off by a small $`SU(2)`$-breaking disturbance inducing a mass term for them. Simple scale considerations suggest that the first order transitions to the antiferromagnetic phase are probably not affected substantially by the Goldstone fluctuations, except for the endpoints. Particularly interesting is the triple point in fig. 2 where the three phases meet. By continuity of the second order lines one expects five massless scalar excitations at this point.
We emphasize that quite generally the possible second order phase transitions between the symmetric and some other phase belong to new interesting universality classes. Long range fermion fluctuations without a gap are present in the symmetric phase and therefore also at the phase transition. They influence the critical exponents and other universal properties.
We hope that our formulation of the colored Hubbard model will be a good starting point for a quantitative renormalization group study of all these interesting questions.
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# The descendents of Lyman break galaxies in galaxy clusters: spatial distribution and orbital properties
## 1. Introduction
The implementation of a simple color selection technique to select efficiently galaxies at redshift larger than 2.5 (Steidel et al. 1996, Madau et al. 1996, Steidel et al. 1999, Fontana et al. 1999 ) revealed a population of blue, actively star forming galaxies at high redshift. Galaxies as bright as those observed are likely hosted inside the most massive halos at high z (Bagla 1998, Baugh et al. 1998, Katz, Hernquist & Weinberg 1999, Coles et al. 1998, but see Somerville, Primack & Faber 1998 and Kolatt et al. 1999 for a slightly different view). These halos are more clustered (a bias factor of the order of 2-5) compared to the general distribution, providing strong support (Adelberger et al. 1998, Giavalisco et al. 1998) to models of biased galaxy formation (Davis et al. 1985). Under the effect of gravitational instability these large halos will merge together and form the massive clusters we see today (Governato et al. 1998). Semi–analytical models (Baugh et al. 1998) further suggested that the present day descendents of LBG in protoclusters would be preferentially giant ellipticals with an old red population of stars.
Bright, red cluster members reside preferentially at the center of clusters and often have been found to have a lower orbital velocity dispersion (Chincarini & Rood 1977, Mellier et al. 1988, Biviano et al. 1992, Whitmore et al. 1993, Biviano et al. 1996, Carlberg et al. 1997) than the global cluster population. Recent results with full redshift information for a large sample of clusters (Adami, Biviano & Mazure 1998, Ramirez, de Souza & Schade 2000) and photometric observations of the Coma cluster (Kashikawa 1998) give support to these claims. Adami et al. , based on a simple theoretical modeling, suggest that orbits of the brightest galaxies have to be circular to explain the decrease in velocity dispersion and at the same time be consistent with the hypothesis of dynamical equilibrium at the cluster center.
Indeed theoretical prejudice would expect to find galaxies formed in massive halos at high redshift to reside preferentially in the central region of clusters. Moore et al. (1998) and White & Springel (1999) showed that, in CDM cosmologies, matter already in virialized objects at high redshift makes a large fraction of the mass within the central regions of present day clusters. Dynamical friction, if acting efficiently on a long enough time scale could further segregate massive halos at the center of clusters (but the effect is likely to be small; see Ghigna et al. 1998 and Colpi, Mayer & Governato 1999, hereafter CMG99). In recent years, numerical and analytical studies of galaxy clusters have rapidly increased in resolution and detail (e.g. Katz & White 1993, Carlberg 1994, Frenk et al. 1996, Fusco–Femiano & Menci 1998, Tormen, Diaferio & Syer 1998, Klypin et al. 1999). Frenk et al. (1996) found mild spatial segregation of the most massive galaxies, they also included gas dynamics and a simple description of star formation processes.
In this work we use a unified approach that couples a state of the art numerical simulation of a galaxy cluster with a detailed, semi–analytical description of galaxy formation inside individual dark matter halos. This method will allow us to study the spatial and orbital distribution of galaxies in a moderately rich cluster with unprecedented detail and to investigate the relation between LBGs at high redshift and present day cluster galaxies.
## 2. Coupling N-body simulations and semi–analytical models
### 2.1. The Cluster simulation
We used a very high resolution N-body (i.e. collisionless) simulation of a galaxy cluster (slightly more massive than Virgo, $`2.3\times 10^{14}h^1M_{}`$ within the virial radius, defined as the radius where $`\rho (r<R)200\rho _{crit}`$) formed in a cluster normalized SCDM cosmology. This cluster contains over 4 $`\times `$ 10<sup>6</sup> particles within the virial radius (it is described in full detail in Ghigna et al. 1998, G98, Ghigna et al. 1999, G99, and Lewis et al. 1999). The effective spatial resolution is of the order of $`1.0h^1\mathrm{kpc}`$, and we are able to resolve substructure halos with circular velocities V<sub>c</sub> down to 50 km/sec and with pericenters larger than $`50h^1\mathrm{kpc}`$, a significant improvement compared with all previous works (here V<sub>c</sub> is defined as $`\sqrt{GM(<r)/r}`$)).
The cluster forms through major mergers at redshift about 0.5 (defined as when its main progenitor has roughly 50% of the final cluster mass), accreting additional mass and galaxies at later times. It is well virialized and close to dynamical equilibrium by the present time (see G98 fig.1). In this high resolution simulation, numerical overmerging (e.g., Moore, Katz & Lake 1996) is likely to be almost negligible, especially for the most massive halos. We can follow the evolution of thousands of halos as they participate in the build up of the cluster and subsequently orbit inside it. Even if severely stripped by the cluster tidal field, virtually all halos maintain their identity once inside the cluster, and only a few get destroyed by the tidal field or decay by dynamical friction at its center (see G99 for a full discussion).
Therefore, within this simulation is possible to follow the descendents of all halos and specifically those associated with LBG (see next subsection) through subsequent outputs to the present time. We first located dark matter halos inside the clusters with the algorithm SKID (see G98 and G99 for details, Springel 1999 for an alternative method) at the final time of the simulation and traced them back to dark matter halos at z = 3. High z halos were identified with ”Friends–of–Friends” (FOF, Davis et al. 1985), using a linking length $`=0.2`$ the initial grid spacing, as FOF gives more robust masses for halos outside larger virialized structures. At $`z=3`$ the region containing the cluster has yet to collapse, but hundreds of smaller halos have already formed within a complex network of filaments. For each halo we measured mass and circular velocity at their virial radius at high z and at their individual tidal radii as imposed by the cluster potential at the present time (again see G98 for details).
### 2.2. Semi-analytical galaxy formation
The growth of dark matter halos can be followed both with N–body simulations and the extended Press & Schechter approach (or PS, see Press & Schechter 1974, Bower 1991, Bond et al. 1991). Within the semi–analytical approach a simple set of equations then describes the cooling of gas inside the dark matter halos and the subsequent star formation history, predicting size, luminosities (including the effects of dust), colors and the morphology of galaxies formed inside these halos (Kauffmann, White & Guiderdoni 1993, Cole et al. 1994, Cole et al. 2000). We used the approach first outlined in Cole et al. 1994 and then further developed in Baugh et al. 1998 and Cole et al. 2000 This approach is remarkably successful in predicting the main properties of high redshift and local galaxies with a minimum set of constrained parameters (Baugh et al. 1998, see also Somerville & Primack 2000). However, the semi-analytical method based on the PS alone cannot recover the full 3D distribution of galaxies, making the full N-body simulations necessary.
For each halo identified at redshift 3 in the N–body simulation and for the whole cluster at the present time we used the semi-analytical approach (and so a PS merging history) to determine their galaxy content. At the present time the halos present inside the cluster are paired with semi-analytical cluster galaxies based on their circular velocity. Statistically this is equivalent to using merger trees obtained directly from N-body simulations (see Governato et al. 1998, Benson et al. 1999). As a test we looked at a set of different merger tree histories for a few clusters of the same mass as the one in our simulation to verify that the scatter introduced by our approach on the average properties as a function of circular velocities of the galaxy population was negligible. In fact, the galaxies produced with the semi-analytical approach show a rather tight luminosity-circular velocity relation, independently of the details of their merging history or of that of the parent cluster. This simple approach is then quite adequate for our purpose of broadly defining the types of galaxies hosted both in massive halos at high z and inside their cluster descendents at the present time as a function of their mass (see Kauffmann et al. 1999, Springel 1999 for an approach based on the full merger trees obtained from N–body simulations).
Our procedure gives the properties of the galaxies hosted inside each given halo complete with full dynamical information (position and velocity inside the cluster). Once galaxies were placed inside dark matter halos, we selected at redshift of 3 those that, applying the same criteria, would have been selected as LBGs (Steidel et al. (1996). At the final time we then compare the properties of their descendents vs. those of the 20 brightest cluster members (comparable to the number of redshift usually measured for a single real cluster) and the whole cluster galaxy population.
## 3. Results: The descendents of Lyman Break Galaxies
At a redshift of 3 there are 12 halos with mass above 10<sup>12</sup> M (the biggest object in the region that will later form the cluster has a mass of $`3.2\times 10^{12}h^1M_{}`$). The semi-analytical approach predicts that each of these halos hosts at least one LBG galaxy, sometimes two. Indeed the N-body simulation already shows significant substructure inside them. There is some intrinsic scatter from one semi-analytical realization to another, depending on the details of the merging histories of individual halos. Sometimes smaller halos (on average less than one per realization) host LBGs, perhaps “observed” while they were at their maximum luminosity. This does not change our results significantly.
Halos containing LBGs are aligned along filaments and are rapidly flowing along them to form massive groups at z $``$ 1.5–0.75 and then merge to form the main progenitor of the cluster by z$`=`$0.5. A large fraction (7 out of 12) merge together to form the central core of the cluster; 90% of the mass contained in their central part (defined as particles within the central $`10h^1\mathrm{kpc}`$) and likely tracers of their stellar component ends up in the inner $`125h^1\mathrm{kpc}`$ of the cluster. Their barionic cores (not present in our simulation that includes only the dark, collisionless component) would then most likely merge together to form the central cD, as the decay time for any remnant of significant mass with orbital apocenters less than $`100h^1\mathrm{kpc}`$ from the cluster center is much shorter than the Hubble time. The five surviving halos have been tidally stripped and orbit in the central part of the cluster. According to the semi–analytical approach the descendents of LBGs are the most luminous ellipticals in the cluster at the present day. This result is independent of the details of the semi–analytical model used. In the approach used e.g. in Kolatt et al. (1999) a large number of Lyman Break Galaxies are small starbursting galaxies. These strong episodes of star formation originate from fly-bys between satellites inside more massive halos. In principle, these satellites could be stripped away from their parent halos and show a different spatial bias, making our conclusions dependent on the analytical modelling.
However, our simulation shows that none of the satellites of the massive halos at z = 3 survive as distinct entities by the present time, having merged with their hosts before the formation of the main body of the cluster.
### 3.1. Orbits and luminosity segregation
At $`z=0.1`$ all LBG descendents can be found within the inner $`0.6h^1\mathrm{Mpc}`$, i.e 60% of the virial radius of the cluster. They are more concentrated than the average cluster population (see Fig. 1). This is more evident in the distribution of the pericentric distances (inset of Fig. 1) and is true for apocenters as well. To strengthen the significance of the signal, we have verified that this holds true at a nearby epoch ($`z=0`$). Using Wilcoxon test, we estimate that the probability of this spatial segregation happening by chance is less than 2%. Also considering that seven LBGs contributed to form the central galaxy, the mass contributed to the cluster by LBG descendents is more centrally concentrated compared with the global cluster population. As halos with large circular velocities are associated with galaxies of higher luminosities than the average galaxy cluster population, this causes a mild luminosity segregation. It is likely that the early formation time of this cluster and its following quiet merging history (it forms slightly earlier than average for its mass in a SCDM cosmology; G98) contributed to this, as recent infall was not substantial enough to accrete massive halos at the outskirts of the cluster.
We then measured the orbital parameters for all galaxies inside our cluster (see Fig. 2). The orbital circularity $`ϵ`$ is defined as the ratio of J, the orbital angular momentum, to J<sub>c</sub>, the angular momentum of a circular orbit with the same energy. (Lacey & Cole 1994, Tormen 1997, G98). Here the orbital energy is defined assuming spherical symmetry for the cluster mass distribution and the most bound particles for its center. There is no obvious difference in the distribution of circularities between descendents of LBG, the twenty brightest objects in the cluster and the rest of the galaxy population. The formal average values of $`ϵ`$ for the three cases are $`0.42`$, $`0.59`$ and $`0.54`$, respectively, with quite similar dispersions around the mean value ($`0.3`$). The dark matter background has similar orbital properties (G98). Results at $`z=0.1`$ and $`0`$ are similar.
This finding confirms results obtained with analytical and numerical models (van den Bosch et al. 1998, CMG99) that dynamical friction is not efficient at circularizing orbits of even the most massive and old galaxies inside clusters. We used the theory of linear response as described in CMG (which agrees extremely well with N-body experiments) to measure the orbital decay predicted for a group sized halo entering the cluster environment at $`z=1`$ (i.e. the formation time of the main progenitor of the cluster itself). Once the effect of tidal stripping are included (see again CMG99) decay times are of the order of several times the Hubble time, and both pericenters and apocenters have changed only by a few percent by the present time. The luminosity segregation putatively observed in real galaxy clusters would then be an imprint of their hierarchical build up rather than the effect of subsequent strong dynamical evolution. This orbital segregation should be present (Springel, 2000 in preparation) or could even be larger in clusters formed in a open or flat cosmology, where clusters would form typically at higher redshift and where the accretion at late times slows down considerably.
Our simulation allows us for the first time to test the dynamical mass estimate based on a complete sample of substructure halos. We estimated the virial mass of the cluster from the galaxies’ projected velocity dispersions, using the classic estimator (Heisler, Bachall & Tremaine, 1985):
$`M_{VT}=(3\pi N/2G)\frac{{}_{i}{}^{}v_{p,i}^{2}}{_{i<j}R_{ij}^1}`$
where v<sub>p,i</sub> is the line of sight velocity and R<sub>ij</sub> the projected separation of a given galaxy pair. This estimator is useful for its simplicity, even if it overestimates the mass inside the virial radius by about 40% (see also Girardi et al. 1998 and references therein). We do not include galaxies in halos outside the virial radius, so that our sample is free of nearby back/foreground interlopers. Our results confirm previous results (Frenk et al. 1996, Tormen 1997) that the use of only the few brightest galaxies as mass estimators results in an underestimate of the cluster mass compared to using the whole galaxy sample, by up to a factor of 2 if the brightest galaxy is included (it has a very small velocity compared to the cluster galaxies as a whole). Additional scatter ($``$30%) is added when considering individual axial projections. Contrary to previous suggestions, this bias is not due to the most massive galaxies being on more circular orbits, but the fact that these galaxies sample only the central part of the cluster mass distribution and therefore have a lower velocity dispersions, as the peak in V<sub>c</sub> for the cluster as a whole is reached only at about $`0.5h^1\mathrm{Mpc}`$, i.e. close to the apocenters of their orbits. Even excluding contamination from background and foreground objects at least 20 galaxies are needed to correctly sample the cluster potential and obtain a reliable mass measurement. Likely, even more redshifts would be needed in case the cluster had significant nearby structures (filaments or rich groups) in the near fore/background.
## 4. Discussion
Using the high resolution simulation of a galaxy cluster coupled with semi–analytical methods of galaxy formation we identify at redshift of 3 a dozen halos hosting at least one Lyman Break Galaxy. At the present time descendents of LBGs can be identified with the central cD galaxy and galaxies hosted in substructure halos with V<sub>c</sub> in the range 200 to 550 km/sec. All 12 LBG descendents end up within the the inner $`0.5h^1\mathrm{Mpc}`$ (or 60% of the cluster virial radius); 7 merged together to form the central galaxy. These descendents are the most bright elliptical galaxies in the cluster. These results are largely independent from the details of the semi-analytical method used. We confirm previous findings obtained with simulations of lower resolution (e.g Frenk et al. 1996) that the most massive galaxies are likely to be centrally segregated and have lower orbital velocity dispersions when compared to the global cluster galaxy population. However, this effect is small, and harder to detect when only limited information (redshifts and positions projected on the sky plane) is available.
Spatial and velocity segregation for bright cluster members has long been observed in Coma (Mellier et al. 1988) and in larger samples of nearby clusters (Biviano et al. 1992, 1996), but the observational picture has been somewhat complicated by the small number of redshift available per cluster and by the fact that they have usually been collected only within the central part of the clusters themselves. Clearly a larger sample of observed and simulated clusters is needed to allow a more quantitative comparison between observations and theoretical predictions. We expect the segregation of bright ellipticals to be larger in well virialized clusters and in cosmologies were recent infall is small (e.g. open or flat CDM cosmologies).
Galaxies in our simulated cluster move on quite eccentric orbits, due to the almost radial infall typical of hierarchical clustering. There is no significant difference in the orbital eccentricity of different galaxy populations and the dark matter background. Also, orbits do not change in shape significantly over time (dynamical friction does not change the orbital eccentricity as shown also in CMG99).
Our analysis shows that to measure the virial mass of the cluster is crucial that a significant number of galaxies is used in order to correctly sample the cluster density profile. A sample, restricted to the most bright cluster members is likely to be biased and underestimate both the cluster total mass and velocity dispersion. In this simulated cluster, about twenty member galaxies sampling the mass distribution out to the virial radius are required for a correct estimate of the cluster total mass. This number could well be higher for a cluster far from virial equilibrium or with significant structures nearby. As our analysis shows, virial mass estimates suffer from an additional scatter of about 30%, due to velocity anisotropies along the cluster projection. This source of scatter cannot easily be removed increasing the number of galaxies.
As clusters likely formed only a few Gyrs ago, dynamical effects like energy equipartion or dynamical friction are very unlikely to have played any significant role in originating the mass/velocity segregation, especially considering that only a small part of the cluster mass is attached to individual galaxies ($`<15\%`$, see G98). If confirmed by a larger sample of real and simulated clusters, the observed segregation of their more massive galaxies would rather be the signature of their hierarchical build–up.
## 5. Acknowledgements
We thank Carlton Baugh and Andrew Benson for providing us the semianalytical galaxy formation models and acknowledge useful discussions with Frank van den Bosch, Rychard Bouwens and Cedric Lacey. This work was partially supported by the EU Network for Galaxy Formation and Evolution, the NASA/ESS programme, and PPARC. SG is a Marie Curie fellow, while BM is a Royal Society Research Fellow. Simulations were run at the Edinburgh supercomputing center and NCSA.
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# DO–TH 00/11 hep-ph/0005270 Signatures of heavy Majorana neutrinos and HERA’s isolated lepton events
## 1 Introduction
Despite the impressive confirmation of predictions from the standard model (SM) it is general believe that we are on the verge of fundamental new discoveries, be it production of new particles or significant deviations of observables in high–precision measurements. Effects of new physics might also be hidden in existing data sets and it is interesting to see what candidates are able to explain any unexpected events or measurements. A first step in this direction was done in terms of observation of nonvanishing neutrino rest masses, most clearly seen in the up–down asymmetry of the atmospheric muon neutrino flux in SuperKamiokande . The smallness of these masses can be related to massive new particles via the see–saw–mechanism . In this respect it seems most natural to look for effects of massive neutrinos, i.e. search for hints of these new particles in high– or low–energy experiments. The theoretical prejudice is that the neutrinos are Majorana particles — be it because they are delivered by see–saw or pop out of almost every GUT — and we shall follow this idea.
Especially for the case of heavy (few 100 GeV) Majorana neutrinos, production at accelerators has been investigated by many authors. The different possibilities include $`e^+e^{}`$ , $`pp`$ , $`p\overline{p}`$ , $`\nu N`$ , $`ep`$ , linear colliders and $`e^{}e^{}`$ or even $`e\mu `$ machines .
Heavy Majoranas have also been studied within the context of low–energy experiments such as neutrinoless double beta decay (0$`\nu \beta \beta `$) or Kaon decays . The respective Feynman diagrams as well as the concrete model differ in most publications and the interested reader might compare the papers with respect to that.
One of the anomalies in existing data is the existence of high $`p_T`$ isolated leptons together with large missing transverse momentum ($`p/_T`$) at HERA. Since the first event was discovered by H1, five more were found and at least 3 of them can not be explained by $`W`$ production or other SM processes. In contrast to that, ZEUS sees no excess in these events , yet, at the present statistical level, there is no contradiction .
In we examined the process (see Fig. 1)
$$e^+p\overline{\nu _e}\text{ }\alpha ^+\beta ^+X\text{ with }\alpha ,\beta =e,\mu ,\tau $$
(1)
and discussed possible signals of this like–sign dileptons (LSD) and high $`p/_T`$ final state. No such events are reported and previously unavailable direct limits on the elements of the Majorana mass matrix were derived . However, it turns out that when the kinematical cuts used in H1’s search for isolated leptons are applied to our process (1), they tend to ignore one of the two leptons. Especially the requirement of $`p_T^{\mathrm{lepton}}>10`$ GeV is often too much for both charged leptons to fulfill. The LSD signal of Eq. (1) is thus reduced to one isolated lepton with high $`p/_T`$. This possibility can be checked by looking for an additional isolated low $`p_T`$ and/or high pseudorapidity lepton. In addition, it is possible that a produced $`\tau `$ decays hadronically<sup>3</sup><sup>3</sup>3We shall use the term electron, muon or tau for both, particle and antiparticle., resulting also in single lepton final states. More than one isolated jet would be a signal for this kind of event. Since process (1) gives LSD with the same sign as the incoming lepton we concentrate on H1’s positive charged muon events, since there are no positron events found. This fact might be explained if one incorporates also limits on mixing of heavy neutrinos, as derived from 0$`\nu \beta \beta `$. With this constraint the expected $`e`$ signal is smaller than the $`\mu `$ signal.
One might argue that direct Majorana ($`N`$) production via a $`e^+NW`$ vertex is more likely to occur since the cross section is larger. At present there is only an analysis in HERA’s $`e^{}p`$ mode available and it was found that detection is only possible if the $`N`$ decays into $`e^+W^{}`$, giving an isolated lepton with different charge than the incoming one. The reason for that is of course the large background from $`W`$ production. However, a general analysis of all channels ($`N\nu Z`$, $`N\mu W,\mathrm{}`$) remains still to be done and it might be interesting to compare the results with our signals in the future. Until that is done, we think that our process is worth considering, inasmuch as we have no restriction to the flavor of the final state lepton.
The paper is organized as follows: In Section 2 we discuss some general features of the process and the diagram and argue in Section 3 that the two lepton signal of Eq. (1) might very well be seen as a one lepton signal. Section 4 sees a discussion of signals of the events and how one might distinguish the genuine final state from the measured one. Finally, Section 5 closes the paper with a conclusion and discussion.
## 2 The process and heavy Majorana neutrinos
We shall work in a mild extension of the SM with no further specification of how heavy Majoranas might be created. The coupling to the usual leptons and gauge bosons is the familiar left–handed weak interaction. The three known light neutrinos $`\nu _\alpha `$ are thus mixtures of light and heavy mass particles, this can be expressed by the replacement
$$\nu _\alpha \mathrm{cos}\theta _\alpha \nu _\alpha +\mathrm{sin}\theta _\alpha N_\alpha $$
(2)
in the (unmixed) Lagrangian for each family. For the sake of simplicity we take $`N_\alpha =N`$. The Lagrangian now reads:
$$\begin{array}{c}=\frac{g}{\sqrt{2}}W_\mu \left\{\mathrm{cos}\theta _\alpha \overline{\nu _\alpha }\gamma ^\mu \gamma _{}l_\alpha +\mathrm{sin}\theta _\alpha \overline{N}\gamma ^\mu \gamma _{}l_\alpha \right\}\\ +\frac{g}{2\mathrm{cos}\theta _W}Z_\mu \left\{\mathrm{cos}^2\theta _\alpha \overline{\nu _\alpha }\gamma ^\mu \gamma _{}\nu _\alpha +\mathrm{sin}2\theta _\alpha \overline{N}\gamma ^\mu \gamma _{}\nu _\alpha \frac{1}{2}\mathrm{sin}^2\theta _\alpha \overline{N}\gamma _\mu \gamma _5N\right\}+\mathrm{h}.\mathrm{c}.\end{array}$$
(3)
where $`\gamma _{}=\frac{1}{2}(1\gamma _5`$) and there is no vector current between $`\overline{N}`$ and $`N`$ due to their Majorana nature. We can keep it for the light neutrinos since for energies much larger than the (light) masses there is hardly a chance to find a difference between $`2\overline{\nu }\gamma ^\mu \gamma _{}\nu `$ and $`\overline{\nu }\gamma ^\mu \gamma _5\nu `$ .
What can we expect for the values of the masses and the mixing parameters? Taking the typical see–saw formula we find
$$m_\nu \frac{m_D^2}{m_N}m_N\frac{(10^5\mathrm{}10^{11})^2}{10^5\mathrm{}1}\mathrm{eV}100\mathrm{}10^{18}\mathrm{GeV}$$
(4)
where we took for the Dirac mass $`m_D`$ every value from electron to top mass and for the light neutrino mass we allowed everything from the vacuum solution in a highly hierarchical scheme ($`\sqrt{\mathrm{\Delta }m^2}m_\nu 10^5`$ eV) to a degenerate scheme (cosmological or also LSND’s mass scale) $`m_\nu `$ few eV (see for a detailed analysis of allowed schemes). For the mixing angle we have
$$\theta _\alpha \mathrm{sin}\theta _\alpha =\frac{m_D}{M_N}\frac{m_\nu }{m_D}10^5\mathrm{}10^{16}.$$
(5)
However, we shall use the current bounds on $`\theta _\alpha `$ which are
$$\mathrm{sin}^2\theta _e6.610^3,\mathrm{sin}^2\theta _\mu 6.010^3\text{ and }\mathrm{sin}^2\theta _\tau 1.810^2.$$
(6)
Note that the lowest value is for the muon sector. Eq. (3) can now be applied to calculate the width of the Majorana, which is dominated by the two–body decays $`NW\alpha `$ and $`NZ\nu _\alpha `$, we find
$$\begin{array}{c}\mathrm{\Gamma }(N)=\underset{\alpha }{}\frac{G_F\mathrm{sin}^2\theta _\alpha }{8\pi \sqrt{2}M_N^3}\{[(M_N^4M_W^4)+M_W^2(M_N^2M_W^2)](M_N^2M_W^2)\\ +\mathrm{cos}^2\theta _\alpha [(M_N^4M_Z^4)+M_Z^2(M_N^2M_Z^2)](M_N^2M_Z^2)\}.\end{array}$$
(7)
Direct searches for heavy neutrinos give typical lower limits on their mass of 70 to 100 GeV, depending on their character (Dirac in general gives a higher bound) and to which lepton family they couple to. Unfortunately, the maximal value of the cross section of process (1) in Fig. 1 is found to lie in that range as well . The dependence on the mass goes as
$$d\sigma \frac{M_N^2}{(q^2M_N^2)^2}\{\begin{array}{cc}M_N^2& \text{ for }M_N^2q^2\\ M_N^2& \text{ for }M_N^2q^2\end{array},$$
(8)
where $`q`$ is the momentum of the Majorana. The standard calculation gives for the matrix element (see Fig. 1 for the attachment of momenta):
$$\begin{array}{c}|\overline{}|^2(e^+q\overline{\nu _e}\text{ }\alpha ^+\alpha ^+q^{})=\mathrm{sin}^4\theta _\alpha M_N^2G_F^4M_W^8\mathrm{\hspace{0.17em}2}^{12}\frac{1}{\left(q_1^2M_W^2\right)^2\left(q_3^2M_W^2\right)^2}(k_1p_2)\\ [\frac{1}{\left(q_2^2M_N^2\right)^2}(k_2p_1)(k_3k_4)+\frac{1}{\left(\stackrel{~}{q_2}^2M_N^2\right)^2}(k_3p_1)(k_2k_4)\\ \frac{1}{\left(q_2^2M_N^2\right)\left(\stackrel{~}{q_2}^2M_N^2\right)}\text{(}(k_2k_3)(p_1k_4)(k_2p_1)(k_3k_4)(k_3p_1)(k_2k_4)\text{)}]\end{array}$$
(9)
and the scattering with an antiquark sees $`k_4`$ interchanged with $`p_2`$. Here $`\stackrel{~}{q_2}`$ denotes the momentum of the Majorana in the crossed diagram, which has a relative sign due to the interchange of two identical fermion lines. In addition one has to include a factor $`\frac{1}{2}`$ to avoid double counting in the phase space integration. For the phase space we called the routine GENBOD and for the parton distributions we applied GRV 98 . In case a $`\tau `$ is produced we additionally folded in its three–body decay. We inserted finite ($`W`$ and $`N`$) width effects in our program and found them to be negligible.
An interesting statistical effect occurs when one considers the relative difference between, say, the $`\mu \mu `$ and the $`\mu \tau `$ final state (mass effects play no significant role): First, there is no factor $`\frac{1}{2}`$ for the latter case. Then, there is the possibility that a $`\tau `$ is produced at the (“upper”) $`e^+\overline{\nu _e}\text{ }W`$ vertex or at the (“lower”) $`qq^{}W`$ vertex. Both diagrams are topologically distinct and thus have to be treated separately. This means, four diagrams lead to the $`\mu \tau `$ final state, whereas only two lead to the $`\mu \mu `$ final state. We see that there is a relative factor 4 between the two cases. Note though that now the interference terms are added to the two squared amplitudes since there is no relative sign between the two. This reduces the relative factor to about 3. However, effects of kinematical cuts and the limits of Eq. (6) wash out this phenomenon. A similar situation occurs when one studies the $`\tau \tau `$ case and lets the $`\tau `$’s decay into different particles (e.g. $`e\nu \nu `$ and $`\mu \nu \nu `$). There is no way to tell into what the “upper” or “lower” tau decays, so one has to include both cases.
A question arises if one can conclude a Majorana mass term if we measure a process like Eq. (1). Here, a simple generalization of arguments first given by Schechter and Valle for neutrinoless double beta decay applies: Assuming we found indubitable evidence for $`e^+q\overline{\nu _e}\text{ }\alpha ^+\beta ^+q^{}`$, then crossing permits the process $`0e^{}\overline{\nu _e}\text{ }\alpha ^+\beta ^+q^{}\overline{q}`$, realized by the “black box” in Fig. 2. Any reasonable gauge theory will have $`W`$’s couple to quarks and leptons, so that a Majorana mass term for $`\nu _\alpha `$ and $`\nu _\beta `$ is produced by coupling one $`W`$ to the $`\alpha ^+e^{}\overline{\nu _e}\text{ }`$ and one $`W`$ to the $`\beta ^+q^{}\overline{q}`$ vertex. Since we do not know which two quarks participate and which neutrino couples to the positron (the Schechter and Valle argument for 0$`\nu \beta \beta `$ works with two pairs of $`u`$ and $`d`$ quarks), this theorem holds for a greater class of models, namely e.g. those with direct $`e^{}\underset{X}{\overset{()}{\nu }}`$ coupling, with $`X`$ being any flavor.
The connection between a neutrinoless double beta decay signal and Majorana masses has been expanded in to supersymmetric theories and it was found that it implies Majorana masses also for sneutrinos, the scalar superpartners of the neutrinos.
However, in contrast to the signal in neutrinoless double beta decay experiments (two electrons with constant sum in energy) the identification of our process will be very difficult and deciding which $`\alpha \beta `$ final state was originally produced remains a hard task. In addition, the number of expected events turns out to be far less than one. Nevertheless, the demonstration of Majorana mass terms will be an exciting and important result, since different models predict different texture zeros in the mass matrix. In some models the $`ee`$ entry in the mass matrix is zero and therefore the only direct information about the mass matrix might come from neutrino oscillations, cosmological considerations, global fits and direct searches, e.g. at LEP. This complicates the situation, since e.g. in oscillations only mass squared differences are measured and the additional phases induced by the Majorana nature are unobservable. To combine all information from the different approaches input from models is required. Thus, the exact form of the matrix is highly nontrivial to find. Therefore, information about non–vanishing entries in the mass matrix is very important and one has to take every opportunity to find out about all elements and the Majorana character in general. In addition, if such a lepton–number violating process is detected, it is surely helpful to know if the “mildly extended” SM can provide the signal or another theory, such as SUSY, has to be blamed.
## 3 Two become one
We applied the same cuts as H1 in their search for isolated leptons:
* Imbalance in transverse momentum $`p/_T25`$ GeV
* Transverse momentum of lepton $`p_T10`$ GeV
* Pseudorapidity of lepton $`|\eta |2.436`$
* Distance between charged lepton and closest jet in $`\eta `$$`\varphi `$ space<sup>4</sup><sup>4</sup>4Actually H1’s value is 1.0 or 0.5, depending on the way they define jets for their respective analysis. We use a general value of 1.5 to account for hadronization effects., $`\mathrm{\Delta }R1.5`$ where $`\varphi `$ is the azimuthal angle
* Angle of the hadronic jet(s) $`4^{}\theta ^X178^{}`$
This has to be compared with our cuts in , $`p/_T10`$ GeV, $`|\eta |2.0`$ for all measured particles and $`\mathrm{\Delta }R0.5`$ between the charged leptons and the hadronic remnants. It turns out that both sets of cuts deliver cross sections in the same ballpark. There are now two possibilities for the original LSD signal to appear as one single lepton:
* One lepton can have high pseudorapidity and/or low $`p_T`$. This lepton then also contributes to the missing transverse momentum.
* If one tau is produced it might decay hadronically, adding one neutrino to the imbalance in $`p_T`$ and also additional hadronic jets.
A detailed analysis of the LSD signal might be done if one finds such events. Collecting all possibilities for the $`\alpha \beta `$ final states and the $`\tau `$ decays results in Fig. 3. We denote with “hadronic” the signal coming from final states which also have additional hadronic activity from a $`\tau `$ decay. We call “leptonic” the signals coming from events in which two final state leptons are produced from which one escapes the identification criteria. Those included therefore most channels, namely all of them except the ones with hadronic tau decay. One can see that muon events have a smaller cross section than the $`e`$ signal, coming from the fact that their mixing with the heavy neutrino has the biggest constraint. Mass effects play no significant role. If the H1 anomaly is indeed explained by heavy Majorana neutrinos, one might ask why only muon events are detected. A possible reason for that might lie in the following fact: The experimental constraint from 0$`\nu \beta \beta `$ on mixing with a heavy Majorana neutrino reads
$$\mathrm{sin}^2\theta _e510^8\frac{m_N}{\mathrm{GeV}}.$$
(10)
Incorporating this in Fig. 3 gives Fig. 4. Now the electron signal is far below the muon signal. In Fig. 5 we plot how the cross section for the production of a $`\mu `$ is composed. The biggest contribution comes from the $`\mu \tau `$ channel, which has its reason in the mentioned factor $`3`$ relative to the $`\alpha \alpha `$ channels and the high hadronic tau branching ratio, BR$`(\tau \nu _\tau \mathrm{hadrons})\frac{2}{3}`$.
An additional reason why the $`\tau \alpha `$ channels add higher contributions is that when the tau decays it distributes its momentum to three particles, i.e. the $`p_T`$ is in general lower and it can therefore escape the $`p_T10`$ GeV cut more easily and thus has a higher cross section. Addmittedly, the process gives only a tiny signal: Multiplying the cross section with the 36.5 pb<sup>-1</sup> luminosity H1 analyzed, gives $`10^8\mathrm{}10^9`$. However, as explained at the end of the previous section, one has to check every possible appearance of Majorana mass terms in order to get information about the mass matrix.
One sees that many LSD signals produce the same single lepton signature. In the next section we will discuss possibilities to distinguish different original final states from the measured ones.
## 4 Signals and observables
Some kinematic quantities of H1’s positive muon events are given in Table 1. In their analysis using $`36.5\pm 1.1`$ pb<sup>-1</sup> luminosity at $`E_p=820`$ GeV and $`E_e`$ = 27.5 GeV, 6 events were found (0 $`e^+`$, 1 $`e^{}`$, 2 $`\mu ^+`$, 2 $`\mu ^{}`$ and 1 $`\mu `$ of undetermined charge), where about 2 $`e`$ and 1 $`\mu `$ are expected from SM processes. From those, the most important ones are $`W`$ production, NC events (for $`e^+`$ events) and photon–photon interactions ($`\mu ^\pm `$). We also include the event with undetermined charge. The $`e^{}`$ and one $`\mu ^{}`$, which also has a $`e^+`$, are very likely to stem from $`W`$ production. In the following we will plot all distributions for $`M_N=200`$ GeV since the qualitative conclusions we draw remain valid for all masses considered. In our analysis it turned out that — with the kinematics from Table 1 — scatter plots of $`p_T`$, $`p/_T`$ and the transverse mass $`M_T=\sqrt{(p/_T+p_T)^2(\stackrel{}{p}/_T+\stackrel{}{p}_T)^2}`$ are most useful. We stress again that for different final states the composition of the missing transverse momentum can be made of two particles (for $`\mu \mu `$ or $`ee`$ final states) to 6 (1 lepton and 5 neutrinos, $`\tau \tau `$ channel with two leptonic decays) thus changing the area in which events populate, say, the $`M_T`$$`p_T`$ space. The transverse momentum is also very sensitive on the original final state since $`\tau `$ decays share the initial momentum to three particles thereby reducing the average $`p_T`$. This is displayed in Fig. 6 ($`\mu \mu `$ final state) and Fig. 7 ($`\tau \tau `$, one hadronic decay). If both taus decay leptonically the distribution looks similar. Obviously, the $`\tau \tau `$ case has in general low $`p_T`$ and $`M_T`$, whereas the $`\mu \mu `$ case displays an uniform distribution with a slight band in the center region indicated.
Turning now to $`p/_T`$ we see in Figs. 8 ($`\mu \mu `$) and 9 ($`\mu \tau `$, hadronic decay) that the situation is not as clear in the “mixed” channels, i.e. channels with two different final state charged leptons. Though the population in $`p/_T`$$`M_T`$ space is different (lower values in the latter case), it is not as obvious as for the $`\mu \mu `$/$`\tau \tau `$ case. Due to its low $`M_T`$, event $`\mu _1`$ seems to be less probable in all figures, though no definite statement can be made.
Now we consider the quantities connected with the hadronic remnants. We found that the muon signal is composed of roughly 1/3 purely leptonic final states and 2/3 events with additional jets. An interesting quantity is $`\delta =E_i(1\mathrm{cos}\theta _i)`$ where the sum goes over all measured final state particles. In Figs. 10 and 11 one sees that the presence of three jets keeps $`\delta `$ in more or less in the same area whereas the transverse momentum of the hadronic system $`p_T^X`$ is shifted towards lower values. Here also $`\mu _1`$ lies in a less crowded area.
What are now the signals of the escaping charged lepton (if there is one)? In Figs. 12 and 13 we plot the pseudorapidity $`\eta `$ of the undetected lepton against its transverse momentum. Again, the $`\mu \mu `$ and the $`\tau \tau `$ case can be distinguished since the latter has far lower $`p_T`$. The $`\tau `$ boosts its decay products more or less in its forward direction so that $`\eta `$ does not alter much. The problem one might encounter is that the escaping lepton is hiding in the hadronic jet. Demanding a distance of $`\mathrm{\Delta }R1.5`$ reduces the cross section by 10 to 15 $`\%`$ but does not change the distributions in Figs. 12 and 13.
We mimicked the hadronic $`\tau `$ decay via two quark jets and ignored effects of modes like $`\tau \pi `$’s $`\nu _\tau `$. Thus, due to the boost of the $`\tau `$, two of the three jets of events with a hadronically decaying tau will be very close together. Fig. 14 displays the normalized distribution of the distance in $`\eta `$$`\varphi `$ space. We denote with $`R_i`$ the distance ordered with ascending value. One distance is centered significantly below one, therefore, probably two jets instead of three will be measured. However, the $`\tau `$ identification is hard to do and Fig. 14 serves only as an indication of how things might work. Information on the jet multiplicity is not given in Ref. , though $`\mu _1`$ and $`\mu _5`$ seem to have additional separated tracks in their event displays as can be seen in Fig. 2 of Ref. .
## 5 Conclusions
In the light of recent data we did a full analysis of the analogue of the neutrinoless double beta decay graph with all possible two charged leptons in the final state. One lepton can escape the identification criteria by either having low transverse momentum and/or high pseudorapidity or (if it is a $`\tau `$) via hadronic decay. Signatures of these events are discussed and compared to H1’s isolated leptons with large missing transverse momentum. All their positive muon events lie in the regions typically populated by the “double beta” process, though $`\mu _1`$ is always in a less crowded area. Due to its high errors $`\mu _5`$ is mostly in the favored region (but for the same reason always in the region populated by SM processes ). Though the process represents an attractive explanation, the smallness of the expected signal might spoil our interpretation. Nevertheless, any information about mass matrix entries and Majorana particles is very important and worth looking for, regardless of the small expected signal. Other extensions of the SM might give larger signals to the discussed final states and it is then helpful to know how the “SM + heavy Majorana” extension contributes.
To confirm our hypothesis direct production of heavy Majorana neutrinos will be the only possibility since cross sections or decay widths of other 0$`\nu \beta \beta `$–like processes are probably too small to be detected . Here, either HERA itself or LHC will be the candidates for this observation. We did not consider the mass reconstruction of the Majorana since the large number of unmeasured particles does not permit that.
Other proposed explanations for the events were FCNC interactions (topologically identical to leptoquark production) or high $`p_T`$ jets from which one fakes a muon signal. Production of supersymmetric particles is suggested in to explain the results. A definite answer regarding all detector/identification issues can only be given by the collaboration itself. From the “new physics” side we believe that massive neutrinos provide one of the most natural possibilities.
Acknowledgments
Its a pleasure to thank M. Flanz for helpful discussions and careful reading of the manuscript. This work has been supported in part (W. R.) by the “Bundesministerium für Bildung, Wissenschaft, Forschung und Technologie”, Bonn under contract No. 05HT9PEA5. Financial support from the Graduate College “Erzeugung und Zerf$`\ddot{\mathrm{a}}`$lle von Elementarteilchen” at Dortmund university (W. R.) is gratefully acknowledged.
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# FZJ-IKP(TH)-2000-12 𝐽/Ψ→ϕ𝜋𝜋(𝐾𝐾̄) decays, chiral dynamics and OZI violation
## I Introduction
The decays of the $`J/\mathrm{\Psi }`$ into a $`\varphi `$ meson and Goldstone boson pair ($`\pi \pi `$ or $`K\overline{K}`$) can be used to investigate the dynamics of the interacting pseudoscalars. In particular, it was argued in Ref. that these data together with data from pion–pion scattering (and from others sources) force the $`f_0(980)`$ to have a pole structure different to the one required by a $`K\overline{K}`$ molecule . This interpretation has been challenged, e.g. in the Jülich meson–exchange model where the $`f_0(980)`$ emerges as a $`K\overline{K}`$ bound state. Furthermore, in this reference the authors are also able to reproduce the data associated with the previous $`J/\mathrm{\Psi }`$ decays within the same formalism than the one employed in Ref. , but making use of their own strong amplitudes. On the other hand, as will be the topic of this investigation, these data can be used to study the violation of the Okubo–Zweig–Iizuka (OZI) rule in the scalar ($`0^{++}`$) channel. This rule is only well founded in the large $`N_c`$ limit of QCD, with $`N_c`$ the number of colors, since OZI violating processes are described by suppressed non–planar graphs . Still, on a purely phenomenological level this rule works astonishingly well, with the exception of the scalar channel, as argued e.g. in Refs.,,. To be more precise, the decay $`J/\mathrm{\Psi }\varphi M\overline{M}`$ (where $`M\overline{M}`$ denotes the pseudoscalar meson pair) is OZI suppressed to leading order, cf. Fig.1a, but has an additional doubly OZI suppressed contribution depicted in Fig.1b. In our approach, both these pieces are taken into account. In fact, it will turn out that the second contribution can not be neglected if one wants to achieve an accurate description of the data. On the other hand, it is mandatory to have a very precise description of the final state interaction in the coupled $`\pi \pi /K\overline{K}`$ system (as indicated by the shaded blob in Fig.1a) before one can ask such detailed questions. As can be seen from Fig.1, the crucial ingredient in the reaction at hand are the expectation values of the scalar–isoscalar condensates in the pion and the kaon, i.e the so-called scalar form factors. These can be calculated at low energies in chiral perturbation theory (CHPT), which is the effective field theory of the Standard Model. In our case, the dimeson system can have energies up to 2 GeV and we thus employ unitarity constraints to get a precise description of these scalar form factors also at higher energies, demanding furthermore matching to the CHPT expressions in the low energy domain. Because of this matching procedure, the large $`N_c`$ suppressed low energy constants $`L_4^r`$ and $`L_6^r`$ of the next–to–leading order effective chiral Lagrangian can be determined in the process we are considering. It has been argued before that so far no direct determinations but rather large $`N_c`$ inspired estimates have been done, see e.g. Refs.,, with the exception of more recent work presented in Refs.. Nevertheless, as we will discuss in much more detail below, a rather definite determination of $`L_4^r`$ can be obtained by considering $`𝒪(p^6)`$ CHPT results .
To be more specific, to address the problem of the final state interactions in the coupled $`\pi \pi `$-$`K\overline{K}`$ system, we make use of the results obtained in Ref.. In this paper it was clearly established that the scattering data of the $`0^{++}`$ $`I=0,\mathrm{\hspace{0.17em}1}`$ ($`I`$ denotes the total isospin of the dimeson system) sectors up to centre–of-mass energies of 1.2 GeV are a reflection of the strong rescattering effects between the lightest pseudoscalars ($`\pi \pi `$, $`K\overline{K}`$ for $`I=0`$ and $`\pi \eta `$ and $`K\overline{K}`$ for $`I=1`$). The approach was based on Bethe-Salpeter equations using the lowest order CHPT amplitudes as the driving potential. The fact that one can generate the resonance states for those channels via loop physics, i.e. rescattering, is a clear signal of the large deviations from OZI rule in the $`0^{++}`$ sector, see also Refs.. Such a mechanism has been advocated since long, for a pedagogic discussion see Ref.. On the other hand, it is well known that there is an on–going controversy concerning the nature of the scalar resonances $`f_0(980)`$ and $`a_0(980)`$. This controversy originates from the observation that there are several different models to deal with the $`I=0,\mathrm{\hspace{0.17em}1}`$ scalar sector, all of them reproducing the scattering data up to some extend, but with different conclusions with respect to the origin of the underlying dynamics. In particular, in Refs. these resonances are considered of preexisting origin while in Ref. they appear as meson–meson resonances originated by a potential. Also in Ref. it is advocated for the solution that the $`a_0(980)`$, $`f_0(980)`$ are exotic resonances, that is, not simply $`q\overline{q}`$, while the preexisting $`q\overline{q}`$ scalar nonet should be heavier, around 1.4 GeV or so. Other interesting approaches to this problem are the relativistic quark model with an instanton induced interaction of the Bonn group , the Jülich meson–exchange approach or the use of QCD sum rules . With respect to this controversy, the contribution of the work in Ref. is very valuable since, at least, the infinite series of diagrams there considered should appear in the whole S–wave partial wave amplitudes calculated to all orders in CHPT. The conclusions of Ref. where generalized in Ref.. In that paper, the most general structure of a partial wave amplitude when the unphysical cuts are neglected was established. In particular, in this paper explicit s-channel resonance exchanges were included together with the lowest order CHPT contribution and the whole $`SU(3)`$ connected scalar sector with $`I=0,1/2,1`$ was studied. In particular, it was shown that the amplitudes of Ref. appear as a particular case when removing the explicit tree level resonance contributions. It was observed that the lightest $`0^{++}`$ nonet is of dynamical origin, i.e. made up of meson–meson resonances, and is formed by the $`\sigma (500)`$, $`\kappa `$, $`a_0(980)`$ and a strong contribution to the physical $`f_0(980)`$. On the other hand, the preexisting scalar nonet would be made up by an octet around 1.4 GeV and a singlet contributing to the physical $`f_0(980)`$ resonance. With respect to this last point, as discussed in Ref., the inclusion of a preexisting contribution to the $`f_0(980)`$ was considered in order to be able to reproduce the data on the inelastic $`\pi \pi K\overline{K}`$ cross section<sup>#3</sup><sup>#3</sup>#3In the last edition of the PDG tables it is argued that, possibly, the previous experiments have a much larger uncertainty than previously given in the corresponding publications. when including also the $`\eta \eta `$ channel. However, if this channel is not considered, one can reproduce the strong scattering data, including also the previous experiments on the inelastic $`\pi \pi K\overline{K}`$ cross section, without including such preexisting contribution. Finally, in Ref. the contribution in the physical region of the unphysical cut contributions were estimated up to $`\sqrt{s}800`$ MeV to be just a few per cents making use of the results of Ref., which apply below that energy. We will use the formalism of Ref. whose partial wave amplitudes have been also tested in many other reactions. As pointed out in Ref., to obtain a consistent picture of the scalar sector, one also has to study other reactions in which the $`0^{++}`$ amplitudes have a possible large influence via final state interactions. In this way one can complement the deficient information coming from the direct strong S–wave scattering data and distinguish between available models. In Ref. all the whole set of photon fusion reactions $`\gamma \gamma \pi ^0\pi ^0`$, $`\pi ^+\pi ^{}`$, $`K^+K^{}`$, $`K^0\overline{K}^0`$ and $`\pi ^0\eta `$ were reproduced in an unified way from threshold up to $`\sqrt{s}1.4`$ GeV making use of Ref. to take care of the final state interactions. The free parameters present in Ref. were fixed by their values from the PDG . In Ref., making use of the formalism set up in Ref. to study the still unmeasured $`\varphi \gamma K^0\overline{K}^0`$ decay, the reactions $`\varphi \gamma \pi ^0\pi ^0`$, $`\gamma \pi ^+\pi ^{}`$ and $`\gamma \pi ^0\eta `$ were predicted. These predictions were nicely confirmed, almost simultaneously, by a recent experiment in Novosibirsk . In this manuscript, we will consider yet another applications of the strong amplitudes calculated in Ref. by studying the $`J/\mathrm{\Psi }\pi \pi (K\overline{K})`$ decays. In this way, the present study together with the whole set of works offer a unique theoretical approach to the scalar sector able to discuss all these reactions in an unified way. This is achieved without including new elements ad hoc for each reaction, because all these processes are related by the use of an effective theory description that combines CHPT and unitarity constraints.
Our manuscript is organized as follows. In section II we develop and justify a simple phenomenological model for the transition of the $`J/\mathrm{\Psi }`$ into the $`\varphi `$ meson and a pseudoscalar meson pair. This model is generalized in section IV by making use of flavor SU(3) symmetry and then applied to the $`\omega `$ case. Section III is devoted to the calculation of the scalar pion and kaon form factors for the non–strange as well as the strange scalar density to one loop. Then we employ the unitarization procedure discussed before to obtain a description of these quantities up to energies of 1.2 GeV. We also perform the matching of these expressions to the one loop CHPT ones to be consistent with the constraints of the broken chiral symmetry of QCD at energies below approximately 0.5 GeV. The results are presented and discussed in section IV where we also generalize section II. Our conclusions and outlook are given in section V.
## II Modeling $`J/\mathrm{\Psi }\varphi \pi \pi ,\varphi K\overline{K}`$ decays
We will calculate the S-wave contribution to the invariant mass distributions of the $`\pi \pi `$ and $`K\overline{K}`$ systems in the $`J/\mathrm{\Psi }\varphi \pi \pi (K\overline{K})`$ decays. Taking care of the final state interactions of three particles can be simplified to a large extend if only two of the final state particles undergo strong interactions, and the third is merely a spectator. We will assume that this is the case in our present problem and we will take the $`\varphi `$ as the spectator. This is certainly very well sounded for the $`\varphi \pi \pi `$ system since the $`\varphi \pi `$ interaction is very weak as required by the OZI rule. On the other hand, the situation is not so clear with respect the $`\varphi `$ and the kaons. Nevertheless, at the energies in which the kaons become important, above the $`K\overline{K}`$ threshold, the experimental mass distribution is completely dominated by the $`f_0(980)`$ resonance and this state is a two body effect emerging from the coupled $`\pi \pi `$ and $`K\overline{K}`$ systems, as discussed already in the introduction. Since we are only considering a small range of energies above the $`K\overline{K}`$ threshold, this approximation should be justified.
We therefore describe the transition from the $`J/\mathrm{\Psi }`$ to the $`\varphi `$+2 Goldstone bosons system by an effective Lagrangian based on the following phenomenological arguments: 1) The already discussed spectator role of the $`\varphi `$ resonance and 2) the $`\pi \pi `$ and $`K\overline{K}`$ invariant event distributions, which will be shown later, seem to be clearly dominated by the S-wave contribution, although these experimental data have not yet been subjected to a partial-wave analysis. These experimental facts, together with Lorentz invariance, can be easily incorporated in the formalism just by writing the interaction vertex of the $`J/\mathrm{\Psi }`$ resonance with the $`\varphi `$ meson and some scalar source $`S`$ with vacuum quantum numbers, $`J^{PC}=0^{++}`$ with spin $`J`$, parity $`P`$ and charge conjugation $`C`$, as:
$$g\mathrm{\Psi }_\mu \varphi ^\mu S$$
(1)
with $`g`$ a real coupling constant. We briefly discuss why other possible structures involving derivatives on the various fields should be suppressed. The $`J/\mathrm{\Psi }`$ is very heavy and thus can be considered a static source. Derivatives acting on the $`\varphi `$ and the scalar source $`S`$ can be combined to the invariant structure $`\mathrm{\Psi }^\mu (_\nu \varphi _\mu _\mu \varphi _\nu )^\nu S`$. This leads to a vertex of the form $`ϵ_\mathrm{\Psi }^\mu ϵ_\mu ^\varphi p_\varphi (q_1+q_2)`$, with $`q_{1,2}`$ the momenta of the two Goldstone bosons and $`p_\varphi `$ the momentum of the $`\varphi `$ meson. However, due to momentum conservation, we have $`p_\varphi (q_1+q_2)=(M_\mathrm{\Psi }^2M_\varphi ^2s)/2`$, with $`s=(q_1+q_2)^2`$ the total two Goldstone boson energy squared. Due to the large value of the $`J/\mathrm{\Psi }`$ mass, this combination of momenta is essentially constant for the Goldstone boson energies considered here, $`\sqrt{s}1.2`$GeV. Such terms become more important at higher di–pion (kaon) energies. Therefore, we can generically write such type of higher order corrections to Eq.(1) in the form
$$ϵ_\mathrm{\Psi }^\mu ϵ_\mu ^\varphi f((q_1+q_2)^2,p_\varphi (q_1+q_2)),$$
(2)
where the function $`f(\mathrm{})`$ essentially only depends on the first argument. Such terms that depend on $`(q_1+q_2)^2`$ can be derived from an interaction of the type $`\mathrm{\Psi }^\mu \varphi _\mu _\nu ^\nu S`$. Such structures lead to a weak $`s`$–dependence of the constant $`g`$ and/or the parameter $`\lambda _\varphi `$ defined below. We have checked that such (weak) energy dependencies do not change any of the conclusions obtained when treating $`g`$ and $`\lambda _\varphi `$ as energy independent. Another possible higher order term of the form $`\mathrm{\Psi }_\mu \varphi _\nu ^\mu ^\nu S`$ giving rise to the vertex $`ϵ_\mathrm{\Psi }^\mu ϵ_\nu ^\varphi (q_1+q_2)_\mu (q_1+q_2)^\nu `$. Such couplings can also have an S–wave contribution, which can be obtained by properly summing over the pertinent polarization vectors. Again, due to the large mass of the $`J/\mathrm{\Psi }`$, such terms are only weakly $`s`$ dependent and can be treated along the lines outlined before. More complicated structures can always be brought into some linear combination of the ones just discussed or have no S–wave component. These considerations not only show that our ansatz Eq.(1) is quite sensible in the energy range considered here but also that corrections to it can be worked out consistently.
From the OZI rule, which can also be seen as a result of the large $`N_c`$ expansion of QCD , and the experimental absence of any clue indicating a non negligible interaction between the $`\varphi `$ and the pions, one should expect that this scalar source $`S`$ would be simply made of strange quarks, i.e. $`S\overline{s}s`$. However, it is known that the $`\varphi `$ also decays into non–strange mesons and furthermore, there are strong arguments to believe that large violations of the OZI rule (and of the large $`N_c`$ limit of QCD) are manifest in the $`0^{++}`$ sector , as discussed in the introduction. As a result, we will consider a more general scalar source $`S`$, that also has a contribution of the form $`\lambda _\varphi \overline{n}n`$, where
$$\overline{n}n=\frac{1}{\sqrt{2}}(\overline{u}u+\overline{d}d)$$
(3)
parameterizes the contribution from the non–strange quarks and $`\lambda _\varphi `$ is just a constant measuring the relative strength of this contribution with respect to the strangeness component $`\overline{s}s`$. Already at this point we stress that the choice $`\lambda _\varphi 0`$ will be justified a posteriori by the results presented below. Therefore, we use
$$S=\overline{s}s+\lambda _\varphi \overline{n}n$$
(4)
Consequently, it follows from Eqs.(1) and (4) that the transition matrix element for the process $`J/\mathrm{\Psi }\varphi M\overline{M}`$ is given by:
$$T=ϵ(\mathrm{\Psi };\rho )_\mu ϵ(\varphi ;\rho ^{})^\mu 0|\left(\overline{s}s+\lambda _\varphi \overline{n}n\right)|M\overline{M}^{}$$
(5)
where $`|0`$ is the vacuum state, $`ϵ(\mathrm{\Psi };\rho )`$ is the polarization four–vector of the $`J/\mathrm{\Psi }`$ resonance with polarization $`\rho `$ and analogously $`ϵ(\mathrm{\Psi };\rho ^{})`$ is the polarization four–vector of the $`\varphi `$ resonance, and the denotes complex conjugation. Note that in this equation we are implicitly assuming that the $`\varphi `$ is a spectator as discussed before. The scalar source $`S`$ couples to the two meson system, in which the rescattering (final state interactions) appear. The anatomy of our model is depicted in Fig.2. As discussed below, invoking SU(3) symmetry, we will also apply this approach to the S–wave contribution of the $`J/\mathrm{\Psi }\omega \pi \pi `$ decay to further constrain the description of the measured event distributions.
## III Coupled channel pion and kaon scalar form factors
As a consequence of Eq.(5), our problem is reduced to calculate the matrix elements $`0|\overline{s}s|M\overline{M}`$ and $`0|\overline{n}n|M\overline{M}`$, which correspond to the strange and non-strange isospin zero ($`I=0`$) scalar form factors.
### A Definition of the scalar form factors
One can define an extended QCD Lagrangian allowing for the presence of external sources. In this way the identification of matrix elements of quark currents can be done easily. For instance, a scalar source can be added simply as:
$$\overline{q}\mathrm{\Sigma }q,$$
(6)
where $`q`$ embodies the three light quarks, $`u`$, $`d`$ and $`s`$. The QCD current quark mass term can be obtained from such a scalar source by setting,
$$\mathrm{\Sigma }=\mathrm{diag}(m_u,m_d,m_s).$$
(7)
This is the standard method of treating explicit chiral symmetry breaking in CHPT (or any similar effective field theory). Consequently, we can work out the scalar quark–antiquark operators,
$`\overline{u}u`$ $`=`$ $`{\displaystyle \frac{_{\mathrm{QCD}}}{\mathrm{\Sigma }_{11}}}={\displaystyle \frac{_{\mathrm{QCD}}}{m_u}},`$ (8)
$`\overline{d}d`$ $`=`$ $`{\displaystyle \frac{_{\mathrm{QCD}}}{\mathrm{\Sigma }_{22}}}={\displaystyle \frac{_{\mathrm{QCD}}}{m_d}},`$ (9)
$`\overline{s}s`$ $`=`$ $`{\displaystyle \frac{_{\mathrm{QCD}}}{\mathrm{\Sigma }_{33}}}={\displaystyle \frac{_{\mathrm{QCD}}}{m_s}}.`$ (10)
One can include in the same way as in the QCD Lagrangian the external sources in the effective CHPT Lagrangian simply based on symmetry arguments. In the lowest order chiral effective Lagrangian, $`_2`$, the scalar source appears in the mass term
$$_2^{\mathrm{mass}}=\frac{1}{4}f^2U^{}\chi +\chi ^{}U,$$
(11)
with $`f`$ the meson decay constant (in the chiral limit), $`\chi 2B_0\mathrm{\Sigma }`$ and $`B_0`$ is a constant not fixed by symmetry. This constant parameterizes the strength of the quark–antiquark condensation in the non–perturbative vacuum, $`B_0=|0|\overline{q}q|0|/f^2`$. The trace in flavor space is denoted by $`\mathrm{}`$. The octet of Goldstone bosons is collected in the matrix–valued unimodular field $`U(x)`$,
$$U=\mathrm{exp}\left(\frac{i\sqrt{2}}{f}\mathrm{\Phi }\right)$$
(12)
with
$$\mathrm{\Phi }=\left[\begin{array}{ccc}\frac{1}{\sqrt{2}}\pi ^0+\frac{1}{\sqrt{6}}\eta _8& \pi ^+& K^+\\ \pi ^{}& \frac{1}{\sqrt{2}}\pi ^0+\frac{1}{\sqrt{6}}\eta _8& K^0\\ K^{}& \overline{K}^0& \frac{2}{\sqrt{6}}\eta _8\end{array}\right].$$
(13)
It is then straightforward to work out the scalar–isoscalar quark-antiquark operators from the effective Lagrangian,
$`\overline{u}u`$ $`=`$ $`{\displaystyle \frac{_2}{\mathrm{\Sigma }_{11}}}=f^2B_0\left[1{\displaystyle \frac{1}{f^2}}\left(\pi ^+\pi ^{}+K^+K^{}+{\displaystyle \frac{(\pi ^0)^2}{2}}+{\displaystyle \frac{\eta _8^2}{6}}+{\displaystyle \frac{\pi ^0\eta _8}{\sqrt{3}}}\right)+\mathrm{}\right]`$ (14)
$`\overline{d}d`$ $`=`$ $`{\displaystyle \frac{_2}{\mathrm{\Sigma }_{22}}}=f^2B_0\left[1{\displaystyle \frac{1}{f^2}}\left(\pi ^+\pi ^{}+K^0\overline{K}^0+{\displaystyle \frac{(\pi ^0)^2}{2}}+{\displaystyle \frac{\eta _8^2}{6}}{\displaystyle \frac{\pi ^0\eta _8}{\sqrt{3}}}\right)+\mathrm{}\right]`$ (15)
$`\overline{s}s`$ $`=`$ $`{\displaystyle \frac{_2}{\mathrm{\Sigma }_{33}}}=f^2B_0\left[1{\displaystyle \frac{1}{f^2}}\left(K^+K^{}+K^0\overline{K}^0+{\displaystyle \frac{2}{3}}\eta _8^2\right)+\mathrm{}\right]`$ (16)
where the ellipsis denotes terms of higher order in the meson fields not needed here. From the last of these equations one concludes that the strangeness component of the pion should be small since it only comes in at higher orders. From this representation of the scalar operators, one can deduce the pertinent expressions for the scalar form factors. For our purposes, it is sufficient to consider pure isospin zero ($`I=0`$) states formed from a pion or kaon–anti-kaon pair, i.e.
$`|\pi \pi `$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{6}}}|\pi ^+\pi ^{}+\pi ^{}\pi ^++\pi ^0\pi ^0,`$ (17)
$`|K\overline{K}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}|K^+K^{}+K^0\overline{K}^0.`$ (18)
Note the extra factor $`1/\sqrt{2}`$ in the definition of the $`I=0`$ $`|\pi \pi `$ state, it is introduced to take care that in the isospin basis states the pions behaves as identical particles. Combining this with Eq.(14), one can easily calculate the lowest order (tree level) CHPT results (remember the normalization of the non–strange quark operator given in Eq.(3)),
$`0|\overline{n}n|\pi \pi `$ $`=`$ $`\sqrt{3}B_0,`$ (19)
$`0|\overline{n}n|K\overline{K}`$ $`=`$ $`B_0,`$ (20)
$`0|\overline{s}s|\pi \pi `$ $`=`$ $`0,`$ (21)
$`0|\overline{s}s|K\overline{K}`$ $`=`$ $`\sqrt{2}B_0.`$ (22)
As anticipated, to leading order the two–pion system has no strangeness component. To all orders, these matrix elements are given in terms of four scalar form factors,<sup>#4</sup><sup>#4</sup>#4We remark that more commonly the definition of these form factors includes the pertinent quark masses, such that e.g. the non-strange scalar form factor of the pion is defined via $`0|\widehat{m}(\overline{u}u+\overline{d}d)|\pi \pi =M_\pi ^2\mathrm{\Gamma }_\pi (s)`$. For our later discussion, the overall normalization does not play a role but should be kept in mind.
$`0|\overline{n}n|\pi \pi `$ $`=`$ $`\sqrt{2}B_0\mathrm{\Gamma }_1^n(s),`$ (23)
$`0|\overline{n}n|K\overline{K}`$ $`=`$ $`\sqrt{2}B_0\mathrm{\Gamma }_2^n(s),`$ (24)
$`0|\overline{s}s|\pi \pi `$ $`=`$ $`\sqrt{2}B_0\mathrm{\Gamma }_1^s(s),`$ (25)
$`0|\overline{s}s|K\overline{K}`$ $`=`$ $`\sqrt{2}B_0\mathrm{\Gamma }_2^s(s).`$ (26)
Here, the following notation is employed. The superscript $`s/n`$ refers to the strange/non–strange quark operator whereas the subscript $`1,2`$ denotes pions and kaons, respectively. In the following we will remove from Eqs.(19,23) the overall factor $`\sqrt{2}B_0`$, since the experimental data on the $`J/\mathrm{\Psi }\varphi M\overline{M}`$ decays are not normalized.
### B Next-to-leading order pion and kaon scalar form factors
The pion scalar form factors $`\mathrm{\Gamma }_1^n(s)`$ and $`\mathrm{\Gamma }_1^s(s)`$ were calculated in Ref. up to one loop in CHPT. Since they were not explicitly given in Ref., we give here the pertinent expressions:
$`\mathrm{\Gamma }_1^n(s)`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}\left\{1+\mu _\pi {\displaystyle \frac{1}{3}}\mu _\eta +{\displaystyle \frac{16m_\pi ^2}{f^2}}\left(2L_8^rL_5^r\right)+8(2L_6^rL_4^r){\displaystyle \frac{2m_K^2+3m_\pi ^2}{f^2}}+f(s)+{\displaystyle \frac{2}{3}}\stackrel{~}{f}(s)\right\},`$ (27)
$`\mathrm{\Gamma }_1^s(s)`$ $`=`$ $`(2L_6^rL_4^r){\displaystyle \frac{8\sqrt{3}m_\pi ^2}{f^2}}+{\displaystyle \frac{1}{\sqrt{3}}}\stackrel{~}{f}(s),`$ (28)
with $`f(s)`$ and $`\stackrel{~}{f}(s)`$ given by
$`f(s)`$ $`=`$ $`{\displaystyle \frac{2sm_\pi ^2}{2f^2}}\overline{J}_{\pi \pi }(s){\displaystyle \frac{s}{4f^2}}\overline{J}_{KK}(s){\displaystyle \frac{m_\pi ^2}{6f^2}}\overline{J}_{\eta \eta }(s)+{\displaystyle \frac{4s}{f^2}}\left\{L_5^r{\displaystyle \frac{1}{256\pi ^2}}\left(4\mathrm{log}{\displaystyle \frac{m_\pi ^2}{\mu ^2}}\mathrm{log}{\displaystyle \frac{m_K^2}{\mu ^2}}+3\right)\right\},`$ (29)
$`\stackrel{~}{f}(s)`$ $`=`$ $`{\displaystyle \frac{3}{4}}{\displaystyle \frac{s}{f^2}}\overline{J}_{KK}(s)+{\displaystyle \frac{m_\pi ^2}{3f^2}}\overline{J}_{\eta \eta }(s)+{\displaystyle \frac{12s}{f^2}}\left\{L_4^r{\displaystyle \frac{1}{256\pi ^2}}\left(\mathrm{log}{\displaystyle \frac{m_K^2}{\mu ^2}}+1\right)\right\},`$ (30)
and $`\overline{J}_{PP}(s)`$ ($`P=\pi ,K,\eta `$) is the standard meson loop function
$$\overline{J}_{PP}(s)=\frac{1}{16\pi ^2}\left(2+\sigma _P(s)\mathrm{log}\frac{\sigma _P(s)1}{\sigma _P(s)+1}\right),$$
(31)
and $`\mu `$ is the scale of dimensional regularization. The quantities $`\mu _P`$ in Eq.(27) are given by
$$\mu _P=\frac{m_P^2}{32\pi ^2f^2}\mathrm{log}\frac{m_P^2}{\mu ^2}.$$
(32)
The scalar kaon form factors at next-to-leading order in CHPT are not given explicitly in the literature. We fill here this gap by performing such a calculation. This implies calculating the diagrams shown in Fig.3, which comprise the lowest order CHPT result, Fig.3a, already derived in section III A, the tadpole contribution, Fig.3b, and the unitarity corrections, Fig.3c. The vertices for these diagrams come from the lowest order CHPT Lagrangian. We note that wave function renormalization diagrams are not depicted in this figure. Finally, in Fig.3d, the local contribution coming from the $`𝒪(p^4)`$ CHPT Lagrangian is depicted. These terms are parameterized in terms of the scale–dependent, renormalized low energy constants $`L_i^r(\mu )`$ (in our case $`i=4,5,6,8`$). Evaluating these diagrams leads to
$`\mathrm{\Gamma }_2^n(s)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\{1+{\displaystyle \frac{4L_5^r}{f^2}}(s4m_K^2)+{\displaystyle \frac{8L_4^r}{f^2}}(2s6m_K^2m_\pi ^2)+L_8^r{\displaystyle \frac{32m_K^2}{f^2}}+{\displaystyle \frac{16L_6^r}{f^2}}(6m_K^2+m_\pi ^2)+{\displaystyle \frac{2}{3}}\mu _\eta `$ (33)
$`+`$ $`{\displaystyle \frac{9s8m_K^2}{36f^2}}J_{\eta \eta }^r(s)+{\displaystyle \frac{3s}{4f^2}}[J_{\pi \pi }^r(s)+J_{KK}^r(s)]\},`$ (34)
$`\mathrm{\Gamma }_2^s(s)`$ $`=`$ $`1+{\displaystyle \frac{4L_5^r}{f^2}}(s4m_K^2)+{\displaystyle \frac{8L_4^r}{f^2}}(s4m_K^2m_\pi ^2)+L_8^r{\displaystyle \frac{32m_K^2}{f^2}}+{\displaystyle \frac{16L_6^r}{f^2}}(4m_K^2+m_\pi ^2)+{\displaystyle \frac{2}{3}}\mu _\eta `$ (35)
$`+`$ $`{\displaystyle \frac{9s8m_K^2}{18f^2}}J_{\eta \eta }^r(s)+{\displaystyle \frac{3s}{4f^2}}J_{KK}^r(s).`$ (36)
As a test of our calculations we have checked that the infinities, associated with the wave function renormalization contributions and the loops in Fig.3c and 3d are properly absorbed by the infinite parts of the pertinent low energy constants and thus the expressions given in Eqs.(33) are finite.
### C Unitarity requirements
We now discuss the constraints that unitarity imposes on the scalar form factors. Of course, at low energies, one can simply work with CHPT and treat unitarity in a perturbative fashion. Here, however, we are interested also at energies of the order of 1 GeV, which requires some resummation technique and also the channel coupling between the $`\pi \pi `$ and the $`K\overline{K}`$ systems has to be taken into account. This has been elaborated in big detail in Ref. and we present here the formalism necessary to discuss the scalar form factors. For convenience, we employ the matrix notation already introduced in the previous subsection, i.e. pions are labeled by the index 1 and kaons by the index 2. ¿From Ref., we have the following expression for the $`T`$-matrix for meson–meson scattering,
$$T(s)=\left[I+K(s)g(s)\right]^1K(s),$$
(37)
where $`s`$ denotes the centre-of-mass energy squared and $`K(s)`$ can be obtained from the lowest order CHPT Lagrangian,
$`K(s)_{11}`$ $`=`$ $`{\displaystyle \frac{sm_\pi ^2/2}{f_\pi ^2}},`$ (38)
$`K(s)_{12}`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}s}{4f_\pi ^2}},`$ (39)
$`K(s)_{22}`$ $`=`$ $`{\displaystyle \frac{3s}{4f_\pi ^2}}.`$ (40)
We remark that because of time reversal, both $`K(s)`$ and $`T(s)`$ are symmetric functions, so that $`K(s)_{21}=K(s)_{12}`$ and similarly for $`T(s)`$. The matrix $`g(s)`$ is also diagonal and given by
$$g(s)_i=\frac{1}{16\pi ^2}\left\{\sigma _i(s)\mathrm{log}\frac{\sigma _i(s)\sqrt{1+\frac{m_i^2}{q_{\mathrm{max}}^2}}+1}{\sigma _i(s)\sqrt{1+\frac{m_i^2}{q_{\mathrm{max}}^2}}1}2\mathrm{log}\left[\frac{q_{\mathrm{max}}}{m_i}\left(1+\sqrt{1+\frac{m_i^2}{q_{\mathrm{max}}^2}}\right)\right]\right\},i=1,2,$$
(41)
where $`\sigma _i(s)=\sqrt{14m_i^2/s}`$, $`f_\pi 93`$ MeV is the weak pion decay constant, $`m_i`$ are the masses of the pions ($`m_1=138`$ MeV) and kaons ($`m_2=495.7`$ MeV) and $`q_{\mathrm{max}}=0.9`$ GeV is a cut-off in three-momentum space. On the other hand, $`g(s)_i`$ can also be calculated in dimensional regularization, using the standard $`\overline{MS}1`$ scheme employed in CHPT ,
$$g(s)_i=\frac{1}{(4\pi )^2}\left(1+\mathrm{log}\frac{m_i^2}{\mu ^2}+\sigma _i(s)\mathrm{log}\frac{\sigma _i(s)+1}{\sigma _i(s)1}\right)=J_{ii}^r(s),$$
(42)
where the last equality follows from the definition of the renormalized two–meson loop function
$$\overline{J}_{ii}(s)\frac{1}{16\pi ^2}\left(1+\mathrm{log}\frac{m_i^2}{\mu ^2}\right)=J_{ii}^r(s).$$
(43)
Note that we have changed the subscript “$`PP`$” appearing in Eq.(31) into “$`ii`$” to conform with our matrix notation. By expanding Eq.(41) in terms of $`m_i/q_{\mathrm{max}}`$ one can easily see, as discussed in appendix 2 of Ref., that the differences between $`g(s)_i`$ in Eq.(41) and Eq.(42) are of higher order in the chiral expansion, i.e. of order $`𝒪(m_i^2/q_{\mathrm{max}}^2)`$, for the following value of the scale $`\mu `$,
$$\mu =\frac{2q_{\mathrm{max}}}{\sqrt{e}}1.2q_{\mathrm{max}}.$$
(44)
For energies above the threshold of the state $`i`$, unitarity implies the following relation between form factors and the $`I=0`$ $`T`$-matrix:
$$\text{Im}\mathrm{\Gamma }_i(s)=\underset{j}{}\mathrm{\Gamma }_j(s)\frac{p_j(s)}{8\pi \sqrt{s}}\theta (s4m_j^2)(T_{ji}^{\text{S-wave}}(s))^{}$$
(45)
with $`p_i(s)=\sqrt{s/4m_i^2}`$ the modulus of the c.m. three-momentum of the state $`i`$, and the strong amplitudes are projected on the S-wave. In the former equation we have suppressed the superscript “$`n`$” or “$`s`$”, appearing in Eq.(23), since the previous equation applies to any of them. Finally, in what follows, we will also remove the superscript “S-wave” with the understanding that any partial wave is projected onto the S-wave. Taking now the complex conjugate on the right-hand-side of Eq.(45) and using the fact that the $`T`$-matrix is symmetric, we can rewrite Eq.(45) in matrix notation as:
$$\text{Im}\mathrm{\Gamma }(s)=T(s)\frac{Q(s)}{8\pi \sqrt{s}}\mathrm{\Gamma }^{}(s)$$
(46)
where
$$Q(s)=\left(\begin{array}{cc}p_1(s)\theta (s4m_1^2)& 0\\ 0& p_2(s)\theta (s4m_2^2)\end{array}\right),\mathrm{\Gamma }(s)=\left(\begin{array}{c}\mathrm{\Gamma }_1(s)\\ \mathrm{\Gamma }_2(s)\end{array}\right).$$
(47)
Substituting in the previous equation $`\text{Im}\mathrm{\Gamma }(s)`$ by $`(\mathrm{\Gamma }(s)\mathrm{\Gamma }(s)^{})/(2i)`$ and $`T(s)`$ by its expression given in Eq.(37), one has:
$`\mathrm{\Gamma }(s)`$ $`=`$ $`\left\{I+\left[I+K(s)g(s)\right]^1K(s)i{\displaystyle \frac{Q(s)}{4\pi \sqrt{s}}}\right\}\mathrm{\Gamma }(s)^{}`$ (48)
$`=`$ $`\left[I+K(s)g(s)\right]^1\left\{I+K(s)g(s)+K(s)i{\displaystyle \frac{Q(s)}{4\pi \sqrt{s}}}\right\}\mathrm{\Gamma }(s)^{}.`$ (49)
Taking into account that the $`K(s)`$-matrix, Eq. (38), is real and that
$$g(s)^{}=g(s)+i\frac{Q(s)}{4\pi \sqrt{s}}$$
(50)
since
$$\text{Im}g(s)=\frac{Q(s)}{8\pi \sqrt{s}}$$
(51)
we can write Eq.(48) as:
$$\left[I+K(s)g(s)\right]\mathrm{\Gamma }(s)=\left[I+K(s)g(s)^{}\right]\mathrm{\Gamma }(s)^{}.$$
(52)
This tells us that the quantity $`\left[I+K(s)g(s)\right]\mathrm{\Gamma }(s)`$ has no cuts since the only one which appears in $`g(s)`$ and $`\mathrm{\Gamma }(s)`$, the right or unitarity cut, is removed. Therefore, we can express $`\mathrm{\Gamma }(s)`$ as:
$$\mathrm{\Gamma }(s)=\left[I+K(s)g(s)\right]^1R(s)$$
(53)
with $`R(s)`$ being a vector of functions free of any singularity. We remark that this procedure of taking into account the final state interactions is based on the work presented in Ref.. In the following, we will fix $`R(s)`$ by requiring the matching of Eq.(53) to the next-to-leading order (one loop) CHPT $`\pi \pi `$ and $`K\overline{K}`$ scalar form factors. These are calculated in the next subsection.
It is worth to stress that Eq.(53), given in terms of a vector of functions $`R(s)`$ without any cut, can be applied to any K-matrix without unphysical cut contributions, as the one derived in ref. . The use of the strong amplitudes calculated from this reference is appealing for several reasons: 1) Because of their simplicity, 2) they have been already successfully used to describe many two meson production processes, as discussed in the introduction, and 3) higher order corrections to the kernel used in ref. are not necessary to match with the next-to-leading order CHPT scalar form factors. In fact, in ref. one can find a detailed comparison between the approach of ref. and the more general ones described in refs. and . The main conclusion is that, apart from the detail of including (or not) the $`\eta \eta `$ channel as already discussed, the unitarity corrections coming from the rescattering of the lowest order CHPT kernel completely dominate the strong S-wave $`I=0`$ scattering amplitudes up to about 1.2 GeV. Thus, one would not expect relevant departures from the use of the strong amplitudes from ref. or from refs. or . In fact, all these approaches give rise to very similar pole positions for the $`f_0(980)`$ and $`\sigma `$ mesons. For higher energies new effects have to be taken into account as e.g. the contributions from a pre-existing octet of scalar resonances around 1.4 GeV <sup>#5</sup><sup>#5</sup>#5Preexisting means here that these resonances with a mass around 1.4 GeV are as “elementary” as the basic fields $`\pi `$, $`K`$ or $`\eta `$. and the increasingly important role played by multiparticle states, basically the $`4\pi `$ intermediate state. In addition, one has to deal with more relevant interaction vertices between the various fields than those given in Eq.(1) as discussed in section II.
### D Matching with chiral perturbation theory
The general expression for the scalar form factors given in Eq.(53) can be further constrained by matching it to the one loop CHPT expression given in Eqs.(27, 33). This ensures that for energies where CHPT is applicable, these form factors fulfill all requirements given by chiral symmetry and the underlying power counting. This matching procedure essentially fixes the vector $`R(s)`$. We remark that since in our unitarization procedure we are not considering the $`\eta \eta `$ channel we thus can not reproduce the chiral logarithms associated with this channel. Therefore, we will only consider the contribution form this channel to the value of the form factors at $`s=0`$ and we will not include any $`s`$ dependence. This approximation should not induce any sizeable theoretical error because the influence of the $`\eta \eta `$ channel was found to be significant only above its threshold as already discussed in the introduction (when comparing the results of Ref. and Ref.).
We only discuss in detail the matching for the form factor $`\mathrm{\Gamma }_2^s(s)`$. The procedure for the other form factors is completely analogous and we thus only give the final results for them. From Eq.(53) one has:
$$\mathrm{\Gamma }^s(s)=\left[I+K(s)g(s)\right]^1R^s(s)=\left[IK(s)g(s)\right]R^s(s)+𝒪(p^4)$$
(54)
where the superscript “$`s`$” in $`R^s(s)`$ indicates that we are considering the $`\overline{s}s`$ form factor. From the former equation and Eq.(19) one has that $`R^s(s)_1=𝒪(p^2)`$ and that $`R^s(s)_2=1+𝒪(p^2)`$. Hence, we can recast Eq. (54) as:
$$\mathrm{\Gamma }^s(s)_2=R^s(s)_2K(s)_{22}g(s)_2+𝒪(p^4)=R^s(s)_2+\frac{3s}{4f^2}J_{KK}^r(s)+𝒪(p^4)$$
(55)
at the regularization scale $`\mu =1.2q_{\mathrm{max}}`$. Comparing this result with the one given in Eq.(33) leads to
$`R^s(s)_2`$ $`=`$ $`1+{\displaystyle \frac{4L_5^r}{f^2}}(s4m_K^2)+{\displaystyle \frac{8L_4^r}{f^2}}(s4m_K^2m_\pi ^2)+L_8^r{\displaystyle \frac{32m_K^2}{f^2}}+{\displaystyle \frac{16L_6^r}{f^2}}(4m_K^2+m_\pi ^2)+{\displaystyle \frac{2}{3}}\mu _\eta `$ (56)
$`+`$ $`{\displaystyle \frac{m_K^2}{36\pi ^2f^2}}(1+\mathrm{log}{\displaystyle \frac{m_\eta ^2}{\mu ^2}}),`$ (57)
using the Gell-Mann–Okubo relation $`m_\eta ^2=4m_K^2/3m_\pi ^2/3`$, the deviations from it being of higher order for our purpose. Proceeding in an analogous way for the other form factors one concludes:
$`R^n(s)_1`$ $`=`$ $`\sqrt{{\displaystyle \frac{3}{2}}}\{1+{\displaystyle \frac{4(L_5^r+2L_4^r)}{f^2}}s+{\displaystyle \frac{16(2L_8^rL_5)}{f^2}}m_\pi ^2+{\displaystyle \frac{8(2L_6^rL_4^r)}{f^2}}(2m_K^2+3m_\pi ^2)`$ (58)
$``$ $`{\displaystyle \frac{m_\pi ^2}{32\pi ^2f^2}}{\displaystyle \frac{1}{3}}\mu _\eta \},`$ (59)
$`R^n(s)_2`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\{1+{\displaystyle \frac{4L_5^r}{f^2}}(s4m_K^2)+{\displaystyle \frac{8L_4^r}{f^2}}(2s6m_K^2m_\pi ^2)+L_8^r{\displaystyle \frac{32m_K^2}{f^2}}+{\displaystyle \frac{16L_6^r}{f^2}}(6m_K^2+m_\pi ^2)+{\displaystyle \frac{2}{3}}\mu _\eta `$ (60)
$`+`$ $`{\displaystyle \frac{m_K^2}{72\pi ^2f^2}}(1+\mathrm{log}{\displaystyle \frac{m_\eta ^2}{\mu ^2}})\},`$ (61)
$`R^s(s)_1`$ $`=`$ $`\sqrt{3}\left\{{\displaystyle \frac{4L_4^r}{f^2}}(s2m_\pi ^2)+{\displaystyle \frac{16L_6^r}{f^2}}m_\pi ^2\right\}.`$ (62)
Notice that $`R^s(s)_1`$ is subleading in large $`N_c`$, i.e. of $`𝒪(N_c^1)`$, while the other quantities in Eqs.(56,58) are of order $`𝒪(1)`$ in this counting. This is expected since the production of pions from an $`\overline{s}s`$ source is subleading in large $`N_c`$ QCD . We also see, as already stressed in section III A that $`R^s(s)_1`$ is $`𝒪(p^2)`$ in the chiral counting. Once the functions $`R^{n,s}(s)`$ have been determined, the final expressions for the form factors are obtained by making use of Eq.(53, 38, 41). Finally, one has to take into account that, when using Eqs.(56, 58), the regularization scale is $`\mu =1.2q_{\mathrm{max}}1.08`$ GeV. Therefore, we have to run the low energy constants $`L_i^r(\mu )`$ to this scale from the usual ones $`\mu =m_\eta `$ or $`\mu =m_\rho `$, with $`m_\eta `$, $`m_\rho `$ the mass of the $`\eta `$, $`\rho `$ meson, respectively, by using the appropriate $`\beta `$–functions given in Ref..
## IV Results
We will first discuss the results for the $`J/\mathrm{\Psi }\varphi \pi \pi (K\overline{K})`$ decays and then we will also consider to some extent the $`J/\mathrm{\Psi }\omega \pi \pi `$ decay. To be more specific, we consider the S–wave contribution to these decay modes.
### A The $`\varphi `$-meson case
Considering the phase space of three particles we can write the unpolarized event distribution for the $`J/\mathrm{\Psi }\varphi \pi ^+\pi ^{}(K^+K^{})`$ reactions as:
$`{\displaystyle \frac{dN(W)_i}{dW}}`$ $`=`$ $`{\displaystyle \frac{𝒞_\varphi ^2G_i^2}{(2\pi )^312m_{J/\mathrm{\Psi }}^2}}|\mathrm{\Gamma }_i^s(s)+\lambda _\varphi \mathrm{\Gamma }_i^n(s)|^2\left[1+{\displaystyle \frac{(m_{J/\mathrm{\Psi }}^2+m_\varphi ^2W^2)^2}{8m_\varphi ^2m_{J/\mathrm{\Psi }}^2}}\right]`$ (63)
$`\times `$ $`\sqrt{[W^24m_i^2]\left[(m_{J/\mathrm{\Psi }}^2W^2m_\varphi )^24m_\varphi ^2W^2\right]},`$ (64)
where $`i=1`$ refers to the $`\pi ^+\pi ^{}`$ and $`i=2`$ to the $`K^+K^{}`$ system, in order. Furthermore, $`W`$ is the total energy in the c.m. of the two pions or kaons, $`G_i`$ is basically a Clebsch-Gordan coefficient equal to 4/3 for pions and 1/2 for kaons, respectively, and $`𝒞_\varphi `$ a normalization constant depending on the experiment, in our case DM2 or MARK-III . In comparing with the experimental data, we will average Eq.(63) over the width of the bin (as given by the corresponding experiment). As discussed in section III D, our calculated form factors depend on the CHPT low energy constants $`L_4^r`$, $`L_5^r`$, $`L_6^r`$ and $`L_8^r`$. From these, only $`L_5^r`$ and $`L_8^r`$ are relatively well determined. Their most recent values, given in Ref. from an $`𝒪(p^6)`$ CHPT analysis of the $`K_{\mathrm{}\mathrm{\hspace{0.17em}4}}`$ form factors, are:
$$10^3L_5^r(M_\rho )=0.65\pm 0.12,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}10}^3L_8^r(M_\rho )=0.48\pm 0.18.$$
(65)
At the scale $`\mu =1.2q_{\mathrm{max}}1.08`$ GeV they are:
$$10^3L_5^r(1.08\mathrm{GeV})=0.15\pm 0.12,\mathrm{\hspace{0.33em}\hspace{0.33em}\hspace{0.33em}10}^3L_8^r(1.08\mathrm{GeV})=0.26\pm 0.18,$$
(66)
On the other hand, $`L_4^r`$ and $`L_6^r`$ are only poorly known and their present values can be considered as stemming more from an estimation of their order of magnitude than from a truly phenomenological fit. According to Ref., it is estimated that for a regularization scale $`\mu `$ between 0.5 and 1.0 GeV one should have $`10^3L_4^r\pm 0.5`$ and $`10^3L_6\pm 0.3`$. A more recent determination gives $`10^3L_4^r(m_\rho )=0.3\pm 0.5`$ and $`10^3L_6^r(m_\rho )=0.2\pm 0.3`$ so that at a scale of 1.08 GeV one has
$$10^3L_4^r(1.08\mathrm{GeV})=0.57\pm 0.5,10^3L_6^r(1.08\mathrm{GeV})=0.36\pm 0.3.$$
(67)
Again, this estimate relies on OZI (large $`N_c`$) arguments. To be more precise, one sets $`L_{4,6}^r`$ to zero at the scale $`\mu =m_\eta `$, which is of course somewhat arbitrary. One can also make use of information about the low energy coupling constant $`\mathrm{}_4^r`$ coming from two flavor CHPT by means of the relation , $`\mathrm{}_4^r(\mu )=8L_4^r(\mu )+4L_5^r(\mu )\nu _K/2+𝒪(p^6)`$ with $`\nu _K=[\mathrm{ln}(M_K^2/\mu ^2)+1]/32\pi ^2`$. The low energy coupling constant $`\mathrm{}_4^r`$ has been determined at $`𝒪(p^4)`$ in ref. with the result $`10^3\mathrm{}_4^r=1.2\pm 6`$ at the scale $`\mu =1.08`$ GeV. However, in ref. making use of the analytically deduced pion scalar form factor $`\overline{u}u+\overline{d}d`$ up to and including $`𝒪(p^6)`$, they update the previous value and give the improved result $`10^3\mathrm{}_4^r=1.8\pm 1.9`$ at the same scale. Although the central value from both determinations is very similar the error is much smaller in the second case. With the value for $`L_5^r`$ given in Eq. (66), we obtain $`10^3L_4^r(1.08\mathrm{GeV})=0.1\pm 0.7`$ when using $`\mathrm{}_4^r`$ from ref. and
$$10^3L_4^r(1.08\mathrm{GeV})=0.19\pm 0.25$$
(68)
in the case of $`𝒪(p^6)`$ SU(2) CHPT . The value obtained in this way for $`L_4^r`$ is also consistent with zero, but on the positive side, in stark contrast to the value given in Eq.(67) whose central value is very far from the more precise determination coming from refs. as given in Eq.(68). Thus, although at present no precise determination of this low energy constant is available, rather strong constraints on its value can be obtained by combining the determination of the low energy constants making use of $`𝒪(p^6)`$ CHPT, both in its SU(2) and SU(3) forms.
More recently, some constraints on the couplings $`L_4^r`$ and $`L_6^r`$ have been reported . The determination of $`L_4^r`$ relies on a comparison of the CHPT series at next- and next-to-next-leading order with a phenomenologically determined scalar form factor. In ref. the value $`10^3L_4^r(1.08\mathrm{GeV})0.14`$ is given without errors and the band of values $`0.1210^3L_4^r(1.08\mathrm{GeV})<0.04`$ is reported in ref. . We note that the second set of values for $`L_4^r`$ is in the lower limit of the value for $`L_4^r`$ given in Eq.(68). Nevertheless, the determination of $`L_6^r`$ is not so well sounded due to strong simplifying working assumptions when computing phenomenologically the quark correlator $`(\overline{u}u+\overline{d}d)\overline{s}s`$ . It is also stated in ref. that the positivity of the fermionic measure gives rise to a lower bound for $`L_6^r`$, $`10^3L_6^r(1.08\mathrm{GeV})0.03`$. However, this bound is somewhat arbitrary since the only necessary requirement to make use of the positivity of the fermionic measure is that all the three light quark masses have to be equal. In ref., they were set equal to the strange quark mass, but they could as well have been taken on another value. In fact, for the average light quark masses $`(m_u+m_d)/2`$, the corresponding lower bound is: $`10^3L_6^r(1.08\mathrm{GeV})0.75`$. The bound based on using the strange quark mass can only be maintained if one assumes the next–to–leading order corrections to be of canonical size, $`cM_K^2/(4\pi f_\pi )^2`$, with $`c`$ a number of order one.<sup>#6</sup><sup>#6</sup>#6We are grateful to Bachir Moussallam for a claryfing discussion on this topic.
While the data of DM2 have been published, this is not the case for the data of MARK-III . On the other hand, both experiments, see Figs.7 and 8, are compatible for the $`\pi ^+\pi ^{}`$ event distribution for 25 and 10 MeV bins<sup>#7</sup><sup>#7</sup>#7The data for the 10 MeV bins of MARK-III has been taken from Ref... However, this is not the case for the $`K\overline{K}`$ event distribution, see Fig.9. Consequently, we will only consider for our fits the data from DM2 for the $`\pi ^+\pi ^{}`$ event distributions in the $`J/\mathrm{\Psi }\varphi \pi ^+\pi ^{}`$ decay, both for the 10 and the 25 MeV bins. In fitting the data for the $`J/\mathrm{\Psi }\varphi \pi ^+\pi ^{}`$ distribution from DM2 we will fix $`L_5^r`$ and $`L_8^r`$ at the values given in Eq.(66). On the other hand, $`L_4^r`$ and $`L_6^r`$ will be taken as free parameters. In this way, our expression for the event distribution of the pions and kaons will have four free parameters: $`𝒞_\varphi `$, $`\lambda _\varphi `$, $`L_4^r`$ and $`L_6^r`$. However, there is still too much freedom in fitting the data with this set of free parameters since the global normalization constant $`𝒞_\varphi `$ can only be determined within a large range and with a sizeable uncertainty. We will further restrict our fit by requiring that we can also describe the pion event distribution in the $`J/\mathrm{\Psi }\omega \pi ^+\pi ^{}`$ decay, at least in the low energy region where the S-wave contribution, the one which we are considering here, is dominant. This extension of our model to the $`\omega `$ case is discussed in the next subsection. Imposing this requirement, $`𝒞_\varphi `$ is fixed,<sup>#8</sup><sup>#8</sup>#8There is approximately a factor 1.11 between the global normalization constant required for the DM2 data with respect the one required for MARK-III. $`𝒞_\varphi =(16\pm 3)`$MeV<sup>-1</sup>, and the fitted values for $`\lambda _\varphi `$, $`L_4^r(1.08\mathrm{GeV})`$ and $`L_6^r(1.08\mathrm{GeV})`$ are:
$$\lambda _\varphi =0.17\pm 0.06,10^3L_4^r(1.08\mathrm{GeV})=0.44\pm 0.11,10^3L_6^r(1.08\mathrm{GeV})=0.38\pm 0.06,$$
(69)
with a $`\chi ^2/\mathrm{dof}=0.92`$. Clearly $`\lambda _\varphi 0`$, in contradiction with the OZI rule. Furthermore, the pion event distribution turns out to be very sensitive to the large $`N_c`$ subleading low energy constants $`L_4^r`$ and $`L_6^r`$. The theoretical uncertainties given in Eq.(69) are obtained in the following way. We have allowed for a relative change of 20% in the global normalization constant $`𝒞_\varphi `$ when considering the data with the $`\varphi `$ and also the ones with the omega in the final state. We consider this estimate of the error in $`𝒞_\varphi `$ as conservative, since the ensuing deviation from the $`\omega \pi ^+\pi ^{}`$ data for such changes in $`𝒞_\varphi `$ is larger than the uncertainty in the omega data by assuming a Poisson distribution. On the other hand, we have also allowed an uncertainty of $`\pm 0.1`$ GeV in the determination of $`q_{\mathrm{max}}`$ from ref. and then we calculated the band of values for $`\lambda _\varphi `$, $`L_4^r`$ and $`L_6^r`$. All these sources of uncertainty are added in quadrature together with the statistical error given by the fitting procedure when using the central values for $`𝒞_\varphi `$ and $`q_{\mathrm{max}}`$.
It is instructive to compare the values for the low energy constants found here, e.g. Eqs.(69), with the ones given in Eq.(67). While $`L_6^r`$ agrees perfectly within error bars, the sign of $`L_4^r`$ is changed. Stated differently, if we evaluate from Eq. (69) $`10^3L_4^r`$ at the rho mass, we find a value of $`0.71`$, which is sizeably larger in magnitude than the central value given in Ref.. Also, it is larger than the positive value deduced from SU(2) information given in Eq.(68), at the scale $`\mu =m_\rho `$ one has for this case $`10^3L_4^r=0.46\pm 0.25`$, although both values are consistent within errors. We reiterate that using large $`N_c`$ arguments, one would expect $`L_4^r`$ to be zero (at a scale somewhere in the resonance region). Therefore, our increased value and also the one from Eq.(68), clearly signals OZI violation. Quite differently, our value for $`10^3L_6^r(m_\rho )=0.22`$ is completely consistent with the previous determination . However, it does not fulfill the positivity constraint $`10^3L_6^r(m_\rho )0.20`$ but it fulfills the other reasonable ‘positivity’ bound previously discussed $`10^3L_6^r(m_\rho )0.59`$. In fact, our value for $`10^3L_6`$ lies in an natural intermediate region between both extreme lower bound. Nevertheless, we have also performed a series of fits enforcing the former constraint. We can fit the $`\varphi `$ data, but on the expense of very large values for $`\lambda _\varphi `$ and $`L_4^r`$. Furthermore, it is not possible to simultaneously get a description of the $`\omega `$ decay data. We think that further study is needed in order to apply the positivity of the Dirac measure and also, we should stress that the LEC $`L_6^r`$ is plagued by the Kaplan–Manohar ambiguity . It is important to point out that one can criticize our determinations of the low energy constants for two reasons. First, our model for the $`J/\mathrm{\Psi }`$ decay with the $`\varphi `$ meson as a spectator is fairly simple, one could e.g. write down higher order transition operators which would complicate the analysis. Given, however, the fact that we can precisely reproduce the data both for the $`\varphi `$ and the $`\omega `$ resonances, it is not obvious a priori that such a modified ansatz would lead to very different results. Second, the use of unitarity to determine the scalar form factors beyond one loop accuracy induces some inevitable model dependence. To overcome this, one could think of doing a pure CHPT analysis on the left wing of the scalar resonance. We believe, however, that the present data in this energy region are not precise enough for an accurate determination of the LECs. Independently of these reservations, our analysis clearly underlines that the OZI rule is strongly violated in the scalar $`0^{++}`$ sector, as indicated e.g. by the large positive value of the LEC $`L_4^r(m_\rho )`$ and also by the non-vanishing value of $`\lambda _\varphi `$. With respect the latter point, see also the footnote in section IV B.
It is also worth to indicate that in ref., making use of the Inverse Amplitude Method (IAM) with complete next-to-leading order CHPT strong amplitudes, a fit to the $`I=0`$ and $`2`$ S-wave and $`I=1`$ P-wave $`\pi \pi `$ and $`K\overline{K}`$ partial wave amplitudes was done in terms of the low energy coupling constants $`L_1^r`$, $`L_2^r`$, $`L_3^r`$, $`L_4^r`$, $`L_5^4`$ and $`2L_6^r+L_8^r`$. This study has in common with the present one that a complete matching to the relevant next-to-leading order CHPT results was given and at the same time fully unitarity amplitudes were derived. The experimental data was very well reproduced up to energies around 1.2 GeV giving rise to the presence of the resonances $`\rho `$ and $`f_0(980)`$. The values obtained for $`L_5^r`$ and $`2L_6^r+L_8^r`$, within errors, are consistent with those recently obtained in ref. . This implies agreement of the results of that reference with our choice for the values of $`L_5^r`$ and $`L_8^r`$ and our presently determined value for $`L_6^r`$. The main difference between the set of values given in ref. and those in ref. corresponds to the value of $`L_2^r`$. While in the former case $`L_2^r2L_1^r`$ as required by Vector Meson Dominance (VMD) <sup>#9</sup><sup>#9</sup>#9There is a very close link between VMD and the IAM in the vector channels ., this relation is only fulfilled within errors by the values given in ref.. Nevertheless, in ref. $`10^3L_4^r(m_\rho )=0.2\pm 0.1`$. This value, although on the positive side, is incompatible with our present one, $`10^3L_4^r(m_\rho )=0.71\pm 0.11`$, and, within errors, is compatible with the rest of analyses presented in this section except for . Summarizing, for $`L_6^r`$ there is a rather good agreement between our present study and refs., , in disagreement with the finding of ref. . On the other hand, for $`L_4^r`$ our present analysis is compatible only with that value of $`L_4^r`$ determined from $`𝒪(p^6)`$ CHPT which also find quite a sizeable central value at $`\mu =M_\rho `$, around $`0.5\times 10^3`$, different from the smaller numbers of refs..
The resulting non–strange and strange normalized scalar form factors of the pion and the kaon are shown in Fig.4, Fig.5 and Fig.6. In the case of the non–strange scalar form factor, we show for comparison in Fig.4 the one– and two–loop CHPT as well as the dispersion theoretical results and the exponentiated two–loop CHPT result . In the latter case, the two–loop CHPT result is improved by making use of an Omnès resummation in terms of the next-to-leading order CHPT phase shifts. Our result is close to the ones obtained by a different method in Ref. and even closer to the exponentiated two–loop CHPT results of ref. . The agreement is worse when comparing our results with the so called modified-Omnès representation of refs. . We also remark that the two–loop representation covers the main feature of this quantity below $`W600`$ MeV, as it is known since long . The strange scalar form factor of the pion is reasonably well described for energies below 350 MeV. In contrast, the strange and the non–strange scalar form factor of the kaon are poorly described at one loop, as expected from the larger mass of the kaon.
In Figs.7 and 8, we show the curves from ours fit to the $`\pi ^+\pi ^{}`$ event distribution in comparison with the experimental data from DM2 and MARK-III for the 25 and the 10 MeV bins, respectively. The data of MARK-III have been multiplied by the factor $`𝒞_{\varphi ,\mathrm{DM2}}^2/𝒞_{\varphi ,\mathrm{MARK}\mathrm{III}}^21.11^2`$ in order to facilitate the comparison between both sets of data. The agreement with the experimental data is very good as indicated by the low $`\chi ^2`$/dof of 0.92. By comparing Fig.7 with the left panel of Fig.6 we see that the event distribution of the two pions is dominated by the strange scalar form factor of the pion. Moreover, in Fig.7 the changes in the results when allowing a change in the cut-off $`q_{\mathrm{max}}`$ by $`\pm 0.1`$ GeV are shown to be quite small. In Fig.9 our prediction for the $`K^+K^{}`$ event distribution is depicted. Incidentally, we find better agreement with the data of MARK-III than with the ones of DM2 with respect to this decay mode. This was also noted in Refs.. In fact, in these references a fit of similar quality to the data of DM2 and MARK-III is also given. The important difference between their method and ours is that we have devised a dynamical approach which means that the parameters that enters in our description of the problem are not specific to it and can be related to many other physical observables. This is particularly true for $`L_4^r`$ and $`L_6^r`$. But even for $`\lambda _\varphi `$ we will see in the next subsection how it can be related to the whole set of U(3) processes that follows from the decays of the $`J/\mathrm{\Psi }`$ resonance to any vector resonance belonging to the lightest nonet of vector resonances and two pseudoscalars. On the other hand, in ref. no attempt was done to describe the $`J/\mathrm{\Psi }\omega \pi \pi `$ decays.
### B The $`\omega `$-meson case
A priori one can expect that the $`J/\mathrm{\Psi }\omega \pi \pi `$ decay requires a rather different dynamical description than that for the mode $`J/\mathrm{\Psi }\varphi \pi \pi `$ considered so far. For instance, the approach of considering the $`\omega `$ as a spectator is by no means so clear as for the $`\varphi `$ case. Note that the Dalitz plot for this decay has very clear bands due to the decays $`b_1(1235)\omega \pi `$ and $`f_2(1270)\pi ^+\pi ^{}`$ . The latter induces a sizeable D-wave contribution so that our approach can only be applied for the first few hundred MeV of the two pion event distribution. However, making use of SU(3) symmetry, we can extend our considerations from section II and we will present our calculated S-wave contribution to the invariant mass distribution of the pions in the $`J/\mathrm{\Psi }\omega \pi \pi `$ decay. This calculation is completely fixed in terms of the parameters already given for the $`\varphi `$ case, cf. Eq.(69), except for the global normalization constant for which we have also used the experimental data from the $`J/\mathrm{\Psi }\omega \pi ^+\pi ^{}`$ decays to further constraint its value.
Let us now be more specific and work out the aforementioned relation between the two cases. An invariant SU(3) Lagrangian involving the octet and singlet of vector resonances $`V_\mu ^{(8)}`$ and $`V_\mu ^{(1)}`$, respectively, and the corresponding ones of the scalar sources $`S^{(8)}`$ and $`S^{(1)}`$ can be written as (making also use of Lorentz invariance)
$$=\widehat{g}\left(\mathrm{\Psi }^\sigma V_\sigma ^{(8)}S^{(8)}+\nu \mathrm{\Psi }^\sigma V_\sigma ^{(1)}S^{(1)}\right),$$
(70)
where $`\widehat{g}`$ is an overall coupling constant whose precise value is not needed in the following. We only need to determine the relative strength of the octet to singlet couplings given in terms of the real parameter $`\nu `$. The symbol $`\mathrm{}`$ refers to the trace over the SU(3) indices of the matrices $`V^{(8)}`$ and $`S^{(8)}`$. These are defined via
$$V_\sigma ^{(8)}=\left[\begin{array}{ccc}\frac{1}{\sqrt{2}}\rho ^0+\frac{1}{\sqrt{6}}V_8& \rho ^+& K^+\\ \rho ^{}& \frac{1}{\sqrt{2}}\rho ^0+\frac{1}{\sqrt{6}}V_8& K^{\mathrm{\hspace{0.17em}0}}\\ K^{}& \overline{K}^{\mathrm{\hspace{0.17em}0}}& \frac{2}{\sqrt{6}}V_8\end{array}\right]_\sigma .$$
(71)
and similarly for the $`S^{(8)}`$ matrix. In the previous equation we have denoted by $`V_8`$ the $`I=0`$ state of the octet of vector resonances. This formalism is in close analogy with the one used in CHPT for the octet of pseudoscalars, compare Eq.(13). Denoting by $`S_8`$ the $`I=0`$ operator of the octet of scalar sources, we can write the terms involving $`V_8`$ and $`V_1`$ of Eq.(70) as:
$$\mathrm{\Psi }^\sigma \left(V_{8;\sigma }S_8+\nu V_\sigma ^{(1)}S^{(1)}\right).$$
(72)
Considering ideal mixing<sup>#10</sup><sup>#10</sup>#10In Ref. the departure from the ideal mixing in the $`\omega \varphi `$ system is thoroughly studied comparing different models, and is described by a rotation of the ideal mixing states with a rotation angle $`|\delta _V|3^0`$. This departure would produce corrections of the order of a 5$`\%`$ with respect the ideal mixing situation considered here. However, this value should be compared with the departure from 1 of the parameter $`\nu `$ which from Eq. (77), taking into account the value of $`\lambda _\varphi `$ given in Eq. (69), is about $`40\%`$, a much bigger effect than any expected deviation from the ideal mixing situation in the $`\omega \varphi `$ system. Note that $`\nu =1`$ is the value expected from U(3) symmetry and for this case $`\lambda _\varphi =0`$ and $`\lambda _\omega =\mathrm{}`$, see Eqs. (77,78) between the $`V_8`$ and the $`V^{(1)}`$ then:
$$V_8=\frac{\omega }{\sqrt{3}}\sqrt{\frac{2}{3}}\varphi ,V^{(1)}=\sqrt{\frac{2}{3}}\omega +\frac{\varphi }{\sqrt{3}}.$$
(73)
In an analogous way we will also introduce the scalar sources $`S_\omega `$ and $`S_\varphi `$ defined by
$$S_8=\frac{S_\omega }{\sqrt{3}}\sqrt{\frac{2}{3}}S_\varphi ,S^{(1)}=\sqrt{\frac{2}{3}}S_\omega +\frac{S_\varphi }{\sqrt{3}}.$$
(74)
Note that in a quark model language, consistently with the transformation properties under SU(3), we can write:
$$S_\varphi =\overline{s}s\text{and}S_\omega =\frac{1}{\sqrt{2}}(\overline{u}u+\overline{d}d).$$
(75)
Rewriting Eq.(72) in terms of $`\omega `$, $`\varphi `$, $`S_\omega `$ and $`S_\varphi `$, we have:
$$\frac{2+\nu }{3}\mathrm{\Psi }^\sigma \varphi _\sigma \left(S_\varphi +S_\omega \frac{\sqrt{2}(\nu 1)}{2+\nu }\right)+\frac{\sqrt{2}(\nu 1)}{3}\mathrm{\Psi }^\sigma \omega _\sigma \left(S_\varphi +S_\omega \frac{1+2\nu }{\sqrt{2}(\nu 1)}\right).$$
(76)
In this way, the parameter $`\lambda _\varphi `$ introduced in Eqs.(3,4) and fitted in the previous subsection, can now be expressed as:
$$\lambda _\varphi =\frac{\sqrt{2}(\nu 1)}{2+\nu }.$$
(77)
¿From this equation we can isolate $`\nu =\nu (\lambda _\varphi )`$ and then predict the corresponding parameter $`\lambda _\omega `$,
$$\lambda _\omega =\frac{1+2\nu (\lambda _\varphi )}{\sqrt{2}(\nu (\lambda _\varphi )1)}$$
(78)
We can also obtain from Eq.(76) the global normalization constant for the $`\omega `$ in terms of the one of the $`\varphi `$ since
$$\frac{C_\omega }{C_\varphi }=\frac{\sqrt{2}(\nu 1)}{2+\nu }=\lambda _\varphi .$$
(79)
In Fig.10 we show the calculated S-wave contribution to the event distribution of pions in the $`\mathrm{\Psi }\omega \pi ^+\pi ^{}`$ decay. This calculation does not introduce any new free parameter since all of them have been fixed in terms of the one of the $`\varphi `$ case by Eqs.(78) and (79). The description of the data of MARK-III up to around 0.7 GeV is very good. For higher energies the D-wave contribution cannot be further neglected.
## V Conclusions
In this work we have addressed the problem of the $`J/\mathrm{\Psi }`$ decays into a vector ($`\varphi ,\omega `$) and two pseudoscalar mesons (Goldstone bosons) measured at DM2 and MARK-III . These processes are considered to be mediated by the corresponding scalar form factors of the pseudoscalar mesons if one considers the emitted vector meson as a spectator. Consequently, these reactions are rather interesting since they are very sensitive to OZI violation physics, in our scheme parameterized by the constants $`\lambda _\varphi `$, see Eq.(4), and the low energy constants $`L_4^r`$ and $`L_6^r`$ of chiral perturbation theory. The first of these constants parameterizes the direct admixture of non–strange quarks to the scalar interpolating field for our model of the $`J/\mathrm{\Psi }`$ decay with the $`\varphi `$ playing the role of a spectator, see Figs.1,2. The two low energy constants enter the one loop description of the pion and kaon scalar form factors. To describe these properly for the range of energies relevant here, we have combined information coming from next-to-leading order (one loop) chiral perturbation theory (CHPT) with the unitarity requirements which are valid to all orders in the chiral expansion. In addition, we also have calculated for the first time the next-to-leading order CHPT kaon scalar form factors, for strange, $`\overline{s}s`$, and non–strange, $`\overline{u}u+\overline{d}d`$, scalar–isoscalar quark densities. The unitarity requirements were imposed by using the strong $`I=0`$ $`\pi \pi `$ and $`K\overline{K}`$ amplitudes derived in Ref.. The amplitudes given in that paper not only describe accurately the S-wave $`I=0`$ and $`I=1`$ strong scattering data but also have been used to successfully reproduce or even predict experimental data for the whole set of reactions listed in the Introduction. With this input, we have successfully described, from threshold up to around 1.2 GeV,<sup>#11</sup><sup>#11</sup>#11In principle, one could also go to higher energies but that is much more difficult due a) to the appearance of multiparticle states and b) due to novel interaction vertices as discussed in section II. the event distribution of the $`\pi ^+\pi ^{}`$ system in the $`J/\mathrm{\Psi }\varphi \pi ^+\pi ^{}`$ decay. We have then predicted, in agreement with the data from MARK-III, the event distribution of kaons in the $`J/\mathrm{\Psi }\varphi K^+K^{}`$ reaction and the low energy part, where the S-wave dominates, of the event distribution of $`\pi ^+\pi ^{}`$ pairs in the $`J/\mathrm{\Psi }\omega \pi ^+\pi ^{}`$ decay. Furthermore, the OZI violation parameter $`\lambda _\varphi `$ comes out different from zero. This also holds for the low energy constants $`L_4^r`$ and $`L_6^r`$. While the value of the latter agrees with previous estimates , our result for $`L_4^r`$ is sizeably larger in magnitude as most previous estimations . However, it is compatible within errors with the quite constraint value derived by combining the information from $`𝒪(p^6)`$ SU(2) and SU(3) CHPT. This offers another indication that the OZI rule does not account for the physics in the scalar $`0^{++}`$ channel, as stressed e.g. in Refs.. The scheme employed here offers a unique approach to describe the scalar sector, which has been at the heart of many investigations over the last decade.
Clearly, to further improve the approach presented here, it would be mandatory to not only have event distributions but rather normalized data. This would allow one to pin down the low energy constants $`L_4^r`$ and $`L_6^r`$ more precisely together with the OZI violation parameter $`\lambda _\varphi `$ as well as the product $`(\widehat{g}m_qB_0)^2`$ (i.e. normalization of the scalar form factors times the strength of the scalar source to vector meson coupling). With respect to the former, one could also try to do a pure one (or even two) loop CHPT calculation for small two–pion invariant masses, i.e. on the left wing of the $`f_0(980)`$ resonance. Clearly, the presently available data are not precise enough for successfully doing that, but such a program is in principle of the similar interest as the study of chiral dynamics in $`\tau `$ decays, see e.g. Refs., specially when referring to the scalar sector.
### Acknowledgments
J.A.O. would like to acknowledge stimulating discussions with T. Barnes. U.-G.M. is grateful to V. Bernard for some pertinent comments. The work of J.A.O. was supported in part by funds from DGICYT under contract PB96-0753 and from the EU TMR network Eurodaphne, contract no. ERBFMRX-CT98-0169.
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# Uniform ergodic theorems on subshifts over a finite alphabet
## 1. Introduction
Ergodic theorems for addditive and subadditive functions play a role in several branches of mathematics and physics. In particular, they are an important tool in statistical mechanics as well as in the theory of random operators (cf. and references therein).
During recent years lattice gas models and random operators on aperiodic tilings have received a lot of interest both in one dimension and in higher dimensions (cf. and references therein). In these cases one rather expects uniform ergodic theorems to hold.
The aim of this paper is to provide a thorough study of the validity of such theorems in the one-dimensional case.
In particular, we show that every strictly ergodic subshift over a finite alphabet admits a uniform additive ergodic theorem. Moreover, we give a necessary and sufficient condition for a minimal subshift to allow for a subadditive theorem. This gives in particular a sufficient condition for unique ergodicity.
The proofs are quite elementary and conceptual. Thus, it is to be expected that a considerable part of the material presented here, can be extended to higher dimensional tiling dynamical systems.
More precisely, we consider the following situation:
Let $`A`$ be a finite set called the alphabet and equipped with the discrete topology. Let $`\mathrm{\Omega }`$ be a subshift over $`A`$. This means that $`\mathrm{\Omega }`$ is a closed subset of $`A^{}`$, where $`A^{}`$ is given the product topology and $`\mathrm{\Omega }`$ is invariant under the shift operator $`T:A^{}A^{}`$, $`Ta(n)a(n+1)`$.
We consider sequences over $`A`$ as words and use standard concepts from the theory of words (). In particular, $`\text{Sub}(w)`$ denotes the set of subwords of $`w`$, the empty word is denoted by $`ϵ`$, the number of occurrences of $`v`$ in $`w`$ is denoted by $`\mathrm{}_v(w)`$ and the length $`|w|`$ of the word $`w=w_1\mathrm{}w_n`$ is given by $`n`$. To $`\mathrm{\Omega }`$ we associate the set $`𝒲=𝒲(\mathrm{\Omega })`$ of finite words associated to $`\mathrm{\Omega }`$ given by $`𝒲_{\omega \mathrm{\Omega }}\text{Sub}(\omega )`$. A word $`w𝒲`$ is called primitive if $`v`$ can not be written as $`v=w^l`$, with $`w𝒲`$ and $`l2`$. For a finite set $`M`$, we define $`\mathrm{}M`$ to be the number of elements in $`M`$.
To phrase our additive ergodic theorem, we need the following definition.
###### Definition 1.1.
Let $`(B,||)`$ be a Banach space. A function $`F:𝒲B`$ is called additive if there exists a constant $`D>0`$ and a nonincreasing function $`c:[0,\mathrm{})[0,\mathrm{})`$ with $`lim_r\mathrm{}c(r)=0`$ s.t. the following holds
* $`|F(v)_{j=1}^nF(v_j)|_{j=1}^nc(|v_j|)|v_j|`$ for $`v=v_1\mathrm{}v_n𝒲`$.
* $`|F(v)|D|v|`$ for every $`v𝒲`$.
Then the additive theorem can be stated as follows.
###### Theorem 1.
Let $`(\mathrm{\Omega },T)`$ be a minimal subshift over the finite alphabet $`A`$ with associated set of words $`𝒲=_{\omega \mathrm{\Omega }}\text{Sub}(\omega )`$. Then the following are equivalent:
* $`(\mathrm{\Omega },T)`$ is uniquely ergodic, i.e. $`lim_{|w|\mathrm{}}\frac{\mathrm{}_v(w)}{|w|}`$ exists for all $`v𝒲.`$
* The limit $`lim_{|w|\mathrm{}}\frac{F(w)}{|w|}`$ exists for every Banach-space-valued additive function $`F`$ on $`𝒲`$.
* The limit $`lim_n\mathrm{}\frac{1}{n}_{j=1}^nf(T^j\omega )`$ exists uniformly in $`\omega \mathrm{\Omega }`$ for every continuous Banach-space-valued function $`f`$ on $`\mathrm{\Omega }`$.
Here, the proof of (iii) $``$ (i) is well known. The proof of (ii) $``$ (iii) follows by rather standard arguments (cf. ) and is in fact valid in arbitrary dimensions. So, the hardest part is the proof of (i) $``$ (ii). There, we use the following strategy:
* A weak form of hierarchic structure exists in arbitrary minimal subshifts.
* This structure suffices to prove a local ergodic theorem such as (ii).
Here, the hierarchy found in the first step generalizes the notion of $`n`$-partition in Sturmian dynamical systems investigated in (cf. as well for application to one-dimensional quasicrystals). It is based on the concept of return word recently introduced by Durand . The second step then uses ideas of .
To introduce the second result of this paper, recall that a function $`F:𝒲`$ is called subadditive if it satisfies $`F(ab)F(a)+F(b)`$. The dynamical system $`(\mathrm{\Omega },T)`$ is said to satisfy (SET), i.e. to admit a uniform subadditive ergodic theorem, if, for every subadditive function $`F`$, the limit $`lim_{|w|\mathrm{}}\frac{F(w)}{|w|}`$ exists and equals $`inf_nF^{(n)}`$, where $`F^{(n)}\mathrm{max}\{\frac{F(v)}{|v|}:v𝒲,|v|=n\}`$.
Define the functions $`l_v:𝒲`$, for $`v𝒲`$, and $`\nu :𝒲`$ by
$$l_v(w)(\text{Maximal number of disjoint copies of }v\text{ in }w)|v|$$
and
$$\nu (v)\underset{|w|\mathrm{}}{lim\; inf}\frac{l_v(w)}{|w|}.$$
Then a subshift $`(\mathrm{\Omega },T)`$ over $`A`$ is said to satisfy uniform positivity of quasiweights (PQ) if the following condition holds:
* There exists a a constant $`C>0`$ with $`\nu (v)C`$ for all $`v𝒲`$.
Our result on subadditive ergodic theorems then reads as follows:
###### Theorem 2.
Let $`(\mathrm{\Omega },T)`$ be a minimal subshift. Then the following are equivalent:
(i) The subshift $`(\mathrm{\Omega },T)`$ satisfies (SET).
(ii) The subshift $`(\mathrm{\Omega },T)`$ satisfies (PQ).
###### Remark 1.
(a) As the validity of a subadditive ergodic theorem immediately implies that the underlying subshift is uniquely ergodic, we see that (PQ) is a sufficient condition for unique ergodicity.
(b) It is not hard to prove (PQ) for subshifts which satisfy the uniform positivity of weights (PW) as well as a highest power condition (HP) given as follows:
* There exists $`E>0`$ with $`lim\; inf_{|w|\mathrm{}}\frac{\mathrm{}_v(w)}{|w|}|v|E`$ for all $`v𝒲`$.
* There exists an $`N>0`$ s.t. $`v^k𝒲`$ implies $`kN`$ (or equivalently: there exists a $`\kappa >0`$ s.t. $`v`$ prefix of $`uv`$ with $`u`$ not empty implies $`|u|\frac{|v|}{\kappa }`$ ).
In particular, (PQ) holds if $`(\mathrm{\Omega },T)`$ is linearly repetitive i.e. if there exists a $`D>0`$ s.t. every $`v𝒲`$ is a factor of every $`w𝒲`$ with $`|w|D|v|`$. This is further discussed in Section 5.
(c) It is not hard to see that (PQ) (or (PW)) implies in particular the minimality of the subshift $`(\mathrm{\Omega },T)`$. In fact, minimality of $`(\mathrm{\Omega },T)`$ is equivalent to $`\nu (v)>0`$ for every $`v𝒲`$. From this point of view (PQ) (or (PW)) can be seen as a very strong minimality condition.
The organisation of the paper is as follows. In Section 2 we review basic notions from the theory of subshifts over finite alphabets. The additive ergodic theorem is then contained in Section 3. The subadditive ergodic theorem is contained in Section 4. Finally, in Section 5 we provide several examples of subshifts satisfying (PQ), thereby giving a precise form to Remark 1 (b).
## 2. Preliminaries
Let $`\mathrm{\Omega }`$ be a subshift over the finite alphabet $`A`$ with associated set $`𝒲`$ of finite words. The following two facts are well known (cf. ):
* The subshift $`\mathrm{\Omega }`$ is minimal, i.e. each orbit is dense if and only if for every $`w𝒲`$ there exists an $`R(w)>0`$ s.t. $`w`$ is a factor of every $`v𝒲`$ with $`|v|R(w)`$.
* The subshift is uniquely ergodic if and only if the limit $`lim_{|w|\mathrm{}}\frac{\mathrm{}_v(w)}{|w|}`$ exists for every $`v𝒲`$.
A return word of $`u𝒲`$ is a word $`w𝒲`$ s.t.
$$wu𝒲,\mathrm{}_u(wu)=2,u\text{ is prefix of }wu.$$
This notion was introduced by Durand in . The set of return words of $`u𝒲`$ will be denoted by $`(u)(u,𝒲)`$. Similarly, one can define for $`n`$ the set $`^n_{v𝒲,|v|=n}(v)`$. If $`(\mathrm{\Omega },T)`$ is minimal, the length of a return word of $`u𝒲`$ is bounded by $`R(u)`$. In particular, the set $`(u)`$ is finite for every $`u𝒲`$ if $`𝒲`$ is minimal.
Set $`m(n)\mathrm{min}\{|x|:x^n\}`$. Recall that a minimal subshift is called aperiodic if one (and then each) of its elements is not periodic.
###### Proposition 2.1.
Let $`(\mathrm{\Omega },T)`$ be a minimal subshift. Then the following are equivalent:
* $`(\mathrm{\Omega },T)`$ is aperiodic.
* $`m(n)\mathrm{}`$, $`n\mathrm{}`$.
Proof. (ii)$``$(i). This is clear.
(i)$``$ (ii). Assume the contrary. Then, there exsists an $`r>0`$ and sequences $`(u_k)`$, $`(v_k)`$ in $`𝒲`$ with
$$|u_k|r,k;|v_k|\mathrm{},k\mathrm{};u_k(v_k).$$
As $`u_k`$ is a return word of $`v_k`$, the word $`v_k`$ is a prefix of $`u_kv_k`$. Thus, $`v_k`$ begins with $`u_k^{l(k)}`$, where $`l(k)=\left[\frac{|v_k|}{|u_k|}\right]`$. Here, $`[x]`$ denotes the largest integer not exceeding $`x`$. By assumption, we have
(1)
$$l(k)\mathrm{},k\mathrm{}.$$
As $`\mathrm{\Omega }`$ is minimal, there exists an $`L>0`$ s.t. every word of length $`r`$ is a factor of every word of length $`L`$. By (1), the word $`u_k^L`$ belongs to $`𝒲`$ for sufficiently large $`k`$. Therefore, the word $`u_k^L`$ contains every word of length $`r`$ for sufficiently large $`k`$. But this implies
$$\mathrm{}\{v𝒲:|v|=r\}\mathrm{}\{v\text{Sub}(u_k^L):|v|=r\}|u_k|r$$
contradicting the aperiodicity of $`\mathrm{\Omega }`$. $`\mathrm{}`$
Consider now a minimal subshift $`(\mathrm{\Omega },T)`$. Let $`u,x𝒲`$ be given. Decomposing $`x`$ according to the occurrences of $`u`$ in $`x`$ gives a unique way of writing $`x`$ as $`x=au_1\mathrm{}u_lb`$ with
$$au_1ϵ,\mathrm{}_u(x)=l,u\text{ prefix of }u_j\mathrm{}u_lb,j=1,\mathrm{},l.$$
This decomposition is called the $`u`$-partition of $`x`$. It should be emphasized that in minimal $`𝒲`$ the word $`u`$ is a factor of every word of length bigger than $`R(u)`$ and therefore $`|a|,|b|R(u)`$ for all such $`w`$. Note that $`u_1,\mathrm{},u_{l1}`$ are return words of $`u`$ whereas $`u_l=u`$. Similarly, there is a unique way of writing an arbitrary $`\omega \mathrm{\Omega }`$ as $`\omega =\mathrm{}u_2u_1u_0u_1u_2\mathrm{}`$ with
* $`\omega (0)`$ belongs to $`u_0`$,
* every occurrence of $`u`$ in $`\omega `$ begins with one of the $`u_j`$, $`j`$.
If $`x`$ is a factor of $`y`$ then the $`u`$-partition of $`y`$ induces a decomposition of $`x`$ respecting $`u_1,\mathrm{},u_{l1}`$.
For a return word $`z`$ of $`u`$ and $`w𝒲`$ with $`u`$-partition $`w=au_1\mathrm{}u_lb`$, the topological number of occurrences of $`z`$ in $`w`$ is defined by
$$p_{z,u}(w)\mathrm{}\{j\{1,\mathrm{},l\}:u_j=z\}.$$
Apparently the following proposition holds.
###### Proposition 2.2.
$`p_{z,u}(w)=\mathrm{}_{zu}(w)`$.
A sequence $`(x_n)`$ in $`𝒲`$ is called partitioning sequence if $`|x_n|<|x_{n+1}|`$. This implies in particular $`|x_n|\mathrm{}`$, $`n\mathrm{}`$.
Given a partitioning sequence $`(x_n)`$, one can consider the series of $`x_n`$-partitions of an $`\omega \mathrm{\Omega }`$. This gives a hierarchic structure in $`\omega `$. This structure is a weak form of the hierarchic structures present in Sturmian models or substitutional models (cf. ).
## 3. Uniform additive ergodic theorems
The following is the key to the proof of Theorem 1
###### Lemma 3.1.
Let $`(\mathrm{\Omega },T)`$ be a strictly ergodic subshift over $`A`$ and let $`B`$ be a Banach space. Let $`F:𝒲B`$ be additive. Then the limit $`lim_{|w|\mathrm{}}\frac{F(w)}{|w|}`$ exists.
Proof. We consider two cases:
Case 1. $`(\mathrm{\Omega },T)`$ is aperiodic.
Let $`(x_n)`$ be an arbitrary partitioning sequence. As $`(\mathrm{\Omega },T)`$ is minimal, there exists only finitely many return words for each $`x_n`$, $`n`$. Therefore, the return words of $`x_n`$ can be denoted by $`y_1^{(n)},\mathrm{}y_{l(n)}^{(n)}`$ with a suitable $`l(n)`$. As $`(\mathrm{\Omega },T)`$ is strictly ergodic, the limits
(2)
$$\nu _j^{(n)}\underset{|w|\mathrm{}}{lim}\frac{p_{y_j^{(n)},x_n}(w)}{|w|}|y_j^{(n)}|=\underset{|w|\mathrm{}}{lim}\frac{\mathrm{}_{y_j^{(n)}x_n}(w)}{|w|}|y_j^{(n)}|$$
exists by Proposition 2.2 for all $`j,n`$ with $`jl(n)`$. Set
$$F^{(n)}\underset{j=1}{\overset{l(n)}{}}\nu _j^{(n)}\frac{F(y_j^{(n)})}{|y_j^{(n)}|}.$$
Of course, it suffices to show that for every $`\epsilon >0`$ there exists an $`n`$ with $`|\frac{F(w)}{|w|}F^{(n)}|\epsilon `$ for all $`w𝒲`$ long enough. So, choose $`\epsilon >0`$ arbitrary. By Proposition 2.1 and the additivity of $`F`$, there exists an $`n`$ with
(3)
$$c(m(|x_n|))\frac{\epsilon }{2}.$$
Considering now a $`w𝒲`$ with $`x_n`$-partition $`w=au_1\mathrm{}u_sb`$, we can estimate
$`\left|{\displaystyle \frac{F\left(w\right)}{\left|w\right|}}F^{\left(n\right)}\right|=\left|{\displaystyle \frac{F\left(w\right)}{\left|w\right|}}{\displaystyle \underset{j=1}{\overset{l\left(n\right)}{}}}\nu _j^{\left(n\right)}{\displaystyle \frac{F\left(y_j^{\left(n\right)}\right)}{\left|y_j^{\left(n\right)}\right|}}\right|`$
$``$ $`\left|{\displaystyle \frac{F\left(w\right)}{\left|w\right|}}{\displaystyle \frac{F\left(a\right)+_{j=1}^sF\left(u_j\right)+F\left(b\right)}{\left|w\right|}}\right|+\left|{\displaystyle \frac{F\left(a\right)+_{j=1}^sF\left(u_j\right)+F\left(b\right)}{\left|w\right|}}{\displaystyle \underset{j=1}{\overset{l\left(n\right)}{}}}\nu _j^{\left(n\right)}{\displaystyle \frac{F\left(y_j^{\left(n\right)}\right)}{\left|y_j^{\left(n\right)}\right|}}\right|`$
$``$ $`{\displaystyle \frac{c\left(\left|a\right|\right)\left|a\right|+c\left(\left|b\right|\right)\left|b\right|}{\left|w\right|}}+{\displaystyle \underset{j=1}{\overset{s}{}}}{\displaystyle \frac{c\left(\left|u_j\right|\right)}{\left|w\right|}}\left|u_j\right|+{\displaystyle \frac{D\left|a\right|+D\left|b\right|}{\left|w\right|}}`$
$`+`$ $`{\displaystyle \underset{j=1}{\overset{l\left(n\right)}{}}}\left|{\displaystyle \frac{p_{y_j^{\left(n\right)},x_n}\left(w\right)\left|y_j^{\left(n\right)}\right|}{\left|w\right|}}\nu _j^{\left(n\right)}\right|{\displaystyle \frac{\left|F\left(y_j^{\left(n\right)}\right)\right|}{\left|y_j^{\left(n\right)}\right|}}`$
$``$ $`{\displaystyle \frac{\left(c\left(0\right)+c\left(0\right)\right)R\left(x_n\right)}{\left|w\right|}}+{\displaystyle \frac{2DR\left(x_n\right)}{\left|w\right|}}+{\displaystyle \frac{\epsilon }{2}}+D{\displaystyle \underset{j=1}{\overset{l\left(n\right)}{}}}\left|{\displaystyle \frac{p_{y_j^{\left(n\right)},x_n}\left(w\right)\left|y_j^{\left(n\right)}\right|}{\left|w\right|}}\nu _j^{\left(n\right)}\right|,`$
where we used (3) in the last equation. By (2), this gives $`\left|\frac{F(w)}{|w|}F^{(n)}\right|\epsilon `$ for all $`w`$ which are sufficiently long and the proof of Case 1 is finished.
Case 2. $`(\mathrm{\Omega },T)`$ is periodic.
There exists a $`p`$ with $`T^p\omega =\omega `$ for every $`\omega \mathrm{\Omega }`$. Fix a word $`v`$ in $`𝒲`$ of length $`p`$. Let $`n`$ be arbitrary. Then every $`w𝒲`$ with $`|w|3|v^n|`$ can be written as $`w=av_1\mathrm{}v_sb`$ with $`v_j=v^n`$, $`j=1,\mathrm{},s,`$ and $`|a|,|b||v^n|`$. Using this partition of $`w`$ instead of the $`x_n`$-partition, one can mimick the proof of Case 1. This gives the existence of the limit in question. $`\mathrm{}`$
###### Lemma 3.2.
Let $`\mathrm{\Omega },𝒲`$ and $`B`$ be as in the foregoing lemma. Let $`f:\mathrm{\Omega }B`$ be a continuous function. Then, $`\frac{1}{n}_{j=1}^nf(T^j\omega )`$ converges uniformly in $`\omega `$ towards a constant function.
Proof. To $`\omega \mathrm{\Omega }`$, we can associate the function $`F_\omega :𝒲B`$ defined by
$$F_\omega (w)\underset{j=0}{\overset{|w|1}{}}f(T^{k_\omega ^w+j}\omega ),$$
where $`k_\omega ^w`$ is minimal with $`\omega _{k_\omega ^w}\mathrm{}\omega _{k_\omega ^w+|w|1}=w`$.
Claim. There exists a nonincreasing function $`c:[0,\mathrm{})[0,\mathrm{})`$ with $`lim_r\mathrm{}c(r)=0`$ s.t.
$$|F_\omega (w)F_\rho (w)|c(|w|)|w|.$$
Proof of the claim. It is clearly enough to show that, for every $`\epsilon >0`$, there exists an $`l=l(\epsilon )`$ with $`|F_\omega (w)F_\rho (w)|\epsilon |w|`$ for all $`\omega ,\rho \mathrm{\Omega }`$ and every $`w𝒲`$ with $`|w|l`$.
As $`f`$ is continuous on the compact space $`\mathrm{\Omega }`$, there exists a $`k`$ s.t. $`\omega _k\mathrm{}\omega _k=\rho _k\mathrm{}\rho _k`$ implies
(4)
$$|f(\omega )f(\rho )|\frac{\epsilon }{2}.$$
Set $`l\mathrm{max}\{\frac{8kf_{\mathrm{}}}{\epsilon },2k\}`$, where $`_{\mathrm{}}`$ denotes the supremum-norm. For $`w𝒲`$ with $`|w|l`$, we can estimate the value of $`D|F_\omega (w)F_\rho (w)|`$ by
$`|{\displaystyle \underset{j=0}{\overset{\left|w\right|1}{}}}(f\left(T^{k_\omega ^w+j}\omega \right)f\left(T^{k_\rho ^w+j}\rho \right)|`$ $``$ $`4f_{\mathrm{}}k+{\displaystyle \underset{j=k}{\overset{\left|w\right|k1}{}}}|(f\left(T^{k_\omega ^w+j}\omega \right)f\left(T^{k_\rho ^w+j}\rho \right)|`$
$`(\text{4})4f_{\mathrm{}}k+\left(\left|w\right|2k\right){\displaystyle \frac{\epsilon }{2}}`$ $`=`$ $`\left|w\right|\left({\displaystyle \frac{4f_{\mathrm{}}k}{\left|w\right|}}+{\displaystyle \frac{\left(\left|w\right|2k\right)}{\left|w\right|}}{\displaystyle \frac{\epsilon }{2}}\right)\left|w\right|\epsilon .`$
The proof of the claim is finished.
The claim implies that
* for every $`\omega \mathrm{\Omega }`$ the function $`F_\omega `$ is additive,
* $`\frac{F_\omega (w)F_\rho (w)}{|w|}c(|w|)`$, with $`lim_{|w|\mathrm{}}c(|w|)=0`$.
Using this and Lemma 3.1, it is straightforward to finish the proof of the lemma. $`\mathrm{}`$
Now, the proof of Theorem 1 can easily be accomplished.
Proof of Theorem 1. The implication (i) $``$ (ii) follows immediately from Lemma 3.1. Similarly, (ii) $``$ (iii) follows immediately from Lemma 3.2. Finally, (iii)$``$ (i) is clear as the function $`w\mathrm{}_v(w)`$ is additive for each $`v𝒲`$. The proof of the theorem is finished. $`\mathrm{}`$
## 4. Uniform subadditive ergodic theorems
Here, we consider subadditive functions on $`𝒲`$. The key results are the following two lemmas. Their proofs use and considerably extend ideas from .
###### Lemma 4.1.
Let $`(\mathrm{\Omega },T)`$ be a minimal subshift over $`A`$ satisfying (PQ). Then, the subshift $`(\mathrm{\Omega },T)`$ satisfies (SET) as well.
Proof. Set $`\overline{F}inf_nF^{(n)}`$. We will show
$$(A)\underset{|w|\mathrm{}}{lim\; sup}\frac{F(w)}{|w|}\overline{F}\text{ and}(B)\underset{|w|\mathrm{}}{lim\; inf}\frac{F(w)}{|w|}\overline{F}.$$
Ad (A): It is clearly enough to show $`lim\; sup_{|w|\mathrm{}}\frac{F(w)}{|w|}F^{(n)}`$ for every $`n`$. But this follows easily form the subadditivity of $`F`$ by chopping each $`w`$ into parts of length $`n`$ and a boundary word and using that the boundary terms tend to zero (cf. ).
Ad (B): Assume the contrary. This implies, in particular, $`\overline{F}>\mathrm{}`$ and that there exists a sequence $`(v_n)`$ in $`𝒲`$ as well as a $`\delta >0`$ with $`|v_n|`$ tending to $`\mathrm{}`$ for $`n\mathrm{}`$ and
(5)
$$\frac{F(v_n)}{|v_n|}\overline{F}\delta $$
for every $`n`$. Moreover, by (A), there exists an $`L_0`$ with
(6)
$$\frac{F(w)}{|w|}\overline{F}+\frac{C\delta }{8}$$
for all $`w𝒲`$, $`|w|L_0`$, where $`C`$ is the constant from (PQ).
Fix $`m`$ with $`|v_m|L_0`$. Using (PQ), we can now find an $`L_1`$ s.t. every $`w𝒲`$ with $`|w|L_1`$ can be written as $`w=x_1v_mx_2v_m\mathrm{}x_lv_mx_{l+1}`$ with
(7)
$$\frac{l2}{2}\frac{C}{4}\frac{|w|}{|v_m|}.$$
Now, considering only every other copy of $`v_m`$ in $`w`$, we can write $`w`$ as $`w=y_1v_my_2\mathrm{}y_rv_my_{r+1}`$, with $`|y_j||v_m|L_0`$, $`j=1,\mathrm{},r+1,`$ and by (7)
$$r\frac{l2}{2}\frac{C}{4}\frac{|w|}{|v_m|}.$$
Using $`(\text{5})`$, $`(\text{6})`$ and this estimate, we can now calculate
$`{\displaystyle \frac{F(w)}{|w|}}`$ $``$ $`{\displaystyle \underset{j=1}{\overset{r+1}{}}}{\displaystyle \frac{F(y_j)}{|y_j|}}{\displaystyle \frac{|y_j|}{|w|}}+{\displaystyle \frac{F(v_m)}{|v_m|}}{\displaystyle \frac{r|v_m|}{|w|}}{\displaystyle \underset{j=1}{\overset{r+1}{}}}(\overline{F}+{\displaystyle \frac{C}{8}}\delta ){\displaystyle \frac{|y_j|}{|w|}}+(\overline{F}\delta ){\displaystyle \frac{r|v_m|}{|w|}}`$
$``$ $`\overline{F}+{\displaystyle \frac{C}{8}}\delta {\displaystyle \frac{C}{4}}{\displaystyle \frac{|w|}{|v_m|}}{\displaystyle \frac{|v_m|}{|w|}}\delta \overline{F}{\displaystyle \frac{C}{8}}\delta .`$
As this holds for arbitrary $`w𝒲`$ with $`|w|=L_1`$, we arrive at the obvious contradiction $`F^{(L_1)}\overline{F}\frac{C}{8}\delta <inf_nF^{(n)}`$. This finishes the proof. $`\mathrm{}`$
###### Lemma 4.2.
Let $`(\mathrm{\Omega },T)`$ be a minimal subshift over $`A`$ satisfying (SET). Then, the subshift $`(\mathrm{\Omega },T)`$ satisfies (PQ) as well.
Proof. Note that, for $`v𝒲`$, the function $`(l_v)`$ is subadditive. Thus, the equation
(8)
$$\nu (v)\underset{|w|\mathrm{}}{lim\; inf}\frac{l_v(w)}{|w|}=\underset{|w|\mathrm{}}{lim}\frac{l_v(w)}{|w|}$$
holds by (SET). The proof will now be given by contraposition. So, let us assume that the values $`\nu (v)`$, $`v𝒲`$, are not bounded away from zero. As the system is minimal, we have $`\nu (w)>0`$ for every $`w𝒲`$. Thus, there exists a sequence $`(v_n)`$ in $`𝒲`$ with
(9)
$$\nu (v_n)>0,\text{and}\nu (v_n)0,n\mathrm{}.$$
As the alphabet $`A`$ is finite, there are only finitely many words of a prescribed length. Thus, (9) implies
(10)
$$|v_n|\mathrm{},n\mathrm{}.$$
Replacing $`(v_n)`$ by a suitable subsequence, we can assume by (9) that the equation
(11)
$$\underset{n=1}{\overset{\mathrm{}}{}}\nu (v_n)<\frac{1}{2}$$
holds. By $`(\text{8})`$, $`(\text{10})`$ and $`(\text{11})`$, we can choose inductively for each $`k`$ a number $`n(k)`$, with
(12)
$$\underset{j=1}{\overset{k}{}}\frac{l_{n(j)}(w)}{|w|}<\frac{1}{2}$$
for every $`w𝒲`$ with $`|w|\frac{|v_{n(k+1)}|}{2}`$. Note that $`(\text{12})`$ implies
(13)
$$|v_{n(k)}|<\frac{|v_{n(k)+1}|}{2}$$
as $`\frac{l_{n(k)}(v_{n(k)})}{|v_{n(k)}|}=1`$. Define the function $`l:𝒲`$ by
$$l(w)\underset{j=1}{\overset{\mathrm{}}{}}l_{n(j)}(w).$$
Note that the sum is actually finite for each $`w𝒲`$. Obviously, $`(l)`$ is subadditive. Thus, by assumption, the limit $`lim_{|w|\mathrm{}}\frac{l(w)}{|w|}`$ exists. On the other hand, we clearly have
$$\frac{l(v_{n(k)})}{|v_{n(k)}|}\frac{l_{n(k)}(v_{n(k)})}{|v_{n(k)}|}1$$
as well as by the induction construction $`(\text{12}),(\text{13})`$
$$\frac{l(w)}{|w|}=\underset{j=1}{\overset{k}{}}\frac{l_{n(j)}(w)}{|w|}<\frac{1}{2}$$
for $`w𝒲`$ with $`\frac{|v_{n(k+1)}|}{2}|w|<|v_{n(k+1)}|`$. This gives a contradiction proving the lemma. $`\mathrm{}`$
Proof of Theorem 2. This theorem follows immediately from the foregoing two lemmas. $`\mathrm{}`$
## 5. Examples
In this section we discuss two classes of examples satisfying condition (PQ). We will need the following proposition proved similarly to Proposition 2.1.
###### Proposition 5.1.
For $`(\mathrm{\Omega },T)`$ the following are equivalent:
(i) $`\mathrm{\Omega }`$ satisfies (HP).
(ii) There exists a $`\kappa >0`$, s.t. the length of every return word of $`v`$ is at least $`\frac{|v|}{\kappa }`$.
Proof. (i) $``$ (ii). Let $`u`$ be a return word to $`v`$. Then $`v`$ is a prefix of $`uv`$. Thus, $`v`$ starts with $`u^l`$, where $`l=[\frac{|v|}{|u|}]`$. By (HP) this implies $`[\frac{|v|}{|u|}]N`$ which in turn yields $`\frac{|v|}{|u|}1N`$. Now, (ii) follows easily.
(ii) $``$ (i). Let $`v𝒲`$ be primitive and assume that $`v^l`$ belongs to $`𝒲`$ as well for a suitable $`l`$, $`l2`$. Then $`v`$ is a return word to $`v^{l1}`$ implying $`|v|\frac{|v^{l1}|}{\kappa }=\frac{l1}{\kappa }|v|`$. This gives immediately $`\kappa l1`$ and the proof of the proposition is finished. $`\mathrm{}`$
From this proposition, we can derive a sufficient condition for (PQ).
###### Proposition 5.2.
If $`(\mathrm{\Omega },T)`$ satisfies (HP) and (PW), then it satisfies (PQ) as well and (SET) holds.
Proof. Choose an arbitrary $`v𝒲`$. By (PW), a word $`w𝒲`$ of sufficient length contains at least $`\frac{E}{2}\frac{|w|}{|v|}`$ copies of $`v`$, where $`E`$ is the constant from (PW). By (HP) and Proposition 5.1 (ii), this implies that a word $`w`$ of sufficient length contains at least $`\frac{E}{2}\frac{|w|}{|v|(\kappa +1)}`$ disjoint copies of $`v`$. This gives $`\nu (v)\frac{E}{2(\kappa +1)}>0`$. Thus, $`(\mathrm{\Omega },T)`$ satisfies (PQ). Moreover, (PQ) implies minimality of $`(\mathrm{\Omega },T)`$ (cf. Remark 1. (c)). Now, the proposition follows from Theorem 2. $`\mathrm{}`$.
###### Corollary 5.3.
If $`(\mathrm{\Omega },T)`$ is linearly repetitive, then a subaddditive ergodic theorem holds.
Proof. It is well known (and easy to see) that a linearly repetitive $`(\mathrm{\Omega },T)`$ satisfies (HP) and (PW). Thus, the corollary follows from the foregoing proposition. $`\mathrm{}`$
###### Remark 2.
(a) In fact, as shown in , linear repetitivity implies a subadditive ergodic theorem in tiling dynamical systems of arbitrary dimension (cf. as well).
(b) The relationship between linear repetitivity and (HP) +(PW) is currently under investigation in .
Acknowledgements. The author would like to thank David Damanik for many stimulating discussions. In fact, this paper would not have been possible without the collaboration leading to .
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# Photometric study of fields of nearby pulsars with the 6 m telescope
## 1 Introduction
In a vast list of more than 1000 radio pulsars there are but a few objects detected in other wavelengths ($`\gamma `$rays, X-rays, far-UV, optical, see e.g., Ulmer, 1998; Becker & Trümper, 1999; Korpela E.J. & Bowyer S., 1998; Mignani, 1998). Multiwavelength observations bring a wealth of a new information on the radiation mechanisms of pulsars which cannot be obtained in a narrow spectral band. For young ($`\begin{array}{c}<\\ \end{array}10^4`$ yrs) pulsars like those in the Crab and Vela nebula the emission in the all spectral ranges is mainly of non-thermal origin; it is produced by relativistic particles generated in magnetospheres of rapidly rotating neutron stars (NSs). Becoming older ($`\begin{array}{c}>\\ \end{array}10^5`$ yrs), pulsars rotate slower, the non-thermal component weakens and one can observe thermal emission from the surface of cooling NSs. According to standard NS cooling models (e.g., Nomoto & Tsuruta, 1987), at this age their surface temperature is about $`10^5`$$`10^6`$K and the maximum of the thermal emission lies in soft X-rays (0.1–2.4 keV) or in far-UV. Thermal emission with spectral temperatures in the above range has been detected at X-ray observations of some middle-aged radio pulsars (see, e.g., Becker & Trümper, 1997). Thermal X-ray emission was also detected from several radio silent objects identified as isolated neutron stars (INSs) (e.g., Neuhäuser & Trümper 1999).
Simulations of INS cooling show that at certain conditions, depending on the NS mass/radius ratio, on equation of state and composition of superdense matter in interiors and in surface layers of the star, (see, e.g., Yakovlev et al. 1999 for a recent review), thermal evolution of a NS may strongly deviate from the standard model. Theoretical investigations combined with comprehensive observational studies of thermal emission from radio pulsars and INSs of different ages enable one to understand which of the evolution scenarios are real. These studies are also of crucial importance for development of realistic models of NSs and deeper understanding of poorly known properties of superdense matter in their interiors.
Optical observations are an important part of the multiwavelength studies of INSs and pulsars. They allow one to constrain the parameters of the thermal emission in the Rayleigh-Jeans spectral region and investigate the properties of the nonthermal radiation in optical bands. Most pulsars and INSs, excepting the young Crab-pulsar and PSR 0540-69, are faint optical objects. Thus, the multicolour photometry is the natural first step to search for the pulsar optical counterparts and to obtain information on their optical spectra. This is the main goal of the program which is carried out at the 6 m telescope during last several years.
The program includes a deep (up to 27<sup>m</sup>) photometry of the fields of some nearby pulsars. The most interesting results reported at present are those obtained for the middle-aged ($`10^5`$ yrs) pulsar PSR B0656+14 (Kurt et al., 1998; Koptsevich et al., 2000). Deep multicolour photometry observations of the pulsar field taken with the 6 m telescope yielded the detection of the pulsar optical counterpart. The pulsar magnitudes in BR<sub>c</sub>I<sub>c</sub>-bands were estimated for the first time. The correspondent fluxes exceed the Rayleigh-Jeans extrapolation of the thermal spectrum seen in the soft X-rays and EUV.
In this paper we report photometry of four other nearby pulsars: Geminga (J0633+1746), PSR B0950+08, PSR J1908+0734 and PSR J0108-1431.
Their parameters are presented in Table 1. The data were taken from the pulsar catalogue (Taylor et al., 1995), excepting Geminga’s DM (Malofeev & Malov, 1997).
Observations and data reduction are described in Section 2, Section 3 summarizes the results, and some conclusions are given in Section 4.
## 2 Observations and data reduction
Optical observations of the pulsar fields were carried out with the 6 m telescope BTA on March, 1997, January, 1998, on January, July and August, 1999 with a CCD detector mounted at the prime focus with the set of filters close to the Johnson-Cousins system. Different CCD detectors were used. Some of their characteristics are given in Table 2. Seeing varied from night to night between $`0.^{\prime \prime }9`$ and $`1.^{\prime \prime }8`$. Table 3 gives details on each observing run.
Standard data reduction, including bias subtraction, account for dark current, correction for non-uniformity of the detector sensitivity (flat-fielding), and removing of cosmic ray events, was performed making use of MIDAS procedures. To compensate fringes we used the so-called superflats, i.e. flats obtained directly from the sky by median-combining several science exposures taken during the night. The individual flat-fielded exposures were stacked together, yielding combined images of the fields under investigation.
For the photometric calibration we used Landolt’s standards (1992). The absolute fluxes $`F_j`$ in \[erg cm<sup>-2</sup>s<sup>-1</sup>Hz<sup>-1</sup>\] were calculated using equations
$$\mathrm{log}F_j=\left(0.4M_j+M_j^0\right),$$
(1)
with the zero-points provided by Fukugita et al.(1995):
| $`M_B^0=19.396,M_V^0=19.445,`$ |
| --- |
| $`M_R^0=19.520,M_I^0=19.623.`$ |
(2)
Astrometric referencing of the images was performed using the coordinates of selected field stars extracted from the USNO catalogue using the ESO Skycat tool. The final accuracy was about $`1^{\prime \prime }`$.
## 3 Results
### 3.1 PSR J0633+1746
The middle-aged Geminga pulsar has been already detected with the ground-based (CFHT, ESO 3.6 m, NTT) and HST telescopes in different bands and presently it is one of the well-studied pulsars in the the optical domain (see, e.g., Bignami & Caraveo, 1996, Mignani et al., 1998). It has been concluded that the optical emission of Geminga seems to be thermal with a broad emission feature at 6000Å. Jacchia et al.(1999) suggested a rough model to explain the observed excess of Geminga’s emission in the V-band as emission of hot ions in the strong magnetic field near the NS surface. The results of Geminga’s spectroscopy (Martin et al, 1998) with tentative detection of an apparent absorption feature at 6400Å are generally consistent with the photometry. In Figure 1 we present our images of the Geminga field in the B, V, R<sub>c</sub> bands taken on March, 1997 and in the I<sub>c</sub> band on January 19, 1999. The pulsar position in the V and R images is indicated with arrow. The objects suggested to be Geminga candidates at the first observations of this field in the optical range (Halpern & Tytler, 1988) are marked with letters G and G.
The photometry both of the pulsar optical counterpart (G<sup>′′</sup>) and the mentioned field objects yields the magnitude estimates, which are presented in Table 4. The obtained Geminga flux distribution is in a good agreement with the data published so far. In Figure 2 we present the broadband spectrum of Geminga based on all optical data available. Filled circles correspond to our results, and arrows with bars to the I and F675W upper limits. Our observations of the Geminga field in the I-band were tentative and we could not get so deep images as had been done before (Bignami et al., 1996). Thus, the fading of Geminga’s flux at these wavelengths is still unconfirmed.
In all the optical bands (excluding the I one) the optical fluxes exceed the Rayleigh-Jeans extrapolation of the thermal spectrum seen in the EUV and soft X-rays (dashed line in Figure 2). Such a spectrum behaviour is similar to that of the PSR B0656+14 and probably can be also fitted by combination of both a magnetospheric and a thermal component.
#### 3.1.1 PSR B0950+08
Like Geminga, the pulsar B0950+08 has been previously detected. The field of this pulsar was observed with the HST in the UV-optical range using the long-pass filter F130LP ($`\lambda \lambda =23104530`$ÅÅ) (Pavlov et al., 1996). The only pointlike source was detected in the image of $`7.^{\prime \prime }4\times 7.^{\prime \prime }4`$ in size ($`m_{F130LP}`$ = 27.1; F = $`0.051\pm 0.003\mu Jy`$). Its location does not coincide with the radio pulsar position. The difference in positions is about $`1.^{\prime \prime }85`$, that may be caused by systematic uncertainties in the HST Guide Star Catalogue (up to $`1^{\prime \prime }`$), which is used for HST astrometric referencing, as well as by uncertainties in the pulsar radio position ($`0.^{\prime \prime }5`$).
The field of this pulsar was observed with the 6 m telescope in the R band of the Cousins system on March, 1997 and January, 1998.
There are not enough reference stars from USNO catalogue in this field (approximately $`2.^{}5\times 2.^{}5`$ in size), that makes an accurate astrometry impossible. To improve the situation we used the image of this field ($`9{}_{}{}^{}\times 9^{}`$) taken with the Zeiss-600 telescope of SAO RAS. This enabled us to increase the accuracy of the astrometrical referencing of the BTA images up to $`1^{\prime \prime }`$. In Figure 3 we present an image of PSR B0950+08 neighbourhood, positions of the radio pulsar (small cross), the UV-opical candidate (big cross) and the possible optical counterpart (in circle) are indicated. An object ($`S/N=4.2`$) has been found in the combined image, and its position differs by $`1.^{\prime \prime }5`$ from the pulsar radio position. In case the detected object is indeed the pulsar, this discrepancy may be due to two reasons: different epochs of the radio and optical observations with the lack of reliable data on the pulsar proper motion, and uncertainties both in our astrometry and the pulsar radio position. The object magnitude is $`R=25.^\mathrm{m}5(3)`$, corresponding flux is $`F_R=0.19\pm 0.04\mu Jy`$.
Figure 4 shows photometric data of the possible PSR B0950+08 counterpart based on the 6 m telescope (R-band) and HST (F130LP) observations.
If the observed object is the pulsar optical counterpart its R-magnitude shows that old pulsars may be brighter in near-IR than in near-UV region. This is in opposite to the expected spectral behaviour and cannot be explained by thermal emission only. More observations are needed.
#### 3.1.2 PSR J0108-1431
PSR J0108-1431 was discovered during a survey of the southern sky for pulsars using the Parkes 64 m radio telescope (Tauris et al., 1994). In accordance with the galactic electron distribution density model (Taylor & Cordes, 1993) the low dispersion measure suggests that this pulsar is within 90 pc from the Sun. Thus this is the closest known neutron star, but it is quite old, $`\tau 1.610^8`$yrs. The pulsar has not been detected in any high-frequency range.
Figure 5 represents images of the PSR J0108-1431 neighbourhood taken with the 6 m telescope on August 1999. No object coinciding with the radio pulsar in position was detected and we obtained the following 3$`\sigma `$-limits for the observed pulsar magnitudes: B $`>25.^\mathrm{m}4`$; V $`>24.^\mathrm{m}7`$; R<sub>c</sub> $`>25.^\mathrm{m}4`$; I<sub>c</sub> $`>24.^\mathrm{m}3`$. Further observations of this field might yield the detection of the INS optical counterpart.
#### 3.1.3 PSR J1908+0746
The pulsar PSR J1908+0746 was discovered at the Arecibo Observatory during the search for pulsars of low luminosity using the 305 m radio telescope (Camilo & Nice, 1995). For its characteristic age this is old pulsar, while the rotational energy loss $`\dot{E}`$ is rather high. In the absence of an unequivocal theory of the high-energy emission of pulsars, a high level of $`\dot{E}/d^2`$ is used as a rough indicator of likelihood that high-frequency emission from a pulsar can be detected (see, e.g., Goldoni & Musso, 1996). Search for the emission in the high-frequency domain from pulsars with the highest values of $`\dot{E}/d^2`$ ($``$40 pulsars including PSR J1908+0746) yielded the detection of X-rays from 27 pulsars (Becker & Trumper, 1997) and the $`\gamma `$-rays from 4 pulsars (Thompson et al., 1994). No high-frequency emission has been found so far from PSR J1908+0746 (Becker & Trumper, 1997).
The field of this pulsar was observed in BVR<sub>c</sub>I<sub>c</sub> bands on July and August, 1999. The obtained images are shown in Figure 6. The photometry yields the following 3$`\sigma `$ upper limits for the observed pulsar magnitudes: B $`>26.^\mathrm{m}0`$, V $`>26.^\mathrm{m}1`$, R<sub>c</sub> $`>25.^\mathrm{m}9`$, I<sub>c</sub> $`>23.^\mathrm{m}4`$. As this rather distant pulsar lies almost in the galactic plane, its dereddened flux values can be derived only after an accurate study of the interstellar extinction towards the pulsar position.
## 4 Discussion and conclusions
The deep photometric study of the PSR B0656+14 at the 6 m telescope has shown that the broadband spectrum of this middle-aged pulsar is significantly of nonthermal origin (Kurt et al., 1998; Koptsevich et al., 2000). Its complicated shape differs from flat and featureless spectra of young Crab-like pulsars and cannot be explained by a simple spectral model (Pavlov et al., 1997; Kurt et al., 1998; Koptsevich et al., 2000). Similar behaviour, confirmed by our observations, shows the spectrum of the slightly older pulsar Geminga. Some doubt still remains how deep is the fall of Geminga’s flux redward of R and whether its depth is restricted by the emission level of the thermal component. More observations, including the IR-range, are needed to address this question. Nevertheless, current stage of the optical studies allows us to suggest that the optical spectra of these middle-aged NSs are likely to be very different from those of young ones and may hint the spectral evolution of the optical emission with the pulsar age.
The detection of the optical emission from the old pulsar PSR B0950+08, if proposed counterpart candidate is confirmed by further observations, may contribute to the above idea. The apparent excess in the R-band may indicate the presence of spectral feature at these wavelengths and raise a question whether thermal component is dominant in the optical emission of old neutron stars ($`>10^6`$ yrs), which is expected from the scenario of INSs evolution. Thus, search for optical counterparts to old pulsars and cooling INSs detected in X-rays may be of great importance.
We have observed the fields of the pulsars J0108-1431 and J1908+0734 for the first time. Although no optical counterparts have been found, the estimated magnitude upper limits suggest deeper observations of the fields of these promising objects. However, the optical emission of PSR J1908+0734 seems unlikely to be detected in the B and V bands due to extinction.
The search for optical emission of more INSs with the subsequent multicolour photometric study is needed to shed light on the mechanisms of their optical emission, to put constraints to the parameters of both thermal and nonthermal components, to enrich our knowledge of the NS interior and draw a real picture of the neutron star evolution.
* The work on pulsar study was partially supported by grant 1.2.6.4 of Program ”Astronomia”, by INTAS (grant 96–0542) and RFBR (grant 99-02-18099). The authors are grateful to S.V. Zharikov, N.A. Tikhonov, A.V. Moiseev and V.R. Amirkhanyan for assistance in observations. VNK thanks A.I. Kopylov, N.A. Tikhonov and I.O. Drozdovsky for fruitful discussions and valuable comments.
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# Regular Coverings in Filter and Ideal Lattices
## Introduction
The purpose of this paper is to further develop and apply the theory of regular coverings in a complete modular lattice, introduced in . The point of that theory is to generalize, to complete modular lattices, some of the nice results available for finite-height modular lattices.
For example, given a finite-height module $`M`$ over a ring $`R`$ (i.e., a module having a finite-height lattice of submodules), a *composition series* of $`M`$ is a (necessarily finite) sequence of submodules
$$\{\mathrm{\hspace{0.17em}0}\}=M_0M_1\mathrm{}M_n=M$$
such that each quotient $`M_i/M_{i1}`$ is simple. The Jordan-Holder Theorem states that any two composition series $`\{M_i\}_{i=1}^n`$, $`\{M_i^{}\}_{i=1}^n^{}`$ are the same length $`n=n^{}`$ and the quotients $`M_i/M_{i1}`$ can be paired with the quotients $`M_j^{}/M_{j1}^{}`$ in such a way that corresponding quotients are isomorphic.
The Jordan-Holder Theorem is an algebraic version of the lattice-theoretic Dedekind-Birkhoff Theorem. The lattice-theoretic correlates of the composition series and the isomorphism of corresponding quotients are maximal chains (maximal linearly-ordered subsets) and projective equivalence of coverings. The Dedekind-Birkhoff Theorem states that in any two maximal chains in a finite-height modular lattice, the lengths of the chains are the same and coverings in the chains can be paired in such a way that corresponding coverings are projectively equivalent.
Unfortunately, the Dedekind-Birkhoff Theorem can fail for infinite-height modular lattices. For example, consider the modular lattice
where $`_ix_i=_iy_i=`$. There are two maximal chains
$$C_1:<\mathrm{}<x_n<\mathrm{}<x_2<x_1<y_1$$
and
$$C_2:<\mathrm{}<y_n<\mathrm{}<y_2<y_1$$
such that the covering $`x_1y_1`$ appears in $`C_1`$, but there is no projectively equivalent covering in $`C_2`$.
The theory of regular coverings was created to try to remedy this situation. In the language of that theory, the covering $`x_1y_1`$ is not *regular*, as defined in Section 2. If it were regular, such behavior would be impossible because of Theorem 4.
Now, note that if we embed the lattice $`L`$ into its lattice of filters $`\mathrm{Fil}L`$, we obtain the lattice
which contains new elements, $`\mathrm{Fg}\{x_1,x_2,\mathrm{}\}`$ and $`\mathrm{Fg}\{y_1,y_2,\mathrm{}\}`$, forming a covering equivalent to $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$. Now any two maximal chains in $`\mathrm{Fil}L`$ have one covering equivalent to $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$. We have *regularized* the covering $`x_1y_1`$ by embedding $`L`$ into $`\mathrm{Fil}L`$. In this paper, we explore this process and strategy of regularization further. Of course, this involves studying coverings in $`\mathrm{Fil}L`$ (and in its dual, the lattice $`\mathrm{Idl}L`$ of *ideals* of $`L`$) and trying to determine whether or not they are regular.
We also examine questions of multiplicity, since in a modular lattice that is not distributive, a maximal chain can contain more than one covering projectively equivalent to a given one.
We give two applications. One application is a generalization of the theory of chief factors of an algebra having a modular congruence lattice. The information supplied by these results is entirely lattice-theoretic; we leave for another time the algebraic correlates such as play roles in the Jordan-Holder Theorem. The other application is a way of defining the steps in the proof of a theorem. Any proof of the theorem from the same premises must *cover*, as we say, the same steps. Also, from any set of instances of rules of inference which covers the steps, a finite subset can be selected and used to construct a proof.
After this introduction and a section of preliminaries, this paper begins in Section 1 with some definitions relating to multiplicities. Given a maximal chain $`C`$ in the modular lattice $`L`$, and a covering $`xy`$, there is a corresponding multiplicity of $`xy`$ in $`C`$ which may vary with $`C`$, except in the important case when $`xy`$ is *weakly regular*. We also define notations for upper and lower bounds on the multiplicity.
Section 2 discusses the theory of *regular coverings*, which, due to a generalization of the Dedekind-Birkhoff Theorem as given in , are also weakly regular.
Section 3 is a preliminary examination of coverings in filter and ideal lattices. As our strategy is to use regular coverings in such lattices for various purposes, we must understand their basic properties before attempting to determine whether or not they are regular. In this section, among other things, we classify filter and ideal coverings into three categories: *atomic*, *quasi-atomic*, and *anomalous* coverings.
Section 4 gives proofs of the stability of regularity, multiplicity when regular, and in some cases the multiplicity upper bound, under the embedding from $`L`$ into $`\mathrm{Fil}L`$ or $`\mathrm{Idl}L`$.
Section 5 proves a relationship between the multiplicity upper bound of a covering $`xy`$ in $`L`$ and the multiplicity lower bound of the corresponding filter or ideal covering. The important consequence of this is that under appropriate conditions, if the multiplicity bound is infinite, then any maximal chain in the filter or ideal lattice will have an infinite number of coverings equivalent to $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$ or $`\mathrm{Ig}\{x\}\mathrm{Ig}\{y\}`$. This complements other theorems which describe the behavior when the multiplicity upper bound is finite.
Section 6 discusses upper regularity of filter coverings (and dually, lower regularity of ideal coverings). We show in this section that anomalous filter and ideal coverings cannot be regular. We also give an example of an atomic filter covering in an algebraic lattice that is not upper regular, and thus not regular.
Section 7 gives a proof that certain filter and ideal coverings are regular. In particular, we show that in a meet-continuous lattice, if the multiplicity upper bound of a covering $`xy`$ if finite, then the corresponding filter covering is regular.
Section 8 discusses the application of these ideas to generalizing the Jordan-Holder Theorem.
Section 9 applies the theory to defining the steps in the proof of a proposition from given premises.
We will talk almost entirely about modular lattices, complete in most cases, except in Section 9, where we will talk about the distributive lattice underlying a boolean algebra $`𝐁`$, and complete distributive lattices constructed from it.
## 0. Preliminaries
The reader should know about *modular lattices* and *distributive lattices*, and that distributive lattices are modular. The reader should also know about *complete lattices*.
We denote the least element of any lattice, if one exists, by $``$, and the greatest element by $``$. If $`xy`$ are elements of a lattice $`L`$, then we denote by $`\mathrm{I}_L[x,y]`$, or simply $`\mathrm{I}[x,y]`$, the *interval sublattice* of elements $`z`$ such that $`xzy`$.
A *covering* is a pair $`x,y`$ of elements such that $`x<y`$ and $`\mathrm{I}[x,y]`$ has only $`x`$ and $`y`$ as elements. We say that $`x`$ is *covered* by $`y`$, or $`xy`$. We will often say $`xy`$ not only to state that $`x`$ is covered by $`y`$, but also to denote a pair $`x,y`$ satisfying the covering relation.
If $`L`$ is a lattice, we say that an element $`mL`$ is *meet-irreducible* if $`x>m`$, $`y>m`$ imply $`xy>m`$. If $`L`$ is complete, then we say that $`m`$ is *strictly meet-irreducible* if for all $`SL`$ such that $`sS`$ implies $`s>m`$, $`S>m`$. Note that if $`m`$ is strictly meet-irreducible, then there is a unique element $`m^{}`$ such that $`mm^{}`$.
If $`x`$, $`y`$, $`z`$, and $`wL`$ with $`xy`$ and $`zw`$, we write $`x,yz,w`$ when $`x,y`$ *transposes up to* $`z,w`$, i.e., when $`yz=x`$ and $`yz=w`$. When pairs $`x,y`$, $`z,w`$, such that $`xy`$ and $`zw`$, are related by the symmetric and transitive closure of $``$, we say that they are *projectively equivalent*, or $`x,yz,w`$.
Projective equivalence classes of coverings in modular lattices will be of fundamental importance to us. The projective equivalence class of a covering $`xy`$ will be denoted by $`[xy]`$.
A lattice $`L`$ is a *chain* if the natural ordering in $`L`$ is a total order. Also, if $`L`$ is a lattice, and $`CL`$, then $`C`$ is called a *chain in $`L`$* if in the ordering inherited from $`L`$, $`C`$ is a chain. If $`C`$ is a chain in $`L`$, then we say $`C`$ is *maximal* if no larger subset of $`L`$ is a chain in $`L`$.
A complete lattice $`L`$ is *meet-continuous* if for all $`aL`$ and $`DL`$ such that $`D`$ is *directed upward* (i.e., $`d`$, $`d^{}D`$ imply there exists $`d^{\prime \prime }D`$ such that $`dd^{\prime \prime }`$ and $`d^{}d^{\prime \prime }`$) we have
$$aD=\underset{dD}{}(ad).$$
A lattice with the dual property is called *join-continuous*.
For some other important concepts of lattice theory that we shall mention–in particular, lattices which are *algebraic* or *coalgebraic*–we refer to texts on lattice theory such at and . We will use the fact that algebraic lattices are meet-continuous, and coalgebraic lattices are join-continuous.
In section 8, we also assume an acquaintance with the basic concepts of Universal Algebra, as defined, for example, in . In particular, the concept of a *congruence* will be used, and that of the *congruence lattice* of an algebra. The reader should know that the congruence lattice of an algebra is always algebraic, and hence, meet-continuous.
The reader should know about cardinal and ordinal numbers, as used in transfinite induction. If $`\kappa `$ is a cardinal number, then $`\mathrm{Succ}\kappa `$ will stand for the smallest cardinal number strictly greater than $`\kappa `$.
## 1. Multiplicity and Multiplicity Bounds
### $`C`$-Multiplicity
If $`L`$ is a modular lattice, $`C`$ is a chain in $`L`$, and $`u`$, $`vC`$ are such that $`uv`$, then we say that the covering $`uv`$ is *in* $`C`$. If $`xy`$ is a covering in $`L`$, $`C`$ is a chain in $`L`$, and the set of coverings $`uv`$ in $`C`$ such that $`uvxy`$ has cardinality $`n`$, then we say that the *$`C`$-multiplicity* of $`xy`$ (in $`L`$), denoted by $`\mu _C[xy]`$, is $`n`$.
### Weak regularity
We say that a covering $`xy`$ is *weakly regular* if $`\mu _C[xy]`$, for maximal chains $`C`$, is a number $`\mu [xy]`$, the *multiplicity* of $`xy`$, independent of $`C`$.
If $`xy`$ is weakly regular, with finite multiplicity, then we can talk not only about the multiplicity of $`xy`$ in $`L`$, but in any interval sublattice $`\mathrm{I}[a,b]`$ of $`L`$ where $`a<b`$:
###### Theorem 1.
If $`L`$ is a modular lattice and $`xy`$ is weakly regular in $`L`$, with finite multiplicity, then given $`a`$, $`bL`$ with $`a<b`$, the number of coverings $`uv`$ equivalent to $`xy`$ in any maximal chain of elements in the interval sublattice $`\mathrm{I}[a,b]`$ is a number $`\mu ^{a,b}[xy]`$ independent of the particular chain $`C`$.
###### Proof.
Any two maximal chains $`C`$, $`C^{}\mathrm{I}[a,b]`$ can be completed to maximal chains in $`L`$ by including the elements of the same maximal chains in $`\mathrm{I}[,a]`$ and $`\mathrm{I}[a,]`$ (where we first adjoin a $``$ and a $``$ to $`L`$ if not already present). Then we use the fact that $`\mu _{\stackrel{~}{C}}[xy]=\mu _{\stackrel{~}{C}^{}}[xy]`$. ∎
### Multiplicity upper bounds and lower bounds
Let $`xy`$ be a covering in a modular lattice $`L`$. We define $`\upsilon [xy]`$, the *multiplicity upper bound* of $`xy`$ in $`L`$, to be the least cardinal number $`\nu `$ such that for every chain $`C`$ in $`L`$, $`\mu _C[xy]<\nu `$. We define $`\lambda [xy]`$, the *multiplicity lower bound* of $`xy`$ in $`L`$, to be the least cardinal $`\nu `$ such that $`\mu _C[xy]=\nu `$ for some maximal chain $`C`$.
###### Proposition 2.
If $`xy`$ is weakly regular in $`L`$, then
$$\lambda [xy]=\mu [xy]$$
and
$$\upsilon [xy]=\mathrm{Succ}\mu [xy].$$
### Distributive lattices
The $`C`$-multiplicity is severely constrained for distributive lattices:
###### Theorem 3.
If $`L`$ is a distributive lattice, $`C`$ is a maximal chain in $`L`$, and $`xy`$ is a covering in $`L`$, then $`\mu _C[xy]`$ is $`0`$ or $`1`$.
###### Proof.
Assume that $`uvzw`$ and $`u^{}v^{}zw`$. Then we claim that $`uu^{}vv^{}`$ and $`uu^{}vv^{}uv`$. For,
$`(vv^{})u`$ $`=(vv^{})(vz)`$
$`=(vz)(v^{}z)`$
$`=uu^{},`$
and
$`(vv^{})u`$ $`=(vu)(v^{}u)`$
$`=v(v^{}u)`$
$`=v(v^{}(vz))`$
$`=v(v^{}v)(v^{}z)`$
$`=v(v^{}v)w`$
$`=v.`$
Similarly, $`uu^{}vv^{}u^{}v^{}`$.
It follows that if $`uvu^{}v^{}`$, then we must have some covering $`zw`$ such that $`zwuv`$ and $`zwu^{}v^{}`$. Therefore, $`u`$, $`v`$, $`u^{}`$, and $`v^{}`$ cannot all be elements of the same chain $`C`$. ∎
###### Remark.
As a result of his theorem, if $`L`$ is a distributive lattice, and $`a<bL`$, then we can talk about the set of weakly regular coverings $`xy`$ in $`\mathrm{I}[a,b]`$. We will do so in the last section of this paper.
## 2. The Theory of Regular Coverings
### Upper regular and lower regular coverings
We say that a covering $`xy`$ in a complete modular lattice $`L`$ is *upper regular* if, for every chain $`I`$, and mapping taking elements $`iI`$ to coverings $`x_iy_i`$ of $`L`$, projectively equivalent to $`xy`$ and such that $`i<j`$ implies $`x_iy_ix_jy_j`$, we have $`_ix_i_iy_i`$ (rather than the only other possibility, for a modular lattice, which would be $`_ix_i=_iy_i`$.) The property dual to upper regularity, we call *lower regularity*. We say that a covering is *regular* if it is both upper regular and lower regular.
Clearly, whether or not a covering is upper regular, lower regular, or regular depends only on the projective equivalence class of the covering.
The importance of the concept of regularity comes from a generalization of the Dedekind-Birkhoff Theorem, proved in :
###### Theorem 4.
If $`L`$ is a complete modular lattice, and $`C`$, $`C^{}`$ are any two maximal chains in $`L`$, then for every regular covering $`xy`$, $`\mu _C[xy]=\mu _C^{}[xy]`$.
Thus, if $`xy`$ is regular, we can drop the $`C`$ from $`C`$-multiplicity and speak of the *multiplicity* $`\mu [xy]`$ of $`xy`$ in $`L`$. In other words, if $`xy`$ is regular, then $`xy`$ is weakly regular.
A partial converse to Theorem 4:
###### Theorem 5.
Let $`L`$ be a complete modular lattice. If $`xy`$ is weakly regular, and furthermore, $`\mu [xy]`$ is finite, then $`xy`$ is regular.
###### Proof.
Let $`I`$ be a chain, and let coverings $`x_iy_ixy`$ be indexed by $`I`$, such that $`i<j`$ implies $`x_iy_ix_jy_j`$. Then, for any arbitrary $`iI`$, consider the chain $`x_i\mathrm{}x_j\mathrm{}_ix_i_iy_i`$ and the chain $`x_iy_i\mathrm{}y_j\mathrm{}_iy_i`$. If we take any refinement of the first of these chains to a maximal chain $`C`$, we can find a refinement of the second chain to a maximal chain $`C^{}`$, by letting $`C^{}`$ consist of the elements of $`C`$ less than or equal to $`x_i`$, the lattice elements $`cy_i`$ for $`cC`$ such that $`x_ic_ix_i`$, and the elements of $`C`$ greater than or equal to $`_iy_i`$. By the modular law, the coverings in $`C`$ between $`x_i`$ and $`_ix_i`$ correspond in a one-to-one fashion with the coverings in $`C^{}`$ between $`y_i`$ and $`_iy_i`$, and corresponding coverings are projectively equivalent. Since $`\mu _C[xy]=\mu _C^{}[xy]`$, and that number is finite, we must have $`_ix_i_iy_i`$, proving that $`xy`$ is upper regular. Lower regularity is proved similarly. ∎
In order to apply Theorem 4, it helps to know which coverings are regular. Some preliminary observations in this direction are as follows: If $`L`$ is finite, or of finite height, then all coverings in $`L`$ are regular. It is easy to see that in any complete, modular, meet-continuous lattice, every covering is upper regular. Dually, in any complete, modular, join-continuous lattice, every covering is lower regular.
## 3. Coverings in Filter and Ideal Lattices
In this section, we will explore coverings in filter and ideal lattices. If $`L`$ is a lattice, a *filter* in $`L`$ is a nonempty subset $`F`$ such that if $`xF`$, and $`yx`$, then $`yF`$, and also, if $`x`$, $`yF`$, then $`xyF`$. If $`F`$ and $`G`$ are filters in $`L`$, we say that $`FG`$ if $`GF`$. With this partial ordering, the filters of a lattice $`L`$ with $``$ form a lattice $`\mathrm{Fil}L`$ which is complete and coalgebraic. We have $`FG=FG`$, while $`FG=\{zL:zxy\text{, for some }xF\text{ and }yG\}`$. If $`SL`$ is a nonempty subset, we write $`\mathrm{Fg}(S)`$ for the smallest (in the sense of set inclusion) filter containing $`S`$, called *the filter generated by $`S`$*. An important special case is $`\mathrm{Fg}\{x\}`$, the *principal filter* generated by $`xL`$, which is $`\{yL:yx\}`$. The mapping $`x\mathrm{Fg}\{x\}`$ is a lattice homomorphism embedding $`L`$ into $`\mathrm{Fil}L`$. As another important example of a filter, if $`m`$ is a meet-irreducible element, then we denote by $`F_{>m}`$ the set of elements of $`L`$ strictly greater than $`m`$. $`F_{>m}`$ is obviously a filter, and is principal iff $`m`$ is not just meet-irreducible, but strictly meet-irreducible.
The dual concept, that of an *ideal*, leads to the lattice of ideals $`\mathrm{Idl}L`$, which is complete and algebraic. If $`SL`$, we write $`\mathrm{Ig}S`$ for the smallest ideal containing $`S`$, and call it *the ideal generated by $`S`$*. If $`xL`$, the *principal ideal* generated by $`x`$, $`\mathrm{Ig}\{x\}=\{yL:yx\}`$ is an important example. $`\mathrm{Idl}L`$ is ordered by inclusion as opposed to $`\mathrm{Fil}L`$, which is ordered by reverse inclusion. The mapping $`x\mathrm{Ig}\{x\}`$ is a lattice homomorphism embedding $`L`$ into $`\mathrm{Idl}L`$.
Both the lattices $`\mathrm{Fil}L`$ and $`\mathrm{Idl}L`$ satisfy every lattice-theoretic identity satisfied by $`L`$; in particular, they are modular if $`L`$ is modular.
For the most part, we will concentrate our attention on filters of a modular lattice $`L`$, leaving to the reader the dualization of the statements and proofs of the theorems to yield similar results about ideals of $`L`$.
### Filter coverings $`FG`$ and $`(FG)`$
If $`L`$ is a lattice and $`SL`$, then we denote by $`(S)`$ the set of maximal elements of $`S`$ (in the partial ordering of $`L`$). If $`FG`$ is a covering in $`\mathrm{Fil}L`$, or, as we say, a *filter covering*, we will be particularly interested in $`(FG)`$. We have
###### Lemma 6.
Let $`L`$ be a lattice, and $`F`$, $`G`$, $`F^{}`$, $`G^{}\mathrm{Fil}L`$ such that $`FG`$, $`F^{}G^{}`$, and $`FGF^{}G^{}`$. Then
1. $`F^{}G^{}=F^{}(FG)`$, and
2. $`(F^{}G^{})=F^{}(FG)`$.
###### Proof.
(1): $`F^{}(FG)=(F^{}F)(F^{}G)=F^{}G^{}`$.
(2): $`x(F^{}G^{})xF^{}(FG)`$ by (1). If, in addition, $`y>x`$, then $`yG^{}`$ which implies $`yG`$. Thus, $`xF^{}(FG)`$. On the other hand, if $`xF^{}(FG)`$, then $`xF^{}G^{}`$ by (1), and $`y>xyGyF^{}G=G^{}`$. Thus, $`x(F^{}G^{})`$. ∎
### Filter coverings and maximal based filters
If $`FG`$ is a filter covering in $`\mathrm{Fil}L`$, then any $`xFG`$ determines a principal filter $`\mathrm{Fg}\{x\}`$. We have $`FG\mathrm{Fg}\{x\}(\mathrm{Fg}\{x\}G)`$. Let $`H=\mathrm{Fg}\{x\}G=\mathrm{Fg}\{x\}G`$. We say that a pair $`x,H`$, such that $`\mathrm{Fg}\{x\}H`$, is a *maximal based filter* with $`x`$ as its *base*. We say that $`x,H`$ is a maximal based filter *determined by* $`FG`$.
We now define separate concepts of $``$ and projective equivalence for maximal based filters. Suppose $`x,G`$ and $`y,H`$ are maximal based filters. We say $`x,Gy,H`$ if $`xy`$, $`yG`$, and $`H=\mathrm{Fg}\{y\}G=\mathrm{Fg}\{y\}G`$. We call the symmetric, transitive closure of this relation *projective equivalence of maximal based filters* and again use the symbol $``$. It is easy to see that the relation of projective equivalence of maximal based filters is a subset of the relation of projective equivalence on filters, restricted to maximal based filters viewed as filter coverings. That is, $`x,Gy,H`$ implies $`\mathrm{Fg}\{x\}G\mathrm{Fg}\{y\}H`$. The converse is also true:
###### Lemma 7.
Let $`L`$ be a lattice. Given two filter coverings $`FG`$ and $`HK`$ in $`\mathrm{Fil}L`$, and given $`xFG`$ and $`yHK`$, $`FGHK`$ iff $`x,\mathrm{Fg}\{x\}Gy,\mathrm{Fg}\{y\}K`$.
###### Proof.
It suffices to prove that if $`FGHK`$, then for any such $`xFG`$ and $`yHK`$, we have $`x,\mathrm{Fg}\{x\}Gy,\mathrm{Fg}\{y\}K`$.
$`yHK`$ implies $`yFG`$, so we have $`F=G\mathrm{Fg}\{y\}`$. Thus, there exists $`gG`$ such that $`gyx`$. Then we have $`gy,\mathrm{Fg}\{gy\}Gx,\mathrm{Fg}\{x\}G`$ and $`gy,\mathrm{Fg}\{gy\}Gy,\mathrm{Fg}\{y\}K`$, whence $`x,\mathrm{Fg}\{x\}Gy,\mathrm{Fg}\{y\}K`$. ∎
###### Corollary 8.
If $`xy`$ and $`zw`$, then $`xyzw`$ in $`L`$ iff $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}\mathrm{Fg}\{z\}\mathrm{Fg}\{w\}`$ in $`\mathrm{Fil}L`$.
###### Proof.
This follows from Lemma 7 and from the fact that $`xyzw`$ iff $`x,\mathrm{Fg}\{y\}z,\mathrm{Fg}\{w\}`$. ∎
### Atomic filter coverings
Let $`x,F`$ be a maximal based filter. We say that $`x,F`$, or a filter covering $`F^{}G^{}`$ such that $`\mathrm{Fg}\{x\}FF^{}G^{}`$, is *atomic* if $`F`$ is principal.
As an example, if $`mL`$ is strictly meet-irreducible, then $`m,F_{>m}`$ is an atomic maximal based filter.
###### Theorem 9.
Let $`L`$ be a modular lattice. The set of atomic maximal based filters in $`L`$, and the set of atomic filter coverings, are closed under projective equivalence. If $`FG`$ is a filter covering, and $`mFG`$ is strictly meet-irreducible, then $`FG`$ is atomic.
###### Proof.
Let $`x,Fy,G`$. If $`F`$ is principal, say $`F=\mathrm{Fg}\{x^{}\}`$, then $`G=\mathrm{Fg}\{y\}\mathrm{Fg}\{x^{}\}=\mathrm{Fg}\{yx^{}\}`$ is also principal.
On the other hand, if $`G`$ is principal, say $`G=\mathrm{Fg}\{y^{}\}`$, then let $`\overline{x}F`$ be such that $`x=\overline{x}y`$, and let $`x^{}=\overline{x}y^{}`$. We cannot have $`x^{}=x`$ because both $`\overline{x}`$ and $`y^{}`$ belong to $`F`$. Thus, by modularity, $`x^{}x`$. But, this implies that $`F=\mathrm{Fg}\{x^{}\}`$. Thus, the set of atomic maximal based filters is closed under projective equivalence, and by Theorem 7, the same is true of the set of atomic filter coverings.
If $`FG`$ and $`mFG`$ is strictly meet-irreducible, then $`FG\mathrm{Fg}\{m\}(G\mathrm{Fg}\{m\})`$. However, $`F_{>m}`$ is the unique cover of $`\mathrm{Fg}\{m\}`$ and is principal. It follows that $`G\mathrm{Fg}\{m\}=F_{>m}`$, and $`FG`$ is atomic. ∎
###### Theorem 10.
Let $`L`$ be a complete, meet-continuous modular lattice, and let $`FG`$ be an atomic filter covering. Then $`(FG)`$ is nonempty and consists of strictly meet-irreducible elements.
###### Proof.
Let $`xFG`$. Then $`FG\mathrm{Fg}\{x\}(\mathrm{Fg}\{x\}G)=\mathrm{Fg}\{x^{}\}`$ where $`x^{}x`$, because $`FG`$ is atomic. The set of elements $`y`$ such that $`yx`$ and $`yx^{}=x`$ is closed under joins of chains, by meet-continuity. Then by Zorn’s Lemma, $`(\mathrm{Fg}\{x\}\mathrm{Fg}\{x^{}\})=\mathrm{Fg}\{x\}(FG)`$ is nonempty. Thus, $`(FG)`$ is nonempty. It is easy to see that because $`FG`$ is atomic, $`(FG)`$ consists of strictly meet-irreducible elements. ∎
### Quasi-atomic filter coverings
We say that a maximal based filter $`x,F`$, or a filter covering $`F^{}G^{}`$ such that $`\mathrm{Fg}\{x\}FF^{}G^{}`$, is *quasi-atomic* if $`F`$ is not principal, but contains an element $`y`$ such that $`x<zy`$ implies $`zF`$.
As an example, if $`mL`$ is meet-irreducible, but not strictly meet-irreducible, then $`m,F_{>m}`$ is a quasi-atomic maximal based filter.
###### Theorem 11.
Let $`L`$ be a modular lattice. The set of quasi-atomic maximal based filters, and the set of quasi-atomic filter coverings, are closed under projective equivalence. If $`FG`$ is a filter covering, and $`mFG`$ is meet-irreducible but not strictly meet-irreducible, then $`FG`$ is quasi-atomic.
###### Proof.
Let $`x,Fy,G`$. If $`x,F`$ is quasi-atomic, then $`F`$ contains an element $`x^{}`$ such that $`x<zx^{}zF`$. Consider the element $`y^{}=x^{}yG=\mathrm{Fg}\{y\}F`$. If $`y<zy^{}`$, then $`x^{}z>x`$, because by modularity, $`y(x^{}z)=(yx^{})z=y^{}z=z>y`$. Thus, $`y<zy^{}zG=F\mathrm{Fg}\{y\}`$, and $`y,G`$ is quasi-atomic because if it were atomic, then $`x,F`$ would also be atomic by Theorem 9.
On the other hand, if $`y,G`$ is quasi-atomic, then there is an element $`y^{}G`$ such that $`y<zy^{}`$ implies $`zG`$. Let $`\overline{x}F`$ be such that $`x=\overline{x}y`$, and let $`x^{}=\overline{x}y^{}`$. We have $`x^{}F`$, so $`yx^{}G`$. If $`x<zx^{}`$, then by modularity, $`z=z(yx^{})=(zy)x^{}`$. But, $`y<zyy^{}`$ because if $`y=zy`$, then $`z=(zy)x^{}=x`$. Thus, $`zyG`$ and $`zF`$. It follows that $`x,F`$ is atomic or quasi-atomic, but $`x,F`$ cannot be atomic, because then $`y,G`$ would also be atomic by Theorem 9.
Now, if $`FG`$ and $`mFG`$ is meet-irreducible, but not strictly meet-irreducible, we have $`FG\mathrm{Fg}\{m\}(\mathrm{Fg}\{m\}G)`$. However, $`F_{>m}`$ is the unique cover of $`\mathrm{Fg}\{m\}`$ in $`\mathrm{Fil}L`$. Thus, $`FG\mathrm{Fg}\{m\}F_{>m}`$, which is quasi-atomic. ∎
###### Theorem 12.
Let $`L`$ be a complete, meet-continuous modular lattice. If $`FG`$ is a filter covering in $`L`$ that is quasi-atomic, then $`(FG)`$ is a nonempty set of elements of $`L`$ that are meet-irreducible, but not strictly meet-irreducible.
###### Proof.
Similar to the proof of Theorem 10. Instead of $`\mathrm{Fg}\{x\}G=\mathrm{Fg}\{x^{}\}`$, we have an $`x^{}\mathrm{Fg}\{x\}G`$ such that $`\mathrm{Fg}\{x\}G=\mathrm{Fg}\{zx<zx^{}\}`$. The set of elements $`y`$ such that $`yx`$ and $`yx^{}=x`$ is again closed under joins of chains by meet-continuity, and nonempty by Zorn’s Lemma. Thus $`(\mathrm{Fg}\{x\}(\mathrm{Fg}\{x\}G))`$ is nonempty, and so is $`(FG)`$ by Lemma 6. ∎
### Anomalous filter coverings
We say that a maximal based filter $`x,F`$, or a filter covering $`F^{}G^{}`$ such that $`\mathrm{Fg}\{x\}FF^{}G^{}`$, is *anomalous* if it is neither atomic nor quasi-atomic.
Recall that $`xL`$ is called *finitely decomposable* if $`x`$ is a finite meet of meet-irreducible elements. For an example of an anomalous filter covering, let $`xL`$ be an element which is not finitely decomposable. (This is possible only if $`L`$ does not satisfy the ascending chain condition.) Let $`G`$ be the filter generated by the set of finitely decomposable elements of $`L`$ that are greater than $`x`$. (This is the same as the filter generated by the set of meet-irreducible elements of $`L`$ that are greater than $`x`$.) We have $`xG`$ because otherwise, $`x`$ would be finitely decomposable. By Zorn’s Lemma, there is a filter $`FG`$ such that $`\mathrm{Fg}\{x\}<FG`$ and $`F`$ is minimal (in the ordering of $`\mathrm{Fil}L`$) for that property. Then by Theorems 10 and 12, $`x,F`$ is an anomalous maximal based filter, because it cannot be atomic or quasi-atomic. The following theorem shows, among other things, that this example is typical:
###### Theorem 13.
The set of anomalous maximal based filters, and the set of anomalous filter coverings, are closed under projective equivalence. If $`FG`$ is an anomalous filter covering, then $`(FG)`$ is empty, and $`FG`$ contains no elements that are finitely decomposable.
###### Proof.
The sets of aomic and quasi-atomic filter coverings are closed under projective equivalence. Since the set of anomalous filter coverings comprises the rest of the filter coverings, it is also closed under projective equivalence. If $`FG`$ and $`x(FG)`$, then clearly $`x`$ is meet-irreducible. Thus, if $`FG`$ is anomalous, $`(FG)`$ must be empty by Theorems 10 and 12. Finally, if $`x`$ is finitely decomposable, then $`x=_{i=1}^nm_i`$ where the $`m_i`$ are meet-irreducible. If $`xFG`$, then $`m_i`$ also belongs to $`FG`$ for some $`i`$, because $`G`$ is closed under finite meets. This would imply that $`FG`$ was atomic or quasi-atomic. ∎
### A counterexample
In working with the $``$ relation and filters, we might make the following conjecture: Let $`L`$ be a modular lattice, and $`F`$, $`G`$, $`H`$, $`K`$ filters such that $`FG`$, $`HK`$, and $`FGHK`$. If $`xFG`$, then there exists $`wHK`$ such that $`xw`$. However, this is false:
###### Example 14.
Consider the modular lattice known as $`M_5`$, with its elements labeled as follows:
Let $`F=M_5`$, $`G=\mathrm{Fg}\{b\}=\{b,\}`$, $`H=\mathrm{Fg}\{c\}=\{c,\}`$, and $`K=\mathrm{Fg}\{\}=\{\}`$. Then $`FGHK`$. Observe that we have $`aFG`$, but no element $`wHK`$ such that $`aw`$.
## 4. Stability Theorems
We will shortly begin to address the question of when a covering in $`\mathrm{Fil}L`$ is regular. First, however, we pose and answer some other important questions, such as, under what circumstances do regular coverings in $`L`$ remain regular after the embedding from $`L`$ into $`\mathrm{Fil}L`$ or $`\mathrm{Idl}L`$? We also examine stability of multiplicity, in case a covering is regular, and of the multiplicity upper bound. We continue to focus on $`\mathrm{Fil}L`$. A lemma:
###### Lemma 15.
Let $`L`$ be a complete modular lattice, and $`xy`$ a covering in $`L`$ which is lower regular. If $`FG`$ is a covering in $`\mathrm{Fil}L`$ and $`FG\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$, then let $`f=F`$ and $`g=G`$; we have
1. $`fg`$,
2. $`fgxy`$, and
3. $`(\mathrm{Fg}\{f\}\mathrm{Fg}\{g\})=(FG)`$.
###### Proof.
For all $`wFG`$, $`\mathrm{Fg}\{w\}G=F`$. It follows that $`wg=f`$ and hence, $`\mathrm{Fg}\{w\}\mathrm{Fg}\{g\}=\mathrm{Fg}\{f\}`$.
We must show that $`fg`$, which we will show by showing that if $`wFG`$, then we cannot have $`gw`$, or in other words, we cannot have $`_\nu g_\nu w`$ for any $`\kappa `$-tuple $`\{g_\nu \}_{\nu <\kappa }`$ of elements of $`G`$, for any cardinal number $`\kappa `$, where $`\nu `$ runs through ordinals less than $`\kappa `$. Assume the contrary, where $`\kappa `$ is the least cardinal possible. By modularity, and the fact that $`FG`$ is atomic, we can assume w.l.o.g. that $`g_\nu wg_\nu `$ for each $`\nu `$. For each ordinal $`\nu \kappa `$, let $`h_\nu =_{\nu ^{}<\nu }g_\nu ^{}`$. By the minimality of $`\kappa `$, we have $`h_\nu w`$ if $`\nu <\kappa `$, so we must have $`h_\nu wh_\nu `$ for $`\nu <\kappa `$. Note $`\kappa `$ cannot be finite, because $`G`$ is a filter. By lower regularity, we have $`h_\kappa wh_\kappa `$, contradicting the assumption that $`h_\kappa w`$.
We have $`\mathrm{Fg}\{g\}F`$, $`\mathrm{Fg}\{g\}G`$, and $`\mathrm{Fg}\{f\}=\mathrm{Fg}\{g\}F`$, whence $`\mathrm{Fg}\{f\}\mathrm{Fg}\{g\}`$, proving (1). Also, $`\mathrm{Fg}\{f\}\mathrm{Fg}\{g\}\mathrm{Fg}\{w\}\mathrm{Fg}\{w\}G`$ for any $`wFG`$. (2) follows by corollary 8.
We have $`(FG)(\mathrm{Fg}\{f\}\mathrm{Fg}\{g\})`$. For, let $`m(FG)`$. Then $`m\mathrm{Fg}\{f\}`$, and we cannot have $`m\mathrm{Fg}\{g\}`$, because then we would have $`gm`$, contradicting the fact that $`mg=f`$. $`m(\mathrm{Fg}\{f\}\mathrm{Fg}\{g\})`$ because $`m`$ is strictly meet-irreducible.
On the other hand, let $`m(\mathrm{Fg}\{f\}\mathrm{Fg}\{g\})`$. Since $`m\mathrm{Fg}\{g\}`$, $`mG`$. It suffices to show $`mF`$, because, $`m`$ being strictly meet-irreducible, $`MFG`$ will imply $`m(FG)`$. Let $`\kappa `$ be the least cardinal number such that some $`\kappa `$-tuple $`\{f_\nu \}_{\nu <\kappa }`$ of elements of $`F`$, where $`\nu `$ runs through ordinals less than $`\kappa `$, satisfies $`_\nu f_\nu m`$. $`m`$ is strictly meet-irreducible because $`\mathrm{Fg}\{f\}\mathrm{Fg}\{g\}`$ is atomic. Let $`m^{}`$ be the unique cover of $`m`$, and for each $`\nu \kappa `$, define $`u_\nu =_{\nu ^{}<\nu }f_\nu ^{}`$. We have $`mu_\nu m^{}`$ if $`\nu <\kappa `$, by the minimality of $`\kappa `$. Thus, for each $`\nu <\kappa `$, we have by modularity
$$mu_\nu m^{}u_\nu mm^{}.$$
If $`\nu ^{}<\nu <\kappa `$, then it is easy to see that
$$mu_\nu m^{}u_\nu mu_\nu ^{}m^{}u_\nu ^{}.$$
Now, if $`\kappa `$ is infinite, then by the lower regularity of $`xy`$, we must have $`mu_\kappa m^{}u_\kappa `$. However, this is absurd because $`u_\kappa m`$. Thus, $`\kappa `$ is finite. It follows that $`mF`$, proving (3). ∎
###### Theorem 16.
Let $`L`$ be a complete, meet-continuous modular lattice, and $`xy`$ a covering in $`L`$ which is regular. Then $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$ is regular in $`\mathrm{Fil}L`$.
###### Proof.
It suffices to prove upper regularity, because $`\mathrm{Fil}L`$ is coalgebraic, which implies that all coverings are automatically lower regular. Suppose given $`F_iG_i\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$, indexed by $`iI`$, for some chain $`I`$, and such that $`F_iG_iF_jG_j`$ for $`i<j`$. We must show that $`_iF_i<_iG_i`$.
For each $`iI`$, let $`M_i=(F_iG_i)`$. We have $`M_jM_i`$ for $`i<j`$ by Lemma 6, and $`_iM_i_iF_i_iG_i`$. For, if $`xM_i`$ for all $`i`$, then $`xF_i`$ for all $`i`$ so $`x_iF_i=_iF_i`$, and $`xG_i`$ for all $`i`$, so $`x_iG_i=_iG_i`$. Thus, it suffices to show that $`_iM_i`$ is nonempty.
For each $`i`$, let $`f_i=F_i`$ and $`g_i=G_i`$. Since $`xy`$ is lower regular, we have $`f_ig_i`$ and $`f_ig_ixy`$ by Lemma 15(1) and (2), and for $`i<j`$, we have $`f_ig_if_jg_j`$. For, $`f_if_j`$, $`f_ig_i`$, and $`g_if_j`$, because for any $`xF_jG_j`$, we have $`F_i=G_i\mathrm{Fg}\{x\}`$ and consequently, $`f_i=g_ix`$. It follows that $`f_i=g_if_j`$. We also have $`g_ig_j`$ and $`f_jg_j`$, whence $`g_j=g_if_j`$.
Since $`xy`$ is upper regular, we have $`fg`$ where $`f=_if_i`$ and $`g=_ig_i`$. This implies that $`(\mathrm{Fg}\{f\}\mathrm{Fg}\{g\})`$ is nonempty, by Theorem 10. Let $`x(\mathrm{Fg}\{f\}\mathrm{Fg}\{g\})`$. Then $`xf_i`$ for all $`i`$ so $`x\mathrm{Fg}\{f_i\}`$ for all $`i`$. On the other hand, $`xg`$ so there is an $`i`$ such that $`x\mathrm{Fg}\{g_i\}`$. If $`j>i`$ then $`g_jg_i`$, so $`x\mathrm{Fg}\{g_j\}`$. If $`y>x`$, then $`yg`$ and $`yG_k`$ for all $`k`$. Thus, $`x(\mathrm{Fg}\{f_j\}\mathrm{Fg}\{g_j\})`$ for all $`ji`$. On the other hand, if $`j<i`$, then by Lemma 6, $`(\mathrm{Fg}\{f_i\}\mathrm{Fg}\{g_i\})=\mathrm{Fg}\{f_i\}(\mathrm{Fg}\{f_j\}\mathrm{Fg}\{g_j\})`$, implying that $`x(\mathrm{Fg}\{f_j\}\mathrm{Fg}\{g_j\})`$ in this case as well. It follows that $`_i(\mathrm{Fg}\{f_i\}\mathrm{Fg}\{g_i\})`$ is nonempty. However, by Lemma 15(3), $`(\mathrm{Fg}\{f_i\}\mathrm{Fg}\{g_i\})=M_i`$ for each $`i`$. Thus, $`_iM_i`$ is nonempty, and $`FG`$ is regular. ∎
Now, we consider the multiplicity:
###### Lemma 17.
If $`L`$ is a complete modular lattice and $`xy`$ is lower regular in $`L`$, then for any maximal chain $`CL`$ and any maximal chain $`\overline{C}\mathrm{Fil}L`$ refining the image of $`C`$ in $`\mathrm{Fil}L`$, $`\mu _{\overline{C}}[\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]=\mu _C[xy]`$.
###### Proof.
Let $`C`$ be a maximal chain in $`L`$, and $`\overline{C}`$ a maximal chain in $`\mathrm{Fil}L`$ refining the image of $`C`$.
If we have a principal filter $`\mathrm{Fg}\{u\}\overline{C}`$, then we must have $`uC`$. For, if $`u\overline{C}`$, then there must exist $`cC`$ such that $`u`$ and $`c`$ are not comparable. But, $`\mathrm{Fg}\{c\}\overline{C}`$, so either $`\mathrm{Fg}\{c\}<\mathrm{Fg}\{u\}`$, implying $`c<u`$, or $`\mathrm{Fg}\{u\}<\mathrm{Fg}\{c\}`$, implying $`u<c`$.
If $`F`$, $`G\overline{C}`$ are principal and such that $`FG`$, say $`F=\mathrm{Fg}\{u\}`$ and $`G=\mathrm{Fg}\{v\}`$, then we must have $`u`$, $`vC`$ with $`uv`$, and if $`FG\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$ then $`uvxy`$.
Let $`F`$, $`G\overline{C}`$ be such that $`FG`$ and it is not true that $`F`$ and $`G`$ are both principal. We will show that $`FG`$ is not projectively equivalent to $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$.
If $`F`$ is principal and $`G`$ is not, then $`FG`$ is not of atomic type, is not projectively equivalent to $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$, and does not count in the multiplicity of $`xy`$.
The case $`G`$ principal and $`F`$ non-principal cannot occur, because $`FG`$.
The only remaining case is that both $`F`$ and $`G`$ are non-principal. We can assume that $`FG`$ is atomic and lower regular, since otherwise $`FG\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$ is impossible. If we had $`cC`$ with $`cFG`$, then we would have $`F<\mathrm{Fg}\{c\}<G`$; thus, we must have $`CF=CG`$ in order to have $`FG`$. Denote this set by $`D`$. Since $`FG`$ is atomic, let $`zFG`$ and $`wG`$ with $`zw`$. For each $`dD`$, we have $`dzdw`$, and for $`d<d^{}D`$, we have $`dzdwd^{}zd^{}w`$. Since $`zw`$ is lower regular in $`L`$, we would have
$$(D)z=\underset{dD}{}dz\underset{dD}{}dw=(D)w.$$
Now, if we had $`DCD`$, we would have $`\mathrm{Fg}\{D\}<FG`$, implying that $`Dz=Dw`$, which is impossible.
The only other possibility is $`DD`$. In this case, we claim that we must have $`Dw`$ and $`DzC`$. For, if $`cCD`$, then $`\mathrm{Fg}\{c\}F`$, implying that $`cDz`$. Then because $`C`$ is maximal, we must have $`DzCD`$ and $`Dw=DD`$. However, this contradicts the fact that $`FC=D`$, because $`DzF`$. It follows that the case $`FG`$ atomic, lower regular, and neither $`F`$ nor $`G`$ principal is impossible. ∎
It follows from the Lemma that we have
###### Theorem 18.
If $`xy`$ is regular in $`L`$, then $`\mu [\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]=\mu [xy]`$.
Finally, we examine the stability properties of the multiplicity upper bound.
###### Theorem 19.
If $`xy`$ is a covering in a modular lattice $`L`$, then $`\upsilon [xy]\upsilon [\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]`$, with equality if $`\upsilon [xy]`$ is finite or countable. In any case, $`\upsilon [xy]`$ infinite implies $`\upsilon [\mathrm{Fg}\{x\}\mathrm{Fg}\{g\}]`$ infinite.
###### Proof.
Let $`C`$ be a chain in $`L`$. Then the image of $`C`$ in $`\mathrm{Fil}L`$ can be refined to a maximal chain $`\overline{C}`$ in $`\mathrm{Fil}L`$, and it is clear that $`\mu _C[xy]\mu _{\overline{C}}[\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]`$.
Now, let
$$F_1G_1F_2G_2\mathrm{}F_nG_n$$
where $`F_iG_i\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$ for all $`i`$. Then $`u_iF_iG_i`$, $`v_iG_i`$ such that $`u_iv_i`$ and $`u_iv_ixy`$. For each $`i`$, define $`c_i=_{ji}u_i`$, $`d_i=v_i_{j>i}u_i`$. Then $`c_i=d_iu_i`$, $`d_iG_i`$, $`c_iF_iG_i`$, and $`d_iv_i`$, implying that $`c_id_ixy`$ and
$$c_1d_1c_2d_2\mathrm{}c_nd_n.$$
It follows that $`\upsilon [xy]n`$, and combined with the fact that $`\upsilon [xy]\upsilon [\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]`$, this implies that $`\upsilon [xy]=\upsilon [\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]`$ if $`\upsilon [xy]`$ is finite or countable, and that $`\upsilon [xy]`$ is infinite if $`\upsilon [\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]`$ is. ∎
## 5. $`\upsilon [xy]`$, $`\lambda [\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]`$, and $`\lambda [\mathrm{Ig}\{x\}\mathrm{Ig}\{y\}]`$
In this section, we consider the relationship between the multiplicity upper bound of a covering, and the multiplicity lower bound of the corresponding filter covering or ideal covering. As usual, we focus on filter coverings, leaving the dual result to be stated by the reader.
Consider the function $`\mathrm{\Lambda }:`$, where $``$ stands for the natural numbers, defined recursively as follows:
$$\mathrm{\Lambda }(n)=\{\begin{array}{cc}0,\hfill & n=0\hfill \\ 1+\mathrm{\Lambda }(\sqrt{n}1),\hfill & n>0.\hfill \end{array}$$
###### Lemma 20.
We have
1. $`\mathrm{\Lambda }(n)0`$ for all $`n`$
2. $`\mathrm{\Lambda }`$ is increasing; i.e., $`n<n^{}\mathrm{\Lambda }(n)\mathrm{\Lambda }(n^{})`$
3. $`lim_n\mathrm{}\mathrm{\Lambda }(n)=\mathrm{}`$
4. $`\mathrm{\Lambda }(n)\sqrt{n}`$ for all $`n`$.
###### Proof.
(1) is clear.
To prove (2), note that we have $`\mathrm{\Lambda }(0)=0`$ and $`\mathrm{\Lambda }(1)=1`$, so (2) is true for $`n<n^{}1`$. If (2) is true for $`n<n^{}\overline{n}>1`$ then for $`n\overline{n}+1`$, $`\sqrt{n}1\overline{n}`$, and the square root function is also increasing, whence $`\mathrm{\Lambda }(n^{})\mathrm{\Lambda }(n)=\mathrm{\Lambda }(\sqrt{n^{}}1)\mathrm{\Lambda }(\sqrt{n}1)0`$ if $`n<n^{}\overline{n}+1`$. Thus, (2) follows by induction.
To prove (3), we use (2) and note that if $`\mathrm{\Lambda }(n)=m`$, then $`\mathrm{\Lambda }((n+1)^2)=m+1`$.
A computation shows that the inequality (4) holds for all $`n10`$. Suppose (4) holds for $`n\overline{n}>10`$, and let us prove it is true for $`n=\overline{n}+1`$. We have by the induction hypothesis
$$\mathrm{\Lambda }(\overline{n}+1)=1+\mathrm{\Lambda }(\sqrt{\overline{n}+1}1)1+\sqrt{\sqrt{\overline{n}+1}1}.$$
Squaring, we have
$$\mathrm{\Lambda }(\overline{n}+1)^2=1+2\sqrt{\sqrt{\overline{n}+1}1}+\sqrt{\overline{n}+1}13\sqrt{\overline{n}+1}\overline{n}+1.$$
Thus, (4) holds for $`n=\overline{n}+1`$, and by induction, for all $`n`$. ∎
###### Theorem 21.
Let $`L`$ be a complete, meet-continuous modular lattice, and $`xy`$ a covering in $`L`$. Let
$$b_1a_1b_2a_2\mathrm{}b_na_n,$$
where $`xya_ib_i`$ for all $`i`$. If $`\overline{C}`$ is a maximal chain in $`\mathrm{Fil}L`$, then $`\mu _{\overline{C}}[\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]\mathrm{\Lambda }(n)`$.
###### Proof.
Let $`m_1L`$ be maximal for the property that $`m_1a_1`$ but $`m_1b_1`$. (By meet-continuity, the set of such elements is closed under joins of chains, so a maximal such element exists by Zorn’s Lemma.) Then $`m_1`$ is strictly meet-irreducible, and has a unique cover, $`m_1^{}=m_1b_1`$.
Now, let $`m_2L`$ be maximal for the property that $`m_2a_2`$, $`m_2b_2`$, and $`m_2m_1`$. We cannot have $`m_2=m_1`$, because $`m_1a_1b_2`$. $`m_2^{}=m_2b_2`$ is the unique cover of $`m_2`$ in the interval $`\mathrm{I}[,m_1]`$, because if $`x>m_2`$ and $`xm_1`$, we must have $`xb_2`$.
Similarly, we successively choose $`m_3`$, $`\mathrm{}`$, $`m_n`$ such that $`m_i`$ is maximal among elements $`x`$ such that $`xa_i`$, $`xb_i`$, and $`xm_{i1}`$, and we obtain covers $`m_i^{}=m_ib_i`$. We have
$$m_1^{}m_1m_2^{}m_2\mathrm{}m_n^{}m_n.$$
For each $`i`$, let $`F_i`$ be the join of all elements of $`\overline{C}`$ containing $`m_i`$, and $`G_i`$ the meet of all elements of $`\overline{C}`$ not containing $`m_i`$. We have $`F_i`$, $`G_i\overline{C}`$ because the maximal chain $`\overline{C}`$ is be closed under joins and meets. Clearly, $`F_iG_i`$ for all $`i`$.
The mapping $`iF_iG_i`$ sends each $`i`$ to the unique covering in $`\overline{C}`$ such that $`m_iF_iG_i`$, and thus partitions the ordered set $`\{\mathrm{\hspace{0.17em}1},\mathrm{},n\}`$ into intervals. If $`\{i,i+1,\mathrm{},j\}`$ is one of these intervals, then we claim that $`m_i^{}G_i`$. For, if $`i=1`$ then $`m_i`$ is strictly meet-irreducible, and so we must have $`m_i^{}G_1`$. If $`i>1`$, then we have $`m_{i1}F_{i1}`$, so $`m_{i1}G_i`$, because, $`\overline{C}`$ being a chain, $`G_iF_{i1}`$. If $`m_i^{}`$ did not belong to $`G_i`$, then there would be an element $`gG_i`$ such that $`m_i^{}g=m_i`$. Then we would have $`m_i^{}(gm_{i1})=m_i`$, but $`m_igm_{i1}`$ because $`G_i`$ is closed under meets and does not contain $`m_i`$. This is impossible, because $`m_i`$ is meet-irreducible in $`\mathrm{I}[,m_{i1}]`$. Thus, the claim that $`m_i^{}G_i`$ is proved. It follows that $`F_iG_i\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$.
Clearly, we have $`m_{i+1}^{}`$, $`m_{i+1}`$, $`\mathrm{}`$, $`m_j^{}`$, $`m_jF_iG_i`$, as well as $`m_i`$. Thus, since we have shown that $`F_iG_i`$ is atomic, $`m_{i+1}`$, $`\mathrm{}`$, $`m_j`$ have unique covers $`\overline{m}_{i+1}`$, $`\mathrm{}`$, $`\overline{m}_jG_i`$. Defining $`\overline{m}_k^{}=\overline{m}_km_k^{}`$ for $`k=i+1`$, $`\mathrm{}`$, $`j`$, we obtain
$$\overline{m}_{i+1}^{}\overline{m}_{i+1}\mathrm{}\overline{m}_j^{}\overline{m}_j$$
and each $`\overline{m}_k\overline{m}_k^{}xy`$.
Now, either the number of intervals is $`\sqrt{n}`$, or the cardinality of the largest interval is $`\sqrt{n}`$. In the first case, we have $`\mu _{\overline{C}}[\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]\mathrm{\Lambda }(n)`$ by Lemma 20(4). In the second case, by induction on $`n`$ (and noting that recursive application of the construction of this proof will find filter coverings above $`F_iG_i`$ and therefore distinct from it), we also have $`\mu _{\overline{C}}[\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]1+\mathrm{\Lambda }(\sqrt{n}1)=\mathrm{\Lambda }(n)`$. ∎
###### Corollary 22.
Let $`L`$ be a complete, meet-continuous modular lattice. If $`\upsilon [xy]`$ is infinite, then so is $`\lambda [\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]`$.
## 6. Upper Regularity of Filter Coverings and Joins of Chains
In this section, we consider the issue of upper regularity of filter coverings, and show anomalous filter coverings cannot be regular, because they cannot be upper regular. We also give an example of an atomic filter covering, in an algebraic lattice, which is not upper regular, showing that upper regularity alone of $`xy`$ does not imply upper regularity of $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$, even if the lattice is meet-continuous.
Filter coverings are always lower regular, because $`\mathrm{Fil}L`$ is coalgebraic, thus join-continuous. Thus, if a filter covering is upper regular, it must be regular.
A necessary condition for a filter covering to be upper regular is easy to state:
###### Theorem 23.
Let $`L`$ be a complete modular lattice, and $`F`$, $`G\mathrm{Fil}L`$ such that $`FG`$. If $`FG`$ is upper regular, then $`FG`$ is closed under joins of chains.
###### Proof.
Suppose $`C`$ is a chain in $`FG`$. For each $`cC`$, define $`F_c=\mathrm{Fg}\{c\}`$ and $`G_c=GF_c=GF_c`$. Then for all $`cC`$, $`FGF_cG_c`$, and if $`c`$, $`c^{}C`$ with $`cc^{}`$, we have $`F_cG_cF_c^{}G_c^{}`$. Since $`FG`$ is upper regular, $`_cF_c_cG_c`$. However, $`_cF_c=\mathrm{Fg}\{C\}`$. If $`CG`$, then we would have $`_cF_c=_cG_c`$. Thus, $`CFG`$. ∎
###### Corollary 24.
If $`L`$ is a complete modular lattice, $`F`$, $`G\mathrm{Fil}L`$ with $`FG`$, and $`FG`$ is anomalous, then $`FG`$ is not upper regular.
###### Proof.
By the Theorem, if $`FG`$ is upper regular, then $`FG`$ is closed under joins of chains. Then, by Zorn’s Lemma, $`FG`$ has maximal elements. However, this is impossible for anomalous $`FG`$ by Theorem 13. ∎
Some sufficient conditions for the preceding necessary condition to hold:
###### Theorem 25.
Let $`L`$ be a complete, modular, meet-continuous lattice, and $`FG`$ a covering in $`\mathrm{Fil}L`$ which is atomic or quasi-atomic. Then $`FG`$ is closed under joins of chains.
###### Proof.
Use meet-continuity as in the proof of Theorem 10 or Theorem 12. ∎
###### Theorem 26.
Let $`L`$ be a complete modular lattice, and $`xy`$ a covering in $`L`$ which is upper regular. If $`HK\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$, then $`HK`$ is closed under joins of chains.
###### Proof.
Let $`C`$ be a chain in $`HK`$, and let $`cC`$. We have $`\mathrm{Fg}\{c\}K=\mathrm{Fg}\{q\}`$ for some $`qK`$ such that $`cq`$ and $`cqxy`$. For each $`c^{}`$, $`c^{\prime \prime }C`$ such that $`cc^{}c^{\prime \prime }`$, we have $`c^{}qc^{}c^{\prime \prime }qc^{\prime \prime }`$. Then $`C=_{c^{}c}c^{}_{c^{}c}(c^{}q)`$ by the upper regularity of $`xy`$. If we had $`CK`$, then we would have $`c=qCK`$, which is absurd. It follows that $`CHK`$. ∎
There follows an example of a meet-continuous lattice, having an atomic filter covering which is not regular:
###### Example 27.
Let $`V`$ be the infinite-dimensional real vector space of sequences of real numbers, only a finite number of which are nonzero. Let $`L`$ be the lattice of subspaces of $`V`$. For each finite set $`S`$ of cardinality $`2`$, consider the subspace
$$A_S=\{a_0,a_1,\mathrm{}V:sSa_s=0\},$$
and the subspace
$$B_S=\{b_0,b_1,\mathrm{}V:s,s^{}Sb_s=b_s^{}\}.$$
Note that $`A_SB_S`$ for each $`S`$, and if $`SS^{}`$ then $`A_S^{}B_S^{}A_SB_S`$.
For each $`n`$, let
$$U_n=\{A_S:sSsn\},$$
and
$$U_n^{}=\{B_S:sSsn\};$$
the sets $`U_n`$, $`U_n^{}`$ are bases for filters $`F_n=\mathrm{Fg}U_n`$ and $`G_n=\mathrm{Fg}U_n^{}`$.
We have $`F_nG_n`$ for all $`n`$. For, if $`HF_nG_n`$, then $`A_SH`$ for some $`S`$ such that $`sSsn`$, but there does not exist an $`S^{}`$ such that $`sS^{}s>n`$ and $`B_S^{}H`$. In particular, $`B_SH`$. Then $`B_SHv`$ where $``$ is the field of real numbers, and $`v`$ has $`1`$ in positions $`s`$ such that $`sS`$, and $`0`$ elsewhere. $`HHv`$, because $`v`$ is an atom of $`L`$.
We have $`A_SB_SH(HB_S`$. Now let $`A_{\widehat{S}}`$ be a basic element of $`F_n`$, and we will show that $`HB_{\overline{S}}A_{\widehat{S}}`$ for some $`\overline{S}`$. If we had $`A_{\widehat{S}}G_n`$ already, this would be trivial. If $`A_{\widehat{S}}F_nG_n`$, however, we have $`B_{\widehat{S}}G_n`$. Then let $`\overline{S}=S\widehat{S}`$. We have $`A_{\overline{S}}=B_{\overline{S}}H=B_{\overline{S}}B_SH=B_{\overline{S}}A_S=A_{\overline{S}}A_{\widehat{S}}`$. Thus, $`F_nG_n`$. This argument has also shown that the covering $`F_nG_n`$ is atomic.
Also, if $`n^{}>n`$ then we claim that $`F_nG_nF_n^{}G_n^{}`$. To prove this, it suffices to prove $`F_nG_nF_{n+1}G_{n+1}`$. We have $`F_n<F_{n+1}`$ and $`G_n<G_{n+1}`$, and if $`S`$ is such that $`sSs>n`$, and $`A_SF_n`$, then either $`A_SF_{n+1}`$, or $`n+1S`$. We may assume that $`\mathrm{card}S>2`$ since the $`A_S`$ for such $`S`$ and such that $`A_SF_n`$ form a base for $`F_n`$. Then $`A_S=B_{\{n+1,j\}}A_{S\{n+1\}}`$, where $`jn+1`$ is any other element of $`S`$. Thus, $`F_n=G_nF_{n+1}`$. On the other hand, if $`HG_{n+1}`$, then there is an $`S`$ such that $`sSs>n+1`$ and $`B_SH`$. $`B_SG_n`$ and $`B_SF_{n+1}`$, since $`A_SB_S`$, so $`B_SG_nF_{n+1}`$. Thus, $`G_nF_{n+1}G_{n+1}`$. The claim follows.
We already proved that $`F_nG_n`$ is atomic for each $`n`$. However, $`F_n=G_n=\{V\}`$. Thus, the coverings $`F_nG_n`$ are not upper regular, and so are not regular.
## 7. Regularity of Filter Coverings
###### Lemma 28.
Let $`L`$ be a complete lattice and $`SL`$, where $`S\mathrm{}`$. Then the following are equivalent:
1. $`\mathrm{Fg}S`$ is principal, and
2. $`S\mathrm{Fg}S`$.
Let $`L`$ be a complete, meet-continuous modular lattice. If $`FG`$ is an atomic or quasi-atomic filter covering, such that the equivalent conditions of the Lemma are satisfied, then we say that $`FG`$ is *principally bounded*.
###### Theorem 29.
Let $`L`$ be a complete, meet-continuous modular lattice, and $`FG`$ a covering in $`\mathrm{Fil}L`$ which is atomic or quasi-atomic. If every filter covering $`HK`$ such that $`HKFG`$ is principally bounded, then $`FG`$ is regular.
###### Proof.
Let $`I`$ be a chain, and $`F_iG_i`$ be filter coverings projectively equivalent to $`FG`$ and such that $`i<j`$ implies $`F_iG_iF_jG_j`$. For each $`i`$, let $`q_i=(F_iG_i)`$, and let $`q=_iq_i`$. We have $`q_iF_i`$ for each $`i`$, so $`q_iF_i`$.
On the other hand, we have $`q_iG_i`$. For, $`q_i`$ some element of $`(F_iG_i)`$, whence $`q_iG_i`$. Also, $`G_jG_i`$ for $`j>i`$, so $`q_iG_j`$ in that case.
If $`j<i`$, then $`G_i=G_jF_i=G_jF_i`$. However, $`q_iF_i`$. Thus, $`q_iG_j`$.
So, $`q_iG_j`$ for all $`i`$ and $`j`$. Now, $`F_jG_j`$ is closed under joins of chains, by Theorem 25. Thus, $`q=_iq_iG_j`$ for all $`j`$. It follows that $`q_iG_i`$.
Thus, $`q_iF_i_iG_i`$, implying that $`_iF_i_iG_i`$, and that $`FG`$ is upper regular, hence regular. ∎
###### Corollary 30.
Let $`L`$ be a complete, meet-continuous distributive lattice. If $`FG`$ is a covering in $`\mathrm{Fil}L`$ which is atomic or quasi-atomic, then $`FG`$ is regular.
###### Proof.
If $`HKFG`$, then $`HK`$ is not anomalous, so $`(HK)`$ is nonempty by Theorem 10 and Theorem 12. Furthermore, by distributivity, it has cardinality one, proving that $`HK`$ is principally bounded. The Corollary then follows from Theorem 29. ∎
###### Theorem 31.
Let $`L`$ be a complete, meet-continuous modular lattice. If $`xy`$ is a covering in $`L`$ such that $`\upsilon [xy]`$ is finite, then $`[\mathrm{Fg}\{x\},\mathrm{Fg}\{y\}]`$ is regular, with multiplicity equal to $`\upsilon [xy]`$.
###### Proof.
If $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$ were not principally bounded, $`(\mathrm{Fg}\{x\}\mathrm{Fg}\{y\})`$ would contain a sequence of strictly meet-irreducible elements $`m_1`$, $`m_2`$, $`\mathrm{}`$ such that for all $`n`$, $`_{i<n}m_i_{in}m_i`$. For each $`i`$, let $`m_i^{}`$ be the unique cover of $`m_i`$. Then if $`yL`$ is such that $`ym_i`$, we have $`ym_iym_i^{}`$. For, we must have $`ym_im_i^{}`$, whence $`ym_i^{}=ym_i`$. If we had $`ym_i=ym_i^{}`$, then the elements $`y`$, $`m_i`$, $`m_i^{}`$, $`ym_i`$, and $`ym_i`$ would form a sublattice isomorphic to the lattice $`N_5`$, which cannot happen in a modular lattice.
It follows that
$$m_1m_2^{}m_1m_2m_1m_2m_3^{}m_1m_2m_3\mathrm{},$$
with $`_{in}m_i(_{i<n}m_i)m_n^{}`$, the general covering in the chain, equivalent to $`xy`$ for all $`i`$. This sequence of elements can be refined to a maximal chain $`C`$ such that $`\mu _C[xy]`$ is infinite. However, this is contrary to the assumption that $`\upsilon [xy]`$ is finite.
Thus, $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$ is principally bounded, and regular by Theorem 29. ∎
## 8. Lattice-theoretic Chief Factors
Suppose we have an algebra $`A`$ (in the sense of universal algebra) which has a modular congruence lattice. Then we can apply the preceding theory to the congruence lattice $`\mathrm{Con}A`$ and talk about the chief factors of $`A`$, obtaining a generalization of the nice multiplicity result seen in the Jordan-Holder Theorem.
### Coverings of rank $``$ and lattice-theoretic chief factors of rank $``$
If we have a regular covering $`xy`$ in the lattice $`(L)`$, where the functor $``$ is some composite of the functors $`\mathrm{Fil}`$ and $`\mathrm{Idl}`$, then we say that $`xy`$ is a *covering of $`L`$ of rank $``$*. Then, if $`L=\mathrm{Con}A`$, we say that a covering in $`L`$ of rank $``$ is a *lattice-theoretic chief factor* of $`A`$ *of rank $``$*.
In the theories of finite groups and finite-height modules, where the lattices involved have finite height, it is standard practice to assign a group or module to a covering, obtaining a *chief factor*, or, in case $`A`$ is a module, a *composition factor*, of $`A`$. In order to do something similar for an arbitrary congruence-modular algebra $`A`$, it is necessary to assign some type of algebraic object to each lattice-theoretic chief factor. We leave to future investigations the question of the manner in which this may be done generally. (We have taken some small steps toward such a theory in , , and .) However, the lattice-theoretic chief factors themselves are of interest, because their multiplicities are invariants of the algebra.
Thus, in the remainder of this section, unless otherwise specified, $`L`$ will denote the lattice $`\mathrm{Con}A`$, for some algebra $`A`$ such that $`\mathrm{Con}A`$ is modular. Since $`\mathrm{Con}A`$ is algebraic, we also are assuming that $`L`$ is meet-continuous.
### The case when $`L=\mathrm{Con}A`$ satisfies the descending chain condition
If $`L`$ satisfies the descending chain condition, then all coverings in $`L`$ are lower regular, all filters are principal, and all filter coverings are of atomic type. In fact, $`\mathrm{Fil}LL`$. Since $`L`$ is meet-continuous, coverings in $`L`$ (and $`\mathrm{Fil}L`$) are also upper regular. Thus, in this situation, all coverings in $`L\mathrm{Fil}L`$ are regular. This result was stated but not proved in ; it must be admitted, however, that as an example for the application of the ideas in that paper, and of lattice-theoretic chief factors, it is vacuous.
### The case when $`L=\mathrm{Con}A`$ is distributive
A better example presents itself when $`L`$ is distributive. Then, Corollary 30 shows that every filter covering $`FG`$ of atomic or strictly quasi-atomic type is regular. We do not know how many coverings there may be in $`\mathrm{Fil}L`$ which are not of anomalous type. However, we note that $`\mathrm{Fil}L`$ is provided with a profuse supply of coverings, by which we mean, somewhat informally, that if $`\alpha `$, $`\beta L`$ with $`\alpha <\beta `$, then there is at least one covering in $`\mathrm{Fil}L`$ between $`\mathrm{Fg}\{\alpha \}`$ and $`\mathrm{Fg}\{\beta \}`$. (This is easy to prove using Zorn’s Lemma.) We can then apply the dual of Corollary 30 to $`\mathrm{Idl}\mathrm{Fil}L`$, because $`\mathrm{Fil}L`$ is coalgebraic. The conclusion is that there is a profuse supply of regular coverings in $`\mathrm{Idl}\mathrm{Fil}L`$, i.e., a profuse supply of lattice-theoretic chief factors of $`A`$ of rank $`\mathrm{Idl}\mathrm{Fil}`$. We have regularized coverings in $`\mathrm{Fil}L`$ by the embedding into $`\mathrm{Idl}\mathrm{Fil}L`$.
Those coverings in $`\mathrm{Idl}\mathrm{Fil}L`$ which are not known to be regular can be regularized by considering them in $`\mathrm{Fil}\mathrm{Idl}\mathrm{Fil}L`$, where they become regular by Corollary 30, as long as $`L`$ is distributive. And so on. By this method, any covering in $`(L)`$, for any functor $``$ which is a nonempty composite of $`\mathrm{Idl}`$ and $`\mathrm{Fil}`$, can be regularized by applying either $`\mathrm{Fil}`$, or $`\mathrm{Idl}`$, and similarly, any covering in any distributive lattice whatever can be regularized by applying either $`\mathrm{Idl}\mathrm{Fil}`$, or $`\mathrm{Fil}\mathrm{Idl}`$.
Because of Theorem 3, all of the regular coverings that arise in this way have multiplicity one.
### Modular but not distributive $`L=\mathrm{Con}A`$
In this case, multiplicities higher than $`1`$ are possible. If $`xy`$ is a covering in $`L`$, such that $`\upsilon [xy]`$ is finite, then $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$ is regular by Theorem 31. Thus, the embedding from $`L`$ into $`\mathrm{Fil}L`$ regularizes such coverings, and the multiplicity of the corresponding filter covering is $`\upsilon [xy]`$ by Theorem 19.
On the other hand, if $`\upsilon [xy]`$ is infinite, then by Corollary 22, so is $`\lambda [\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}]`$. Thus, any maximal chain $`\overline{C}`$ in $`\mathrm{Fil}L`$ has an infinite number of coverings in it that are equivalent to $`\mathrm{Fg}\{x\}\mathrm{Fg}\{y\}`$, and we can say that the multiplicity of the filter covering is infinite, even though we may not be able to say that that multiplicity is a well-defined cardinal number.
## 9. The Steps in the Proof of a Theorem
Another application of these ideas is a method of formalizing the steps necessary and sufficient to prove a given proposition from given premises. Our treatment of this will use a simple Logic framework.
Suppose we have a set $`𝐏`$ of “propositions,” which can in principal be determined to either be true, or not. Then we have two truth values $`𝐓`$ and $`𝐅`$, and for any proposition $`P`$ we can say that $`𝒯(P)`$ (the truth value of $`P`$) takes values $`𝐓`$ and $`𝐅`$. Given any $`n`$ propositions $`P_1`$, $`\mathrm{}`$, $`P_n`$, and any $`n`$-ary function $`f`$ with arguments consisting of truth values, we can formulate a new, synthetic proposition $`f(\stackrel{}{P})`$ with truth value $`f(𝒯(P_1),\mathrm{},𝒯(P_n))`$. If we consider the truth values as elements of the two-element boolean algebra $`\{𝐓,𝐅\}`$, then the functions obtainable by compositions of the ordinary logical connectives give us this, because the two-element boolean algebra has the property of being *primal*–i.e., the property that every finitary function can be constructed from the basic operations. We will use the symbols $``$, $``$, $`\neg `$, $``$ with their usual meanings, along with $`𝐓`$ and $`𝐅`$. In fact, it is convenient to replace our original set of propositions $`𝐏`$ by a boolean algebra $`𝐁`$ free on $`𝐏`$ as set of generators. (Or, if $`𝐏`$ already has some or all of the logical connectives, by a quotient of such a free boolean algebra.) The assignment $`𝒯`$ of truth values can then be extended to a boolean algebra homomorphism from $`𝐁`$ to the two-element boolean algebra. Henceforth, *proposition* shall mean an element of $`𝐁`$.
We will write $``$ and $``$ for the maximum and minimum elements of $`𝐁`$. The *underlying lattice* of $`𝐁`$ is just $`𝐁`$, forgetting the unary operation $`\neg `$. A *filter* of $`𝐁`$ is the same as a filter of the underlying lattice.
Given some sort of calculus of proving propositions, consisting of finitary rules of inference, we assume that the rules of inference include a small set of *trivial* rules of inference, and otherwise we call them *nontrivial* rules of inference. The trivial rules of inference are the rule that we can infer $`PQ`$ from $`P`$ and $`Q`$, for any elements $`P`$, $`Q𝐁`$, and the rule that for any $`P`$, $`Q𝐁`$, if $`PQ`$, then we can infer $`Q`$ from $`P`$. Note that *modus ponens*, the rule that we can infer $`Q`$ from $`P`$ and $`PQ`$ (or $`\neg PQ`$) will thus be considered a trivial rule of inference, because $`P(\neg PQ)=(P\neg P)(PQ)=PQQ`$.
Note that if the ordering of $`𝐁`$ provides that $`PQ`$ whenever $`Q`$ can be proved from $`P`$, then the second trivial rule of inference would actually subsume all the rules of inference, rendering our analysis of the situation vacuous. Thus, we want to consider a situation where $`𝐁`$ does not have such an ordering.
We say $`PQ`$ ($`SQ`$, where $`S𝐁`$) if $`Q`$ can be proved from $`P`$ (from elements of $`S`$) using both trivial and nontrivial rules of inference.
###### Theorem 32.
If $`T𝐁`$, then the following are equivalent:
1. $`T`$ is a filter of the underlying lattice of $`𝐁`$, and
2. $`T`$ is closed under application of the trivial rules of inference, and contains $``$.
We call a set $`T`$, satisfying the equivalent conditions of Theorem 32, a *pretheory*. A *theory* is quite often defined as a set of propositions closed under the rules of inference. We use the term *pretheory* to suggest that a filter in $`𝐁`$ is a forerunner of a theory, and will not have occasion to mention theories further.
Now, $`𝐁`$ may or may not be complete or algebraic when viewed as a lattice, and may or may not have any coverings at all, but $`\mathrm{Fil}𝐁`$ is complete and coalgebraic, and has a profuse supply of coverings. Thus, by the dual of corollary 30, $`\mathrm{Idl}\mathrm{Fil}𝐁`$ has a profuse supply of regular coverings. (We use the word *profuse* in the informal sense of the previous section.) We call regular coverings in $`(\mathrm{Fil}𝐁)`$, where $``$ is a functor as used in section 8, *steps of order $``$*. That is, an step of order $``$ is a regular covering in $`𝐁`$ of rank $`\mathrm{Fil}`$. For example, if $`T`$ and $`T^{}`$ are pretheories such that $`TT^{}`$, $`\mathrm{Ig}\{T\}\mathrm{Ig}\{T^{}\}`$ is regular (by the dual of Corollary 30) and is an step of order $`\mathrm{Idl}`$.
If $`T`$ is a pretheory, $`P`$ is a proposition, and $``$ is a functor as before, then we call the steps of order $``$ in $`\mathrm{I}[\varphi (T)\varphi (\mathrm{Fg}\{P\}),\varphi (T)]`$, where $`\varphi :\mathrm{Fil}𝐁\mathrm{Fil}𝐁`$ is the natural embedding, the *steps of order $``$ in the proof of $`P`$ from $`T`$.* Note that we do not assume $`TP`$.
If we have an instance of a rule of inference which infers $`Q`$ from $`P_1`$, $`\mathrm{},P_n`$, then we say that that instance *covers* the set of steps of order $``$ which occur in $`\mathrm{I}[\varphi (\mathrm{Fg}\{Q_iP_i\}),\varphi (\mathrm{Fg}\{_iP_i\})]`$. If we have a set $`𝐍`$ of instances of rules of inference, then we say that $`𝐍`$ *covers* the union of the sets of steps covered by the individual instances $`N𝐍`$. We also say that $`𝐍`$ covers any smaller set of steps.
Recall that an *ultrafilter* is a cover of $``$ in $`\mathrm{Fil}𝐁`$.
###### Theorem 33.
Let $`\mathrm{Ig}\{T_1\}\mathrm{Ig}\{T_2\}`$ be an step of order $`\mathrm{Idl}`$, where $`T_1T_2`$ (i.e., an step of order $`\mathrm{Idl}`$, of atomic type), and let $`PT_1T_2`$. Then $`T_2\mathrm{Fg}\{\neg P\}`$ is an ultrafilter of $`𝐁`$, and $`T_2\mathrm{Fg}\{\neg P\}T_1T_2`$. This is a one-one correspondence of ultrafilters with projective equivalence classes of steps of order $`\mathrm{Idl}`$, and the steps of order $`\mathrm{Idl}`$ in the proof of $`P`$, from a pretheory $`T`$, correspond to the ultrafilters that contain $`T`$ but not $`P`$.
###### Proof.
If $`T_2\mathrm{Fg}\{\neg P\}`$ were $``$, there would be $`QT_2`$ such that $`Q\neg P=`$. Then we would have to have $`QP`$, implying that $`PT_2`$ which is not true. Thus, $`T_2\mathrm{Fg}\{\neg P\}`$ is an atom, i.e., an ultrafilter. The ultrafilters $`U`$ of $`𝐁`$ all form coverings $`U`$ which determine distinct projective equivalence classes, because given two ultrafilters $`U`$ and $`U^{}`$, if we had $`UU^{}`$, the multiplicity of $`U`$ in $`\mathrm{I}_{\mathrm{Fil}𝐁}[,UU^{}]`$ would be $`2`$, and this is impossible. ∎
###### Theorem 34.
The set of steps (of any order $``$) covered by an instance of a trivial rule of inference is empty.
###### Theorem 35.
If $`T`$ is a pretheory and $`P`$ is a proposition, then any proof of $`P`$ from $`T`$ covers the steps (of any order $``$) in the proof of $`P`$ from $`T`$.
###### Proof.
Let $`N_i`$, $`i=1`$, $`\mathrm{}`$, $`n`$ be the instances of rules of inference in a proof, in order. Let pretheories $`T_i`$, $`i=0`$, $`\mathrm{}`$, $`n`$ be defined by $`T_0=T`$, $`T_i=T_{i1}\mathrm{Fg}\{Q_i\}`$ for $`0<in`$, where $`Q_i`$ is the conclusion of $`N_i`$. For each $`i>0`$, let the set of steps of order $``$ in $`\mathrm{I}[\varphi (T_i),\varphi (T_{i1})]`$ be $`E_i`$, and the set of steps of order $``$ covered by $`N_i`$, by $`E_i^{}`$.
We have $`E_iE_i^{}`$. For, if $`N_i`$ infers $`Q_i`$ from the finite set of propositions $`S_i`$, we have
$$T_i,T_{i1}\mathrm{Fg}\{\overline{Q}_i\},\mathrm{Fg}\{\overline{Q}_i\}T_{i1},$$
where $`\overline{Q}_i=Q_iS_i`$. However, $`\mathrm{Fg}\{\overline{Q}_i\}T_{i1}\mathrm{Fg}\{S_i\}`$, because, the $`n`$-tuple $`N_1,\mathrm{},N_n`$ being a proof, $`S_iT_{i1}`$.
Thus, $`_iE_i_iE_i^{}`$, but the left side is the set of steps of order $``$ in the proof of $`P`$ from $`T`$, and the right side is the set of steps of order $``$ covered by the proof. ∎
Let $`T`$, $`T^{}`$ be pretheories such that $`T^{}T`$, and let $`𝐍`$ be a set of instances of rules of inference. For each $`N𝐍`$, let $`S_N`$ be the (finite) set of premises of $`N`$, and $`Q_N`$ the conclusion. If $`T^{}`$ is the join (intersection) of all pretheories $`\overline{T}T`$ such that $`N𝐍`$ and $`S_N\overline{T}`$ imply $`Q_N\overline{T}`$, then we say that $`𝐍`$ *generates $`T^{}`$ from $`T`$*. In this case, $`T^{}`$ consists of all propositions provable from $`T`$ using the elements of $`𝐍`$ as the only instances of nontrivial rules of inference:
###### Theorem 36.
Let $`T`$, $`T^{}`$ be pretheories with $`T^{}T`$, and let $`𝐍`$ be a set of instances of rules of inference which generates $`T^{}`$ from $`T`$. Then $`T^{}`$ is the set of propositions $`P`$ such that there is a finite sequence of elements of $`𝐍`$ that can be refined to a proof of $`P`$ from $`T`$ by adding instances of trivial rules of inference.
###### Proof.
Let $`\stackrel{~}{T}`$ be that set of propositios, and we will show that $`T^{}=\stackrel{~}{T}`$. Since $`T^{}`$ is generated from $`T`$ by $`𝐍`$, $`T^{}`$ is the intersection (join) of all pretheories $`\overline{T}T`$ such that $`N𝐍`$ and $`S_N\overline{T}`$ imply $`Q_N\overline{T}`$.
Clearly, $`N𝐍`$ and $`S_N\stackrel{~}{T}`$ imply $`Q_N\stackrel{~}{T}`$, because we can construct a proof of $`Q_N`$ from proofs of the elemtns of $`S_N`$. Thus, $`\stackrel{~}{T}T^{}`$.
On the other hand, suppose that $`\overline{T}T`$ is such that $`N𝐍`$ and $`S_N\overline{T}`$ imply $`Q_N\overline{T}`$, and let $`P\stackrel{~}{T}`$. The existence of a proof of $`P`$ from $`T`$ using instances from $`𝐍`$ implies that $`P\overline{T}`$. Thus, $`\overline{T}\stackrel{~}{T}`$, so $`T^{}\stackrel{~}{T}`$.
Thus, $`T^{}=\stackrel{~}{T}`$. ∎
Finally, a theorem which shows that covering the steps in the proof of $`P`$ from $`T`$ is not only necessary, but sufficient:
###### Theorem 37.
Given pretheories $`T`$, $`T^{}`$ such that $`T^{}T`$, a set $`𝐍`$ of instances of rules of inference that generates $`T^{}`$ from $`T`$, and a proposition $`P`$, then we have $`PT^{}`$ iff $`𝐍`$ covers the steps in the proof of $`P`$ from $`T`$.
###### Proof.
If $`PT^{}`$, then the conclusion follows from Theorems 35 and 36.
If $`PT^{}`$, then we have
$$\mathrm{Fg}\{P\}T^{},T^{}\mathrm{Fg}(P)T,(\mathrm{Fg}\{P\}T)T^{}$$
and we have $`(\mathrm{Fg}\{P\}T)T^{}T`$. Thus, the steps in the interval $`\mathrm{I}[\mathrm{Fg}\{P\}T^{},T^{}]`$ are a subset of the set of steps in the proof of $`P`$ from $`T`$. By Zorn’s Lemma, there is an ideal $`J\mathrm{Idl}\mathrm{Fil}𝐁`$ such that $`\mathrm{Ig}\{\mathrm{Fg}\{P\}T^{}\}<J\mathrm{Ig}\{T^{}\}`$. The covering $`\mathrm{Fg}\{J\}\mathrm{Fg}\{\mathrm{Ig}\{T^{}\}\}`$ is an step (of order $`\mathrm{Fil}\mathrm{Idl}`$) in the proof of $`P`$ from $`T`$ that is not covered by $`𝐍`$. ∎
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# Near-infrared adaptive optics observations of galaxy clusters: Abell~262 at z=0.0157, J1836.3CR at z=0.414, and PKS~0743-006 at z=0.994
## 1 Introduction
Adaptive optics systems using natural and laser guide stars are an important observational tool that allow large ground based telescopes to operate at or close to the diffraction limit. Considerable improvements and successes have been obtained in installing such systems at several sites (Davies et al. Davies1999 (1999), Davies et al. Davies1998 (1998), Glindemann et al. Glindemann (1997), Quirrenbach et al. Quirrenbach (1997), Drummond et al. Drummond (1998), Hubin Hubin (1997), Max et al. Max (1997), Arsenault et al. Arsenault (1994)). However, it remains challenging to use them efficiently especially in the field of extragalactic observations. We have concentrated on three galaxy clusters at different redshifts for which adaptive optics observations were possible due to the presence of sufficiently bright reference stars in the corresponding fields.
The galaxy cluster Abell 262 (R.A.(2000) = 01<sup>h</sup>52.1<sup>m</sup>, DEC(2000) = 3540) is one of the most conspicuous condensations in the Pisces-Perseus super cluster. It has a systemic velocity of 4704 km s<sup>-1</sup> ($`z=0.0157`$, Giovanelli and Haynes Giovanelli1985 (1985)) and an Abell radius of $`r_A`$ = 1.75. It has been extensively studied in X-rays and in the radio. It is a spiral-rich cluster, characterized by the presence of a central X-ray source positioned on the D galaxy NGC 708 right at the center of the cluster. The distribution of galaxies in projection on the sky as well as in redshift space have been studied by Melnick and Sargent (Melnick (1977)), Moss and Dickens (Moss1977 (1977)), Gregory et al. (Gregory (1981)), and Fanti et al. (Fanti (1981)). The large number of spirals in this cluster as well as the presence of a central X-ray source and its low redshift make Abell 262 an ideal candidate to study the properties of member galaxies, such as HI content and star formation activity. As in many other rich galaxy clusters the member galaxies of Abell 262 show an HI deficiency towards the center of the cluster. For Abell 262 this phenomenon has been investigated by Giovanelli et al. (Giovanelli1982 (1982)), Giovanelli & Haynes (Giovanelli1985 (1985)) and others.
We used the new MPIA-MPE ALFA adaptive optics system at the Calar Alto 3.5 m telescope to observe two of the Abell 262 cluster members – UGC 1344 and UGC 1347 – at subarcsecond resolution. For UGC 1347 the observations were carried out using the ALFA laser guide star (LGS) and a nearby natural guide star (NGS) for tip-tilt correction. To our knowledge UGC 1347 is the first extragalactic source for which LGS assisted observations have been performed. For the UGC 1344 observations we used a nearby NGS as a wavefront reference.
As the two higher redshift clusters we selected J1836.3CR (R.A.(2000) = 13<sup>h</sup>45<sup>m</sup>, DEC(2000) = $``$0053) at a redshift of $`z=0.414`$ (Couch et al. Couch (1998)) and an area around the quasar PKS 0743-006 at a redshift of $`z=0.994`$ (Hewitt & Burbidge Hewitt (1993)). Both fields contain bright guide stars with $`m_\mathrm{V}12`$ mag. that can be used as natural guide stars for adaptive optics observations.
Couch et al. (Couch (1998)) presented a catalogue of faint southern galaxy clusters identified on high-contrast film derivatives of a set of Anglo-Australian Telescope photographic plates. The cluster J1836.3CR is one of them. A bright star ($`m_\mathrm{V}`$ = 12 mag) is located about 60<sup>′′</sup> north of 4 prominent cluster members for which redshifts have been determined. For three of four galaxies spectroscopy (Couch et al. Couch (1998)) indicates a redshift of $`z=0.415\pm 0.003`$ and one galaxy has a redshift of $`z=0.319`$. Couch et al. (Couch (1998)) use a cluster redshift of $`z=0.414`$. At this redshift 1<sup>′′</sup> corresponds to a linear distance of 6.8 kpc.
PKS 0743-006 is a quasar (R.A.(2000) = 07<sup>h</sup>45<sup>m</sup>53.37<sup>s</sup>, DEC(2000) = $``$004411.4<sup>′′</sup>) of visual magnitude $`m_\mathrm{V}`$ = 17.1 mag at a redshift $`z=0.994`$ (Hewitt & Burbidge Hewitt (1993)). The radio spectrum has a convex shape with possible variability around the peak occurring between 5 and 10 GHz. Tornikoski et al. (Tornikoski (1993)) find this source strongly variable at 90 GHz. Variability by a few tenths of a magnitudes is also reported in the NIR (White et al. White (1988)). On the milliarcsecond angular resolution scale at cm wavelengths this object shows a classical core-jet structure (Stanghellini et al. Stanghellini (1997)). Within the errors, the whole radio flux density is accounted for by this structure. A natural guide star for adaptive optics observations is located at only about 12.2<sup>′′</sup> northeast of the quasar.
In section 2 we describe the observations and data reduction as well as the adaptive optics systems we used. In section 3 we present the observational results and the data analysis for UGC 1347 (section 3.1) and for UGC 1344 (section 3.2) in conjunction with data available in the literature. In section 3.3 we outline the results we obtained for a sample of 11 spiral galaxies in the Abell 262 cluster and 15 spiral galaxies in the Abell 1367 cluster. In section 4 then we discuss the star formation activity in the observed cluster galaxies and give a summary and conclusions in section 5.
## 2 Observations and data reduction
Our new high spatial resolution observations in Abell 262 were carried out using the OMEGA-CASS camera mounted to the laser guide star adaptive optics system ALFA at the Calar Alto 3.5 m telescope. The observations for the two higher redshift clusters were obtained with the ESO AO system ADONIS. In the following we give a brief description of the two systems.
### 2.1 ALFA
The performance goal of ALFA is to achieve a 50% Strehl-ratio at 2.2 $`\mu `$m under average seeing conditions (0.9<sup>′′</sup>), with good sky coverage. The adaptive optics (Glindemann et al. Glindemann (1997)) and the sodium laser guide star (Quirrenbach et al. Quirrenbach (1997), Davies et al. Davies1998 (1998)) have been designed and built as a joint project between MPIA in Heidelberg and MPE in Garching, both in Germany. The system is installed at the German/Spanish 3.5 m telescope on Calar Alto near Almeria, Spain.
The laser used for generating the artificial guide star is a high power continuous-wave dye laser. It is installed in the coudé lab of the telescope, and the laser beam is fed along the coudé train until it is picked off near the primary mirror and directed into a 50-cm launch telescope. The launched laser power is around 3 W, and produces a $`m_\mathrm{V}=`$ 9–10 mag sodium guide star. The tip-tilt correction is achieved using a natural guide star, currently with a limiting magnitude $`m_\mathrm{V}15`$ mag. The laser can be used for high order wavefront correction. In the wavefront sensor there are several lenslet arrays which can be interchanged, and the positions of the resulting laser beacon centroids in the Shack-Hartmann sensor are determined and used to derive coefficients of Zernike or Karhunen-Loeve modes which are then used to control a 97-actuator deformable mirror. The loop was closed on the laser guide star in September 1997, and it was first used to improve an image in December 1997 at a sampling rate of 60 Hz and correcting 7 modes plus tip and tilt.
OMEGA-CASS is a near-infrared camera for the Cassegrain focus of the 3.5 m telescope at Calar Alto, which is specialized for use at high spatial resolution and has been developed at the MPIA, Heidelberg. It is based around a Rockwell 1024<sup>2</sup> pixel HAWAII array, and has capabilities for broad and narrow band imaging, spectroscopy, and polarimetry over the 1.0–2.5$`\mu `$m wavelength range. When used in conjunction with ALFA ($`f`$/25), the pixel scales available are 0.04<sup>′′</sup>, 0.08<sup>′′</sup>, and 0.12<sup>′′</sup> per pixel.
### 2.2 ADONIS
For the observations of the two higher redshift clusters we used the ESO adaptive optics system ADONIS (Beuzit et al. Beuzit (1994)). This system is operated at ESO’s 3.6 m telescope at La Silla, Chile, and includes the SHARP II+ camera built at MPE. The atmospheric wavefront distortions are measured with a Shack-Hartmann sensor at visible wavelengths and are corrected by a deformable mirror with 52 piezo actuators. This mirror is driven by a closed control loop with a correction bandwidth of up to 17 Hz. The natural guide star within the near-infrared isoplanatic patch must be brighter than about m$`{}_{\mathrm{V}}{}^{}=13`$ mag. The SHARP II+ camera (Hofmann et al. Hofmann (1992), Eisenhauer et al. Eisenhauer (1998)) is based on a 256<sup>2</sup> pixel NICMOS III detector. The wavelength range of our observations covers the atmospheric J, H and K bands. Compared with the Johnson K band, we used a somewhat narrower K filter (1.99 – 2.32 $`\mu `$m) in order to reduce the thermal background.
### 2.3 The data
Goal of our investigation was to exploit structural information on galaxy cluster members from adaptive optics and seeing limited images and interpret the results making use of all available quantities and known correlations. The photometric quality of the data is of the order of 0.10<sup>m</sup> in the K- and H-band and 0.15<sup>m</sup> in the J-band. The sources were mainly selected on the basis of availability during the allocated observing time, the presence of bright AO reference stars, and the availability of additional literature data. The sample described at the end of section 2.3.1 was selected on the basis HI deficiency and beeing located within or outside the cluster‘s Abell radii.
#### 2.3.1 Abell 262 and Abell 1367
The broad-band J, H, K images as well as first K-band adaptive optics data of the cluster member UGC 1347 were taken on November 10 and 12, 1997. K-band AO data of UGC 1344 and UGC 1347 as well as the direct imaging data of other cluster members were obtained on December 6 and 7, 1997. Many individual 5-second exposures were taken in all three bands and coadded after sky subtraction, flat-fielding and correcting for bad pixels. The integration times, pixel scales and angular resolutions of the final co-added images are listed in Tab. 1.
Calibration of the NIR data was accomplished by observation of the standard star $`\xi `$<sup>2</sup>Ceti. Sky data were taken separately 120<sup>′′</sup> east of UGC 1344 and UGC 1347. For the other cluster members a median sky was obtained from the 4 to 5 settings taken with different offset positions from the target sources. The galaxies UGC 1344 and UGC 1347 had sufficiently bright reference stars nearby to observe them with the ALFA adaptive optics system. In order to estimate the image improvement we took data in open loop before and after the closed loop exposures. For UGC 1347 the wavefront data on the laser guide star were taken at a sampling rate of 60 Hz through a 3$`\times `$3 lenslet array with field sizes of 3<sup>′′</sup> diameter. Correcting for a total of 7 Zernike modes plus tip and tilt a disturbance rejection bandwidth of up to 5 Hz was achieved. The tip and tilt information was derived from a nearby natural reference star.
In Fig. 1 we show an image through the TV-guider shortly after the LGS-supported AO observations of UGC 1347 were made.
The image is focused on the stars – so the LGS appears as a defocused image at the tip of the Rayleigh cone. We had placed the LGS between the tip-tilt reference star and the nucleus of UGC 1347, such that the image of the star which is also in the field of view of the NIR camera could be used as the point spread function (PSF) with the same degree of correction as the nucleus and most of the galaxy. On the reference star for UGC 1344 the loop was closed using the same lenslet array and a camera frame rate of 100 to 200 Hz resulting in a slightly higher rejection bandwidth. Although we did not reach the diffraction limit due to the low sampling rate and the small number of subapertures, definite improvements in angular resolution were achieved. The corresponding full-width-half-maximum (FWHM) values are given in Tab. 1. In the case of UGC 1347 the image improvement could independently be monitored via a star in the same field at approximately the same separation from the reference star as the target object. The two stellar images agreed very well with each other indicating that all the sources were well within the isoplanatic patch and that the images of the reference stars can safely be taken as the PSF to clean the galaxy images.
In addition to the adaptive optics data we took seeing limited images with exposure times of 10 minutes each of 9 galaxies in Abell 262 and 15 spirals in Abell 1367 (see section 3.3 for further details).
#### 2.3.2 J1836.3CR and PKS 0743-006
We observed these clusters using the SHARPII+ camera together with the ESO adaptive optics system ADONIS on the 3.6 m telescope on La Silla, Chile. The observations were conducted during the nights from April 26 till May 1, 1996. A bright star was used to lock the AO system. We took a series of 60 second exposures in a dither mode in the near-infrared J, H, and K bands, using pixel scales of 0.05 and 0.10 ”/pixel. The total integration times and the angular resolution measured on a PSF reference are listed in Tab. 2 and Tab. 3.
## 3 Results
In the following we will present the results obtained for the two galaxies UGC 1347 and UGC 1344 that were observed with the ALFA adaptive optics system as well as for a sample of galaxies located in the inner and outer part of the Abell 262 and Abell 1367 clusters. We also describe source properties at other wavelengths as well as quantities that we derived from them. The description of this derivation is given in detail for UGC 1347. For UGC 1344 we have used the same approach and only summarize the results. We regard the corresponding analysis as an important consistency check between our own data and the data and correlations available in the literature. Data on external galaxies even at lower or medium redshift will always be sparse and it is required to make use of all the knowledge available to allow for a full comparison to what is known in the local universe and at different redshifts.
Although our NIR data has a subarcsecond resolution we extracted K-band and H$`\alpha `$ fluxes in larger apertures to conduct a starburst analysis in section 4. The reason for this is that the radio and bolometric luminosities for individual source components especially the nucleus, disk, and southern component in UGC 1347 can only be estimated indirectly and can probably only be attributed to larger regions. We have chosen a circular aperture of 7.2<sup>′′</sup> diameter corresponding to a linear size of 2.2 kpc.
### 3.1 UGC 1347
UGC 1347 is an almost face-on SBc galaxy located at R.A.(2000) = 01<sup>h</sup>52<sup>m</sup>45.9<sup>s</sup> and DEC(2000) = 363709<sup>′′</sup> approximately 57 north of the center of Abell 262, well within the region in which the largest amount of HI deficiency is observed. There is a bright field star (PPM 1111, $`m_\mathrm{V}`$ = 11.5 mag) located about 37<sup>′′</sup> to the southeast of the galaxy. The HI content of UGC 1347 was first studied by Wilkerson (Wilkerson (1980)). Velocity fields and intensity maps were obtained in HI by Bravo-Alfaro (Bravo-Alfaro (1997)) and in H$`\alpha `$ by Amram et al. (Amram (1994)). Amram et al. (Amram (1994)) quote an inclination of $`i=30^{}`$. Oly and Israel (Oly (1993)) measured the 327 MHz radio continuum flux density of UGC 1347, and the far-infrared flux densities as measured by IRAS can be found in the IRAS point source catalogue (Lonsdale et al. Lonsdale (1985)). The HI and H$`\alpha `$ data indicate a systemic velocity of UGC 1347 of 5524 km s<sup>-1</sup> (Wilkerson Wilkerson (1980)) and 5478 km s<sup>-1</sup> (Amram et al. Amram (1994)), respectively. Here we assume that the difference of approximately 800 km s<sup>-1</sup> between the cluster velocity of 4704 km s<sup>-1</sup> and the systemic velocity is due to the motion of the galaxy within the cluster. We therefore adopt for UGC 1347 the cluster distance of 63 Mpc or a redshift of $`z=0.0157`$ assuming $`H_0`$ = 75 km s<sup>-1</sup> Mpc<sup>-1</sup>. At this distance 1<sup>′′</sup> corresponds to about 310 pc.
#### 3.1.1 Near-infrared emission from UGC 1347
The NIR emission of UGC 1347 is dominated by two almost equally bright components at a separation of 8.85<sup>′′</sup> (or 2.74 kpc) oriented approximately north-south. In Fig. 2 we show the NIR continuum emission from UGC 1347 together with the digitized sky survey V-band image in Fig. 3.
In Fig. 4 we show NIR intensity profile cuts of the nucleus, the bright off-nuclear source, and the star PM 1111 in both open loop and closed loop. We have used the image of PPM 1111 as the point spread function and deconvolved the NIR continuum image of UGC 1347 with a Lucy-Richardson algorithm (Lucy Lucy (1974)). From the comparison of the images in Figures 2 and 3 and the deconvolved image it is evident that the northern component coincides with the nucleus of UGC 1347 and has an extent of about 1<sup>′′</sup> (corrected for the FWHM of the PSF; see Tab. 1) corresponding to a diameter of about 310 pc. The southern compact component is located at the southern tip of the galaxy bar and is unresolved compared to the 0.40<sup>′′</sup> FWHM PSF. We estimate an upper limit to its angular extent of 0.15<sup>′′</sup> corresponding to less than 45 pc.
The J, H, and K flux densities listed in Tab. 4 were measured in 4.8<sup>′′</sup>, 3.6<sup>′′</sup>, 2.4<sup>′′</sup>, 1.2<sup>′′</sup> and 0.72<sup>′′</sup> diameter circular apertures centered on each component.
In Fig. 5 we show the locations of the multi-aperture data in the J $``$ H, H $``$ K two color diagram. The graph indicates that the nuclear colors are in agreement with a stellar disk population reddened with an $`A_\mathrm{V}`$ of about 4 mag corresponding to about 0.6 magnitudes of reddening in the K band. The value of $`A_\mathrm{V}`$ = 4 mag is large with respect to color variations due to differences in stellar population, age, and metallicity from a “normal“ Sc disk population. Here we assumed a screen model for the extinction. In case of a mixed model the extinction may be even large. In addition to simple reddening there may also be some contribution from hot dust to the nuclear NIR emission. In the case of the southern component, however, the colors from a normal stellar disk population are apparently more influenced by additional emission from hot dust and an extinction of $`A_\mathrm{V}`$ 2 mag (corresponding to the arrows in the figure ledgend). In both components the reddening and the influence from hot dust emission increase with decreasing aperture size. This indicates a decrease in dilution by a surrounding or underlying stellar population unaffected by reddening. The southern component is probably similar to the red knot found in the nearby spiral NGC 7552 (Schinnerer et al. Schinnerer (1997)) and is likely region of recent active star formation in the disk. However, its emission may very well be contaminated by red super-giants (see results of our starburst analysis in section 4.2). The contribution of hot dust in the nucleus may be indicative for star formation activity there as well.
In order to analyze the NIR data in conjunction with other data taken from the literature we used a starburst model as described in section 4.
#### 3.1.2 The K-band luminosity
Our K-band images allow us to calculate the nuclear K-band luminosity $`L_\mathrm{K}`$ and compare it to an estimate of the K-band luminosity from the extended disk and southern component of UGC 1347. The results are listed in Tab. 5 and Tab. 6.
The K-band luminosity is calculated via
$$L_\mathrm{K}[L_{}]=1.14\times 10^4\times D[\mathrm{Mpc}]^2\times S_\mathrm{K}[\mathrm{mJy}]$$
(1)
where $`D`$ is the distance in Mpc and $`S_\mathrm{K}`$ is the 2.2 $`\mu `$m flux density in mJy (Krabbe et al. Krabbe (1994)). We found K-band flux densities of 5.5 mJy for the extended nucleus and 3.7 mJy for the southern compact component measured in 7.2<sup>′′</sup> diameter apertures. In the previous section we have shown that the K-band flux density of the southern component is clearly contaminated by emission of hot dust. The J $``$ H colors are less effected by hot dust emission and correspond to that of an un-reddened stellar population. We have therefore calculated the stellar K-band flux density from the H-band magnitude using the standard spiral disk colors with H $``$ K = 0.21 (e.g. Frogel et al. Frogel (1978)). The K-band flux density of the southern component corrected for the contribution of hot dust is then reduced to about 3 mJy ($`m_\mathrm{K}`$ = 13.3 mag) resulting in a K-band luminosity of 1.5$`\times `$10<sup>8</sup> $`L_{\mathrm{}}`$ .
In order to estimate the K-band luminosity of the disk we used the total H-band magnitude of 10.48 mag measured by Gavazzi et al. (Gavazzi1996b (1996)) For a mean H $``$ K color of 0.21 (Frogel et al. Frogel (1978)) this results in a total K-band flux of the order of 49 mJy corresponding to a total K-band luminosity of $`L_\mathrm{K}`$ = 2.1$`\times `$10<sup>9</sup> $`L_{\mathrm{}}`$ . Correcting for the contribution of the nucleus and the southern component we obtain a K-band flux and luminosity for the disk of UGC 1347 of about 40 mJy and $`L_\mathrm{K}`$ = 1.6$`\times `$10<sup>9</sup> $`L_{\mathrm{}}`$ .
#### 3.1.3 The Lyman continuum luminosity
For the overall galaxy as well as for individual components an estimate of the H$`\alpha `$ luminosities have been obtained using the continuum and H$`\alpha `$ line images kindly provided by Amram et al. (Amram (1994); P. Amram and M. Marcelin 1998, private communication). Both the continuum and H$`\alpha `$ line images were taken simultaneously with the same spectral resolution. They have a field of view of $`4.9^{}\times 4.9^{}`$ and include all of UGC 1347. Main purpose of these observations was to derive the H$`\alpha `$ velocity field. In order to determine flux densities we used an estimate of the total H$`\alpha `$ continuum flux density of UGC 1347 to calibrate the data. The H$`\alpha `$ continuum flux density was obtained via a linear interpolation between the flux densities derived from the total H-band and V-band magnitudes as given in Gavazzi & Boselli (Gavazzi1996a (1996)). Although uncertain by probably 30% this estimate allows us to further probe the consistency of the available data with our own measurements and starburst analysis. The corresponding calibration factor between the measured and calculated H$`\alpha `$ continuum was applied to the H$`\alpha `$ line data. Finally the Lyman continuum luminosity $`L_{\mathrm{Lyc}}`$ then was derived from the H$`\alpha `$ flux density $`F_{\mathrm{H}\alpha }`$ and the source distance $`D`$ via:
$$L_{\mathrm{Lyc}}[L_{}]=5.6\times 10^{17}\times F_{\mathrm{H}\alpha }[\mathrm{ergs}^1\mathrm{cm}^2]\times D[\mathrm{Mpc}]^2$$
(2)
(following Osterbrock Osterbrock (1989)). Without extinction correction we obtain about 2.3$`\times `$10<sup>8</sup> $`L_{\mathrm{}}`$ for the overall Lyman continuum luminosity of UGC 1347. In 7.2<sup>′′</sup> diameter apertures centered on the nucleus and on the southern component we obtain approximately 10<sup>7</sup> $`L_{\mathrm{}}`$ and 2$`\times `$10<sup>7</sup> $`L_{\mathrm{}}`$ , respectively. In the following we assume the following values for the extinction, which are based on the JHK measurements and (for the disk) on comparisons to other galaxies: nucleus $`A_\mathrm{V}`$ = 3 mag, southern component $`A_\mathrm{V}`$ = 1 mag, and disk $`A_\mathrm{V}`$ 1 mag. Furthermore we use $`A_{\mathrm{H}\alpha }=0.8A_\mathrm{V}`$ (Draine Draine (1989)). The corresponding extinction-corrected Lyman continuum luminosities are given in Tab. 5 and Tab. 6.
#### 3.1.4 The gas content and total mass
The HI content of UGC 1347 was measured by Wilkerson (Wilkerson (1980)) using the Arecibo telescope with a 3.2 beam. The distribution of atomic hydrogen in UGC 1347 was studied by Bravo-Alfaro (Bravo-Alfaro (1997)) using the Westerbork Synthesis Radio Telescope. The interferometric measurements show that the HI gas extends well beyond the optical disk although the FWHM of the distribution is in approximate agreement with the optical extent. The HI (Bravo-Alfaro Bravo-Alfaro (1997)) and H$`\alpha `$ (Amram et al. Amram (1994)) rotation curves are in good agreement. The image suggests a slight HI line flux enhancement to the south-east and extensions to the west and north. Assuming a distance of 67 Mpc Wilkerson (Wilkerson (1980)) obtained 3.5$`\times `$10<sup>9</sup> $`M_{\mathrm{}}`$ of atomic hydrogen gas containing most of the overall HI content of the galaxy. Scaled to our adopted distance of 63 Mpc this results in M<sub>HI</sub> = 3.2$`\times `$10<sup>9</sup> $`M_{\mathrm{}}`$ .
No direct measurement of the molecular gas mass in UGC 1347 is available. As an estimate we use the IRAS far-infrared (FIR) flux densities and a dust temperature of 22 K (see the following section) and the molecular hydrogen mass to L<sub>FIR</sub> correlation by Young & Scoville (Young (1991)). We estimate a molecular hydrogen mass of approximately $`M_{\mathrm{H}_2}`$ = 3.2$`\times `$10<sup>9</sup> $`M_{\mathrm{}}`$ . The resulting $`M_{\mathrm{H}_2}/M_{\mathrm{HI}}`$ ratio of 1 is then consistent with the value expected for late type spirals (Young & Scoville Young (1991)). Although the determination of the molecular gas mass is uncertain, it is very unlikely that it is wrong by a large factor, say 10, since the age derived from $`(M_{\mathrm{tot}}M_{\mathrm{HI}}M_{\mathrm{H}_2})/L_\mathrm{K}`$ is fully consistent with the $`L_{\mathrm{bol}}/L_{\mathrm{Lyc}}`$ and $`L_\mathrm{K}/L_{\mathrm{Lyc}}`$ ratios (see section 4.2) for tis age. We note that we implicitly assume a standard $`N_{\mathrm{H}_2}/I(\mathrm{CO})`$ conversion factor. We also note, that $`M_{\mathrm{H}_2}/M_{\mathrm{HI}}`$ is of the order of 0.1 for optical selected galaxies, while $`M_{\mathrm{H}_2}/M_{\mathrm{HI}}1`$ is found for infrared selected samples. This is consistent with UGC 1347 being listed in the IRAS point source catalogue (see below).
The measured HI line width of 144 km s<sup>-1</sup> (Wilkerson Wilkerson (1980)) and the H$`\alpha `$ velocity field corrected for the inclination of $`i=30^{}`$ (Amram et al. Amram (1994)) indicates a full velocity width covered by the rotation curve of $`\mathrm{\Delta }v_0300`$ km s<sup>-1</sup>. Following Shostak (Shostak (1978)) and Heckman et al. (Heckman (1978)) this allows to estimate the total dynamical mass of UGC 1347 as $`M_{\mathrm{dyn}}`$ = (0.5–1.0)$`\times `$10<sup>11</sup> $`M_{\mathrm{}}`$ . The resulting total gas to dynamical mass ratio ($`M_{\mathrm{HI}}+M_{\mathrm{H}_2})/M_{\mathrm{dyn}}`$ of about 5% to 10% is in agreement with typical values found for spiral galaxies (Shostak Shostak (1978)).
#### 3.1.5 The FIR luminosity
The FIR luminosity can be derived using the 60 $`\mu `$m and 100 $`\mu `$m IRAS flux densities of $`S_{60}`$ = 1.40 Jy of $`S_{100}`$ = 3.84 Jy as listed in the IRAS point source catalogue. At both wavelengths UGC 1347 is fully contained in the large IRAS beams. Following the formalism given by Lonsdale et al. (Lonsdale (1985)) and Fairclough (Fairclough (1985)) we find a total $`L_{\mathrm{FIR}}=1.18\times 10^{10}h^2`$ $`L_{\mathrm{}}`$ . From the 60 $`\mu `$m and 100 $`\mu `$m data we calculate a dust color temperature of 22 K assuming an emissivity proportional to $`\lambda ^1`$ (Hildebrand Hildebrand (1983)) and a silicate to graphite ratio of 7:3 (Whittet Whittet (1981)).
We can also estimate how large the disk and nuclear contributions to the FIR luminosity are. Here we assume that a dominant fraction of the disk FIR emission originates in diffuse interstellar dust and gas clouds which are heated by the interstellar UV radiation field of the stellar disk population. These clouds have correspondingly large volume and disk area filling factors. To obtain a first order estimate of the disk and nuclear contributions we assume that for this population of clouds the FIR flux densities can be estimated by adopting the relation found for “cirrus” emission in our Galaxy (De Vries et al. De Vries (1987), Helou Helou (1986)). These relations have already successfully been applied to extragalactic objects by Eckart et al. (Eckart1990 (1990)) for Centaurus A and by Jackson et al. (Jackson (1991)) for NGC 2903. Following De Vries et al. (De Vries (1987)) the far-IR flux density $`S_{100}`$ of the “cirrus” emission in Ursa Major can be obtained via
$$S_{100}=aN_{\mathrm{HI}}+bI_{{}_{}{}^{12}\mathrm{CO}(10)}+S_{100,\mathrm{BG}},$$
(3)
where $`S_{100,\mathrm{BG}}`$ is the flux density of the background emission, $`N_{\mathrm{HI}}`$ the HI column density, and $`I_{{}_{}{}^{12}\mathrm{CO}(10)}`$ the integrated <sup>12</sup>CO(1$``$0) line flux. The authors determine the constants $`a`$ and $`b`$ as $`a=(1.0\pm 0.4)\times 10^{20}`$ MJy sr<sup>-1</sup> cm<sup>2</sup> and $`b=(1.0\pm 0.5)`$ MJy sr<sup>-1</sup> K<sup>-1</sup> km<sup>-1</sup> s. Here we assume that IRAS point source data do not have to be corrected for significant contributions of any background emission. In order to obtain a lower limit to the FIR contribution of the disk to the overall FIR luminosity we just calculate the contribution expected from the atomic HI gas which is mostly distributed throughout the disk. The disk diameter of about 24.5 kpc and the adopted HI mass of $`M_{\mathrm{HI}}`$ = 3.2$`\times `$10<sup>9</sup> $`M_{\mathrm{}}`$ result in a 100 $`\mu `$m flux density contribution of about 8.6 MJy sr<sup>-1</sup>. Integrated over the disk this gives a flux density of $`S_{100}`$ = 1.26 Jy. For a total molecular gas mass of 3.2$`\times `$10<sup>9</sup> $`M_{\mathrm{}}`$ we find a 100 $`\mu `$m flux density contribution of the order of 0.3 mJy. The total disk flux density at 100 $`\mu `$m thus amounts to 1.56 Jy. With a mean ratio between the 60 and 100 $`\mu `$m flux density contribution for cirrus clouds of log($`S_{60}/S_{100})=0.65`$ (Helou Helou (1986)) the expected 60 $`\mu `$m flux density from the disk is $`S_{60}`$ = 0.35 Jy. These disk values can now be used to derive the nuclear FIR flux densities and luminosity (Tab. 5 and Tab. 6). Here we assume that the contribution of the southern compact component to the FIR luminosity is negligible because of the small filling factor of the source in the IRAS beam.
Although the presented decomposition of the FIR flux densities is very indirect, the FIR luminosities of the nucleus and disk are consistent with what one expects from the relation between the radio continuum and the FIR luminosity (Wunderlich & Klein, Wunderlich (1988)).
#### 3.1.6 The supernova rate
Oly and Israel (Oly (1993)) measured a 327 MHz flux density of UGC 1347 of $`S_{327\mathrm{MHz}}`$ = 23.6 mJy in a 55<sup>′′</sup> beam. The difference between peak and integrated flux density is only $`+`$0.5 mJy. The disk size of UGC 1347 is of the order of 90<sup>′′</sup> diameter. If the radio flux density were be dominated by the central 10% (30%) of the disk this would result in an approximately 3% (30%) deviation between the two quantities. In addition the source shows no clear indications for extended emission in the NRAO VLA Sky survey at 1.4 GHz (Condon et al. Condon1996 (1996)) with a beam size of 45<sup>′′</sup>. Therefore we assume that almost all of the flux density can be attributed to the nuclear region and that less than 1/20 of the radio emission originates in the disk of UGC 1347. Using a mean spectral index ($`S\nu ^\alpha `$) between 327 MHz and 1420/5000 MHz of $`0.71\pm 0.05`$, which the authors obtained for a sample of 35 UGC galaxies, we estimate a 5 GHz flux density of $`S_{5\mathrm{GHz}}`$ = 3.41 mJy. This value can be used to calculate the supernova rate $`\nu _{\mathrm{SN}}`$ via
$$\nu _{\mathrm{SN}}[\mathrm{yr}^1]=3.1\times 10^6\times S_{5\mathrm{G}\mathrm{H}\mathrm{z}}[\mathrm{mJy}]\times D[\mathrm{Mpc}]^2$$
(4)
(Condon Condon1992 (1992)). For UGC 1347 we find a supernova rate of $`\nu _{\mathrm{SN}}`$ = 0.044 yr<sup>-1</sup>. This value for the nuclear region of UGC 1347 is of the same order as the estimated overall supernova rate in the Milky Way of (0.025 $`\pm `$ 0.006) yr<sup>-1</sup> (Tammann et al. Tammann (1994)). For the disk of UGC 1347 we adopt an upper limit of $`\nu _{\mathrm{SN}}`$ = 0.002 yr<sup>-1</sup>. This determination of the SNR assumes that there is no major contribution to the radio flux density by an active nucleus. In all star burst analyzes were such a contribution could be excluded or was unlikely the derived SNR have shown to be consistent with other measurements of the star formation activity in the framework of the star burst model calculations.
### 3.2 UGC 1344
UGC 1344 is a SBc galaxy with an inclination of about $`i=60^{}`$ (data in UGC catalogue and visual inspection) located at R.A.(2000) = 01<sup>h</sup>52<sup>m</sup>34.8<sup>s</sup> and DEC(2000) = 363002<sup>′′</sup> approximately 21 north of the center of Abell 262. Like UGC 1347, it is well within the inner region in which the largest amount of HI deficiency is observed. A bright field star (GSC 2319-0343, $`m_\mathrm{V}`$ = 11.0 mag) is located about 23<sup>′′</sup> to the south. In Tab. 5 we list all parameters and estimates for UGC 1344 that were derived in a similar way as described above for UGC 1347.
Fig. 6 shows the NIR continuum emission from UGC 1344 together with the digitized sky survey V-band image in Fig. 7. We have used the image of the nearby reference star as the PSF and deconvolved the NIR continuum image of UGC 1344 with a Lucy-Richardson algorithm. In Fig. 8 we show intensity profile cuts through the nuclear component of UGC 1344 and the reference star, for both open and closed loop as well as deconvolved.
The K-band flux density distribution of UGC 1344 is very centrally peaked but smoothly connects to the extended disk. The central 3<sup>′′</sup> of the bulge contain about half of the K-band flux density in a 10<sup>′′</sup> aperture. For the nuclear component we measured a K-band flux density of 12 mJy in a 7.2<sup>′′</sup> aperture. This is about a factor of 1.7 more than for UGC 1347. As for UGC 1347 the K-band disk luminosity can be estimated from the deep H-band image presented by Gavazzi et al. (Gavazzi1996b (1996)) as $`m_\mathrm{H}`$ = 9.81 mag. For a mean H $``$ K color of 0.21 (Frogel et al. Frogel (1978)) this results in a total K-band flux of the order of 91 mJy. Correcting for the contribution of the nucleus we obtain a K-band flux and luminosity for the disk of UGC 1344 of about 81 mJy and $`L_\mathrm{K}`$ = 3.7$`\times `$10<sup>9</sup> $`L_{\mathrm{}}`$ .
Amram et al. (Amram (1994)) did not detect UGC 1344 in H$`\alpha `$ using a similar integration time as for UGC 1347 (P. Amram and M. Marcelin 1998, private communication). Based on the data obtained for UGC 1347 we adopt here an upper limit for the Lyman continuum luminosity of 10<sup>8</sup> $`L_{\mathrm{}}`$ for the entire galaxy.
The HI content of UGC 1344 has been studied by Wilkerson (Wilkerson (1980)). A weak line has been detected at a systemic velocity of 4155 km s<sup>-1</sup> and a width of 103 km s<sup>-1</sup> (Wilkerson Wilkerson (1980)). Again we assume that the difference of approximately 550 km s<sup>-1</sup> between the cluster velocity of 4704 km s<sup>-1</sup> and the systemic velocity are due to the motion of the galaxy within the cluster and adopt the same distance of 63 Mpc as for UGC 1347. From the HI detection Wilkerson (Wilkerson (1980)) derived an upper limit of the HI mass of 3.3$`\times `$10<sup>8</sup> $`M_{\mathrm{}}`$ (scaled to the adopted distance of 63 Mpc).
Oly and Israel (Oly (1993)) measured the 327 MHz radio continuum flux density of UGC 1347. The difference between the integrated 327 MHz flux density of 6.2 mJy and the peak flux density of 4.99 mJy in a $`55^{\prime \prime }\times 68^{\prime \prime }`$ beam indicates that not all of the radio emission can be associated with the nuclear component but that at least 1 mJy is due to extended emission. If the radio emission had a Gaussian distribution the angular size would be of the order of about 30<sup>′′</sup> to 40<sup>′′</sup> suggesting that this possible extended emission is distributed over the entire disk of UGC 1344.
From an inspection of the IRAS all sky survey we estimate an upper limit of the flux density at a wavelength of 100 $`\mu `$m of 0.5 Jy. Assuming a FIR spectrum similar to that of UGC 1347 this results in an upper limit of the far-infrared luminosity of 4$`\times `$10<sup>9</sup> $`L_{\mathrm{}}`$ .
### 3.3 Near-infrared imaging of the sample
In addition to the adaptive optics observations of UGC 1344 and UGC 1347 we took seeing-limited images of additional 9 spiral galaxies (a total of 11) in the Abell 262 cluster and 15 spiral galaxies in the Abell 1367 cluster. They were selected according to their HI deficiency and separation from the cluster center as given by Giovanelli et al. (Giovanelli1982 (1982)) and listed in Tab. 7 and Tab. 8.
A sample of 6 galaxies in Abell 262 and 7 galaxies in Abell 1367 was selected from the central parts of the clusters. In both cases they have separations from the cluster center less than 0.55 Abell radii ($`r_A`$ = 1.75 for Abell 262 and $`r_A`$ = 1.40 for Abell 1367) and HI deficiencies (as defined in Giovanelli et al. Giovanelli1982 (1982)) ranging between 0.06 and $`>`$1.18. The only exception is UGC 1347 with a deficiency of $`0.07`$ which is HI rich for a galaxy within the Abell radius. A second sample of 5 galaxies in Abell 262 and 8 galaxies in Abell 1367 was selected from the outer cluster regions. Here the separations from the cluster center range from 1.2 to 5.2 Abell radii and the HI deficiencies range from 0.03 to 0.6. Our K-band images of these sources contain in almost all cases reference stars that allowed us to accurately estimate the seeing. If no star was contained in the image we used the stars in adjacent exposures as a reference. From radial averages centered on the galaxy nuclei and on the stars we extracted the sizes as FHWM values (as listed in Tab. 7 and Tab. 8) and deconvolved the measurements on the galaxies with the stellar data via quadratic subtraction assuming Gaussian flux distributions of the sources. A comparison to our AO results on UGC 1347 and UGC 1344 shows that this procedure gives reliable estimates of the bulge sizes. In Fig. 9 we plot the deconvolved nuclear source sizes against the distance from the cluster center. To combine the data sets we scaled the results according to the mean radial velocities of the two galaxy clusters and corrected the Abell 1367 data to the distance of the Abell 262 cluster. For the galaxies close to the cluster center we find a median nuclear FWHM and median deviation from that value of 0.77<sup>′′</sup> $`\pm `$ 0.07<sup>′′</sup>. For the outer sources the result is 1.10<sup>′′</sup> $`\pm `$ 0.18<sup>′′</sup>. The difference between the two median values is 2.6 times the mean of the two median deviations. A Kolmogorov-Smirnov test shows that the two distributions are different at the 85% level. This result provides evidence that the K-band flux density distribution of the galaxy bulges in the inner part of the cluster are systematically smaller than those in the outer part. This result is discussed in section 4.3.
## 4 Discussion
The data on UGC 1344 and UGC 1347 indicate ongoing or recent nuclear star formation. The implications from that provide a framework in which the properties of our Abell 262 and Abell 1367 sample can be explained as well.
### 4.1 The starburst model
To derive the properties of a starburst from the observed continuum and line intensities we have used the starburst code STARS. This model has been successfully applied to NGC 1808 (Krabbe et al. Krabbe (1994), Tacconi-Garman et al. Tacconi-Garman (1996)), NGC 7469 (Genzel et al. Genzel (1995)), NGC 6764 (Eckart et al. Eckart1996 (1996)) and NGC 7552 (Schinnerer et al. Schinnerer (1997)). A description of the model can be found in the appendices of Krabbe et al. (Krabbe (1994)) and Schinnerer et al. (Schinnerer (1997)). The model is similar to other stellar population synthesis models (Larson & Tinsley Larson (1978), Rieke et al. Rieke1980 (1980), Gehrz et al. Gehrz (1983), Mas-Hesse & Kunth Mas-Hesse (1991), Rieke et al. Rieke1993 (1993), Doyon et al. Doyon (1994)) and includes the most recent stellar evolution tracks (Schaerer et al. Schaerer (1993), Meynet et al. Meynet (1994)).
We assume power-law initial mass functions (IMFs) which vary as $`M^\alpha `$ between a lower and upper mass cut-off, $`M_l`$ = 1 $`M_{\mathrm{}}`$ and $`M_u`$ = 100 $`M_{\mathrm{}}`$ , with an index $`\alpha `$ = 2.35 (Leitherer Leitherer (1996), Salpeter et al. Salpeter (1955)). STARS has as output observable parameters such as the bolometric luminosity $`L_{\mathrm{bol}}`$, the K band luminosity $`L_\mathrm{K}`$, the Lyman continuum luminosity $`L_{\mathrm{Lyc}}`$ and the supernova rate $`\nu _{\mathrm{SN}}`$, as well as the diagnostic ratios between these quantities. The adopted values for $`L_{\mathrm{bol}}`$, $`L_{\mathrm{Lyc}}`$, $`L_\mathrm{K}`$, and $`\nu _{\mathrm{SN}}`$ have been derived from observed properties in the previous sections. All relevant quantities as well as the diagnostic ratios that can be calculated from them for the present analysis are listed in Tab. 5 and Tab. 6. All the ratios are measures of the time evolution and the shape of the IMF, with slightly different dependencies on $`\alpha `$ and $`M_u`$.
### 4.2 Nuclear star formation in UGC 1347 and UGC 1344
UGC 1347: The diagnostic ratios (Tab. 5 and 6) and the framework of the starburst model now allow us to discuss the star formation for the whole of UGC 1347 as well as for the nucleus, the southern component and the disk. In summary the overall disk data of UGC 1347 are consistent with a high age and constant star formation, whereas the data for the nucleus and the southern compact component indicate more recent or ongoing star formation activity.
Disk: The $`M_{\mathrm{tot}}`$/ $`L_\mathrm{K}`$ ratio of 22 to 43 that we obtained for the overall galaxy indicates an age of the stellar population in the disk of the order of several 10<sup>9</sup> to 10<sup>10</sup> years (see Fig. 10).
This is also supported by the ratios $`L_{\mathrm{bol}}/L_{\mathrm{Lyc}}`$ = 6.4 and $`L_\mathrm{K}/L_{\mathrm{Lyc}}`$ = 3.2 combined with 10$`{}_{}{}^{9}\nu _{\mathrm{SN}}^{}/L_{\mathrm{Lyc}}<0.004`$ that results from the upper limit on the extended radio flux.
Nucleus: The diagnostic ratios $`L_{\mathrm{bol}}/L_{\mathrm{Lyc}}`$ = 53, $`L_\mathrm{K}/L_{\mathrm{Lyc}}`$ = 2.3 combined with 10$`{}_{}{}^{9}\nu _{\mathrm{SN}}^{}/L_{\mathrm{Lyc}}`$ = 0.27 indicate that the nuclear light is currently dominated by a decaying star formation event that happened a few times 10<sup>7</sup> years ago. This assumes that the upper mass cutoff is 100 $`M_{\mathrm{}}`$ and that the star burst itself did not last longer than a few times 10<sup>6</sup> years (see Fig. 10). The initial star formation rate was 18 $`M_{\mathrm{}}`$ yr<sup>-1</sup> and the current star formation rate would then be of the order of 0.04 $`M_{\mathrm{}}`$ yr<sup>-1</sup>.
The southern component: The southern component and the nucleus belong to the brightest regions in H$`\alpha `$ line emission. The only diagnostic ratio that we can calculate is that between the K-band luminosity corrected for the contribution of hot dust and the Lyman continuum luminosity. We obtain a value of $`L_\mathrm{K}/L_{\mathrm{Lyc}}`$ = 4. In addition the compact source of hot dust revealed by our AO measurements and its location at the tip of the bar suggest that the star formation activity in that region may be as high and recent as we find it for the nucleus (see Fig. 10).
UGC 1344: A ratio of $`M_{\mathrm{tot}}/L_\mathrm{K}`$ = 2.2 to 3.3 indicates an age of 10<sup>9</sup> years (see Fig. 10). This is supported by the lower limit $`L_\mathrm{K}/L_{\mathrm{Lyc}}37`$. The fact that the far-infrared luminosity is lower than that of UGC 1347 may imply that the molecular gas mass is below 10<sup>8</sup> $`M_{\mathrm{}}`$ . In addition the the system is deficient in neutral hydrogen as well. This indicates that UGC 1344 is deficient overall in fuel for star formation. The fact that UGC 1344 shows only a weak and narrow HI line may indicate that the HI has largely been stripped and the HI line width cannot necessarily be taken as a measure for the total dynamical mass. If the dynamical mass has been underestimated by a factor of 10 the log($`M_{\mathrm{tot}}/L_\mathrm{K}`$) ratio will be similar to the one obtained for UGC 1347 and the age of the dominant stellar population in the disk is then most likely also of the order of several 10<sup>9</sup> to 10<sup>10</sup> years.
### 4.3 Nuclear star formation and bulge sizes in the sample
The detailed discussion of the data available for UGC 1344 and UGC 1347 has shown that there is strong evidence for nuclear star formation activity and that not all of the K-band nuclear flux density can be explained by the presence of an old stellar population alone. Additional K-band flux density may originate from hot dust (see Fig. 5) or a population of supergiants or AGB stars, both of which are indicative for recent or ongoing star formation activity. The spatial distribution of the sources responsible for additional K-band flux density may be looked upon as independent of the distribution of the old stellar population forming the bulge of the galaxies. The relative flux density contribution of the additional sources will then have an influence on the measured size of the bulge.
That star formation activity is an important quantity for the appearance and classification of a galaxy has also been pointed out by Kennicutt et al. (Kennicutt (1994)). They have combined H$`\alpha `$ and UBV measurements of 210 nearby Sa-Irr galaxies with new photometric synthesis models to reanalyze past and future star formation timescales in the disks. The authors find that the pronounced change in the photometric properties of spiral galaxies along the Hubble sequence is predominantly due to changes in the star formation histories of disks, and only secondarily to changes in the bulge/disk ratio.
It is well known that there is a strong morphological segregation in clusters of galaxies, with most of the ellipticals in the center of the cluster while the spirals are more dispersed. Also early-type spirals seem more concentrated than late-types. Although Abell 262 is an extremly spiral rich cluster this segregation will most likely also effect the bulge to disk ratio as a function of location in the cluster. It is, however, not self-evident that this effect has any influence on the measured bulge sizes. In the following we will address this problem.
Our finding can be discussed in the framework of recent investigations of the bulge-to-disk luminosity ratio. In a sample of 3114 galaxies Solanes et al. (Solanes (1989)) analyzed the luminosity of bulge and disk components of disk galaxies and their possible correlations with morphological type and local density. Independently of the local environment no evidence is found for any bulge segregation among disk galaxies. Instead they find that disks appear to be less luminous with increasing local density. They find that the absolute brightness difference $`M`$(bulge)$`M`$(total) corresponds to about 3 for Sc galaxies, 2 for Sb and 1 for Sa and S0 galaxies. A similar trend is also observed in the dependency of the near-infrared concentration index $`C_{31}`$ (defined in Gavazzi et al. Gavazzi1990 (1990)) on the absolute H-band magnitude in a sample of 297 galaxies investigated by Gavazzi et al. (Gavazzi1996b (1996)). In this sample the bulge-to-disk ratio systematically increases with decreasing H-band luminosity. The morphological segregation reported by Solanes et al. (Solanes (1989)) is (if at all) only weakly indicated.
Both quantities, $`M`$(bulge)$`M`$(total) and $`C_{31}`$, are concentration parameters that simply describe the bulge versus disk brightness. In a scenario in which bulge components of identical brightness are located in the centers of disks of varying brightness the FWHM of the NIR light, measured with respect to the combined peak of the disk and bulge component will result in smaller values for lower disk luminosities. Therefore the observed disk luminosity antisegregation could in principle be responsible for our observed tendency that the FWHM of galaxies within the Abell radii of the Abell 262 and Abell 1367 clusters are smaller than those outside the Abell radii.
However, it is also possible that a luminosity antisegregation as observed for the disks is also present for the bulges, but is just compensated for by the additional contribution to the bulge luminosity due to enhanced star formation triggered by the effects of a higher cluster density environment.
Theoretical studies (Moore et al. Moore (1996), Fujita Fujita (1998)) show that the velocity perturbation induced by a single high-speed encounter is in most cases too small to affect the star formation rate of a disk galaxy. However, several successive high-speed encounters between galaxies (galaxy harassment) may lead to gas inflow and strong star formation activity (Fujita Fujita (1998)). This picture is consistent with the cluster crossing and star formation time scales. From the sample of 84 Abell 262 cluster members listed in Giovanelli & Haynes (Giovanelli1985 (1985)) we derive a velocity dispersion of 750 km s<sup>-1</sup>. The Abell radius of 1.75 then indicates that several 10<sup>9</sup> years are required to cross a significant fraction of the central part of the cluster. However, the time scale for formation of supergiants plus their life time amounts to only several 10<sup>7</sup> years. This indicates that in gas rich spiral galaxies star formation can easily be triggered via galaxy-galaxy interactions while passing through the central part of the cluster.
### 4.4 Star formation in the high redshift clusters
In Fig. 11 we show a K-band image of the J1836.3CR field and in Fig. 12 a K-band image of the PKS 0743-006 field.
Model values for the JHK flux densities were calculated using the GISSEL stellar population models (Bruzual & Charlot Bruzual (1993)), and compared with the measurements for each field and appropriate redshift. For better comparison we only show the tracks in Fig. 13, and then superimpose the observed data in Fig. 14 and Fig. 15.
All the models have been calculated for a passively evolving population after a 1 Gyr starburst. The Padova initial spectral energy distribution is used with a Salpeter IMF. Continued star formation, renewed star formation or initial mass functions truncated at their high mass end all tend to move galaxies towards the starts of the tracks. Since the SHARP II+ data were taken in the K filter we transformed the K magnitudes to K band values using the following relation: $`\mathrm{K}=\mathrm{K}^{}+0.2\times (\mathrm{H}\mathrm{K})`$ (Wainscoat and Cowie Wainscoat (1992)). For the color correction we used the mean observed H $``$ K colors of those objects which we could identify as being extragalactic and close to the redshift of the corresponding cluster. For the PKS 0743-006 cluster these mean colors are H $``$ K $`0.55`$ and for the J1836.3CR cluster we find H $``$ K $`0.69`$. In Tab. 9 and Tab. 10 we give the magnitudes and flux densities and JHK colors for the objects in our field of view towards both clusters.
In all cases the flux densities were derived from sky-subtracted images, taking into account possible contaminations by neighboring sources. The mean apertures for J1836.3CR and PKS 0743 were 3.95<sup>′′</sup>(K), 3.80<sup>′′</sup>(H), 3.62<sup>′′</sup>(J), and 4.13<sup>′′</sup>(K), 3.60<sup>′′</sup>(H), 3.33<sup>′′</sup>(J), respectively.
J1836.3CR: In Fig. 14 we show the colors of all sources in our field of view towards the J1836.3CR cluster at $`z=0.414`$. All galaxies are located at the tip of the evolutionary track for a passively evolving galaxy at that redshift. This is consistent with an age of 10 Gyrs (or more) and some intrinsic reddening with a mean value of the order of $`A_\mathrm{V}`$ = 2–3 mag.
Source $`a`$ is the brightest galaxy in our field of view of the J1836.3CR cluster. For this source our data allow us for the first time to obtain reliable color information close to the diffraction limit of the 3.6 m telescope. For this purpose the three 0.10<sup>′′</sup> pixel scale maps were first deconvolved with the PSF and then reconvolved with a Gaussian clean beam to be at the same resolution. The final resolution of the H and K maps is 0.2<sup>′′</sup> and that of the J map is 0.3<sup>′′</sup>. From the flux density calibrated maps we therefore calculated a J $``$ H color map at a resolution of 0.3<sup>′′</sup> and a H$``$K color map at a resolution of 0.2<sup>′′</sup>. As a result over a 1.9<sup>′′</sup> diameter aperture centered on source $`a`$ the colors turned out to be fairly uniform with mean colors corresponding to the values given in Tab. 9 and color variations of $`\pm `$ 0.15. The nucleus of source $`a`$ which can be clearly distinguished in the individual images at the three wavelengths does not appear particularly red or blue. This indicates that there are little variations in extinction or in spectral type across the source.
PKS 0743-006: In Fig. 15 we show the colors of all sources in our field of view towards the quasar PKS 0743-006 at $`z=0.994`$. Compared to the cluster J1836.3CR this field is more sparsely populated. Clearly not all of the objects have colors corresponding to high-redshift galaxies. Two of the objects ($`e,c`$) are at a location in the JHK color-color diagram that is populated by local giant or dwarf stars. For comparison – and as a convenient check of our calibration – we have also added the position of the bright reference star on which the adaptive optics loop was locked. This star is of type A0, in agreement with the measured colors. All other objects are located close to the middle of the evolutionary track for a passively evolving galaxy at a redshift of $`z=0.994`$. This is consistent with an age closer to 2 Gyrs than 10 Gyrs. The age may even be lower under the assumption of intrinsic reddening. This would indicate that the light of these objects is dominated by a reasonably young, blue stellar population.
Although the identification of these objects as $`z=0.994`$ galaxies has to be confirmed spectroscopically, an alternative explanation for the nature of these object is difficult to find. Even taking the uncertainties of the measured colors into account these sources are positioned well to the lower right of the mean colors of spiral galaxies or those of local dwarfs and giant stars. The only other objects that are located in this area of the JHK-color diagram are local HII regions, if their near-infrared emission is un-reddened and dominated by free-free radiation. Having 3 of these objects in the same field as the quasar (with similar colors) is very unlikely.
The best information on the sizes of the sources in this field is from the data taken with a 0.05<sup>′′</sup>/pixel sampling. While the quasar itself is unresolved as compared to radial and tangential cuts with respect to the direction of the AO reference, we find one of the brighter sources ($`b`$) that is close to the quasar as clearly extended in all directions. This source has an angular separation to the quasar of only 3.8<sup>′′</sup>. Deconvolving its measured size with the size of the unresolved quasar by subtracting the values in quadrature its deconvolved source size is $`0.22^{\prime \prime }\times 0.16^{\prime \prime }`$ at a position angle of about 45. At a redshift of $`z=0.994`$ this corresponds to a linear size of 1.9 kpc $`\times `$ 1.4 kpc. The extent of the source, its small angular separation from the quasar, and its colors make this object the best candidate for a cluster galaxy which is associated with the quasar PKS 0743-006. The previously described blue colors with respect to the values expected for a passively evolving galaxy at a redshift of $`z=0.994`$ then imply that we look at a galaxy with a region of enhanced star formation that extends over an area of 2.7 kpc<sup>2</sup>.
### 4.5 Bulge sizes and structures in J1836.3CR
Our best K-band AO images of J1836.3CR can be used to derive structural information on individual cluster members. We determined radial profiles in the direct and deconvolved images. The radial profiles of the deconvolved images have the advantage that they have been corrected for the non-diffraction limited part of the PSF and allow a much clearer view on the detailed distribution of light in each object. In Fig. 16 we show the profiles of 8 sources in the field including a profile of the guide star after the same number of Lucy deconvolution iterations ($`1000`$).
To determine the profile of source $`b`$ we first subtracted a radially averaged image of source $`a`$ as determined from its non-contaminated section to the north. This process revealed a further, even weaker source $`b2`$ (see Fig. 17). Both sources $`b`$ and $`b2`$ were subtracted from the image before calculating the average radial profile of $`a`$.
Following Kormendy (Kormendy (1977)) the radial profiles were fitted via a reduced $`\chi ^2`$ test as a combination of a $`r^{1/4}`$ spheroid (De Vaucouleurs De Vaucouleurs (1959)) with a surface brightness
$$\mu _\mathrm{V}=\mu _\mathrm{V}(0)+8.325[(r/r_0)^{1/4}1]$$
(5)
and an exponential intensity fall-off with a sharp inner cutoff radius (Kormendy Kormendy (1977)). with a surface brightness
$$\mu _\mathrm{E}=\mu _\mathrm{E}(0)+1.0857[\alpha r+(\beta /r)^3].$$
(6)
Here $`r_0`$ and $`\beta `$ are the corresponding cutoff radii, $`1/\alpha `$ is the exponential scale length and $`\mu _\mathrm{V}(0)`$ and $`\mu _\mathrm{E}(0)`$ are the central surface brightnesses for the $`r^{1/4}`$ spheroid and the exponential disk (if it had not been cut-off) in units of magnitudes per square arcseconds. The surface brightness $`\mu (r)`$ and intensity distribution $`I(r)`$ are linked via $`\mu (r)=2.5\times \mathrm{log}I(r)`$. None of the sources, with the exception of $`a`$, can be fit via an $`r^{1/4}`$ law alone. The distributions at small radii are too flat with respect to their decrease at larger radii; they can only at larger radii be approximated with a De Vaucouleurs law. This indicates that these objects are not typical elliptical galaxies which are usually found amongst the brightest cluster members (Thuan and Romanishin Thuan (1981)). All sources – except object $`a`$ – are best fit with exponential profiles suggesting that they are more similar to spiral disks. However, no central bulge could be detected with an upper size limit of 0.6–0.9 kpc. This could be due to the fact that the angular resolution of 0.15<sup>′′</sup> is still not sufficient to clearly discriminate between the bulge and disk components. Alternatively it could imply that the bulges are intrinsically weak with respect to the disks. The statistical investigation of bulge to disk intensity ratios by Solanes et al. (Solanes (1989)) would then suggest that the galaxies in the center of J1836.3CR are of type Sa and S0 and that the center of this cluster represents a dense environment.
Source $`a`$ shows a radial profile of a typical cD galaxy that is expected to be present in the centers of rich galaxy clusters. The cD galaxies in centers of clusters with low richness can usually be described as giant ellipticals (Morgan et al. Morgan (1975), Albert et al. Albert (1977)). However, cD galaxies in rich clusters show excess flux density in their profiles at large radii with respect to an $`r^{1/4}`$ law (Oemler Oemler (1997), Dressler Dressler (1979)). This is also the case for source $`a`$. The profile can only be fit by a De Vaucouleurs law for radii $`<0.6^{\prime \prime }`$. The resulting scale length (see Tab. 11) is in agreement with the lower limit of the range of scale lengths of $`r^{1/4}`$ laws found for the central regions of bright cD galaxies in rich Abell clusters (Oemler Oemler (1997)). The extended flux at larger radii of the central cD galaxy in J1836.3CR may be indicative of either tidal disturbances due to galaxy-galaxy interactions of cluster members in which the debris is sinking onto the central massive galaxy (Gallagher & Ostriker Gallagher (1972)) or even the result of recent merger processes (Ostriker & Tremaine Ostriker (1975)). These activities are not likely to be important in clusters of low richness (Thuan and Romanishin Thuan (1981)) but can very well be responsible for the structures of objects in rich clusters.
### 4.6 Luminosity of the cluster members
In order to check how representative the individual galaxies that we studied are we compared their luminosities to those of $`L_{}`$ galaxies at the corresponding redshifts of the three clusters at $`z=0.0157`$, $`z=0.414`$, and $`z=0.994`$. From a subsample of the optically selected Anglo-Australian Redshift-Survey, Mobasher et al. (Mobasher (1993)) find the best infrared Schechter luminosity function parameters as $`M_\mathrm{K}^{}=25.1\pm 0.3`$ with $`\alpha =1.0\pm 0.3`$ for $`H_0=50`$ km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_0=0.02`$. From their sample they also find that E/S0 and spirals have identical infrared luminosity functions within the errors. Their parameters are very much in agreement with recent determinations of $`M_\mathrm{K}^{}`$ by De Propris et al. (De Propris (1998)) and Gardner et al. (Gardner (1997)). For $`h=H/H_0=0.5`$ and $`q_0`$ = 0.5 they find $`M_\mathrm{K}^{}=24.8`$ with bright galaxy slopes of $`\alpha =0.78`$ and $`M_\mathrm{K}^{}=24.6\pm `$ 0.1 with $`\alpha =0.91\pm `$ 0.1, respectively. For the purpose of comparison we use the same values for $`H_0`$ and $`q_0`$ and obtain the following results for our targets:
Abell 262: For our low redshift cluster both objects we study here – UGC 1344 and UGC 1347 – are $`L_{}`$ galaxies to within less than 0.5 magnitudes.
J1836.3CR: For this cluster we find that the bright source $`a`$ is just 0.6 magnitude brighter and that the other sources are about 1 to 2 magnitudes fainter than a typical $`L_{}`$ galaxy at that redshift of $`z=0.414`$.
PKS 0743-006: For the 3 presumed galaxies close to the quasar PKS 0743-006 at $`z=0.994`$ we find that two are about 0.7 magnitudes brighter than an $`L_{}`$ galaxy and the weakest one is 2 magnitudes fainter than $`L_{}`$. This is consistent with the statement that the brighter members of the observed clusters are on the K-$`z`$ relation which, as emphasized by Lilly (Lilly (1988)), shows only a small scatter of $`\sigma 0.3`$ mag at redshifts below $`z=1.5`$. In addition it should be mentioned that the quasar PKS 0743-006 itself shows very similar colors compared to the galaxy candidates in the observed field. However, due to the nuclear contribution it is – as expected (see data in Dunlop et al. Dunlop (1993) and Lehnert et al. Lehnert (1992)) – between 1.5 and 2 mag brighter than what is predicted by the galaxy K-$`z`$ relation. We also find that the measured JHK fluxes of this variable quasar were at the time we took the data about 0.7 to 1.0 magnitudes below what is given in the literature (White et al. White (1988), Lépine et al. Lepine (1985)). Since quasar host galaxies have been found to be often as bright as $`5\times L_{}`$ (McLeod & Rieke McLeod1995 (1995), Sánchez et al. Sanchez (1997)) the underlying host galaxy may contribute substantially to the overall flux density of the quasar in this low state.
This indicates that the source properties that we determined in the previous sections are representative for typical $`L_{}`$ cluster members at the corresponding redshifts.
## 5 Summary and Conclusions
We have presented high angular resolution NIR observations of three galaxy clusters at different redshifts using adaptive optics. In the case of the barred spiral UGC 1347 in Abell 262 we presented the first adaptive optics data using the laser guide star provided by the ALFA system.
The diagnostic ratios for the nucleus of UCG 1347 indicate recent and ongoing star formation activity. In addition to the resolved NIR nucleus in UGC 1347 we found a bright and compact region of recent and enhanced star formation at one tip of the bar. The $`L_\mathrm{K}/L_{\mathrm{Lyc}}`$ ratio as well as the V $``$ K color of that region imply that a starburst happened about 10<sup>7</sup> years ago. For UGC 1344 we found that the overall star formation activity is low and that the system is deficient in the fuel for star formation.
The comparison of seeing-corrected nuclear bulge sizes of a sample of 26 cluster galaxies within and outside of the Abell radius of Abell 262 and Abell 1367 indicates that the galaxies in the inner part of the cluster show a tendency for more compact bulges than those outside. This phenomenon can tentatively be ascribed to an increased star formation activity due to interactions of cluster members inside the Abell radius. Such an increase of central activity is also indicated at other wavelengths. Scodeggio and Gavazzi (Scodeggio (1993)) find in a 21 cm survey of spiral galaxies in clusters that about 30% of them show extended radio continuum emission and that a substantial fraction of those (but not all) show indications of interaction. Moss and Whittle (Moss1993 (1993)) find from an H$`\alpha `$ survey of cluster spirals that interacting spirals show a strong tendency to have compact nuclear H$`\alpha `$ emission which the authors conclude to be most likely due to tidally induced star formation from galaxy-galaxy interactions, since interactions are more likely to happen close to the cluster center. Several successive high-speed encounters between galaxies may lead to gas inflow and strong star formation activity (Fujita Fujita (1998)). This would imply that dynamically induced star formation is more important in the center than the outer parts of a cluster, although current investigations have not yet convincingly shown the obvious presence of such an correlation. Future observations of larger samples of cluster members are clearly needed to substantiate these correlations.
Since the spiral content of galaxy clusters at higher redshifts is about as large as the spiral content in the field at $`z=0`$ (Oemler et al. Oemler (1997)), detailed observations of galaxies in low-redshift spiral-rich clusters may provide essential information of the cluster evolution at higher redshifts in general. In particular it would allow us to study in detail the influence of galaxy harassment. Abell 262 and Abell 1367 are spiral-rich clusters. With a ratio of spirals to S0 and E0 galaxies of 47%/53% and 43%/57%, respectively, their spiral content is similar to that in the field at $`z=0`$ of 55%/45% and to that in clusters at a redshift of $`z=0.4`$ of 40%/60% (Oemler Oemler (1997)).
UGC 1347 and UGC 1344 could very well be taken as examples of the blue and red fraction of clusters at higher redshift. Couch et al. (Couch (1998)) reports on 3 galaxy clusters at $`z=0.3`$ measured with the HST. In these clusters he finds the fraction of spirals at least 3 times higher than at $`z=0`$ which would approximately correspond to the spiral fraction of Abell 262 and Abell 1367. About 20% of all galaxies show signs of interaction. The blue fraction of the cluster population shows morphologies similar to Sb-Sdm/Irr galaxies with compact knotty regions of star formation. These knots may be very similar to the bright star formation region we found for UGC 1347. The red fraction of the cluster population they find is 1-2 Gyr past the last major star formation event and has morphological similarities to S0-Sb disks. This may be very similar to UGC 1344 and source $`a`$ in J1836.3CR.
However, the radial profile of source $`a`$ in J1836.3CR shows indications for recent or ongoing interactions between cluster members. Here the enhanced flux above the De Vaucouleurs fit to the data at larger radii identifies source $`a`$ as the central cD galaxy in a rich cluster environment in which interactions between cluster members are probably still of importance for their further evolution.
From an investigation of NIR colors Hutchings & Neff (Hutchings (1997)) find for 5 quasars at redshifts ranging from $`z=0.06`$ to 0.3 that they are located in mostly evolved groups of galaxies with an indication for average extinction values of $`A_\mathrm{V}`$ = 2–3 mag quite similar to J1836.3CR at $`z=0.414`$ discussed in this paper. However, on the basis of large-scale investigations (Butcher & Oemler, Butcher (1984), Ellis Ellis (1997)) as well as studies of individual galaxy clusters (e.g. Morris et al. Morris (1998)) one finds that both in clusters and in the field the blue fraction of galaxies generally increases towards higher redshifts (the Butcher Oemler effect, BOE). This is consistent with our finding of relatively blue NIR colors in the galaxy candidates associated with PKS 0743-006 at $`z=0.994`$.
Based on an investigation of Abell 2390 at $`z=0.2279`$ and a comparison to the cluster MS 1621.5+2640 at $`z=0.4274`$ Abraham et al. (Abraham (1996)) and Morris et al. (Morris (1998)) suggest that the star formation process is shut down by a combination of gas stripping followed by gas exhaustion via star formation. This truncation of star formation activity may explain both the BOE and the large fraction of S0 galaxies in clusters. This also suggests that truncated star formation induced by infall does not play a major role in driving cluster galaxy evolution at lower redshifts although this mechanism may have played a role in earlier history. The barred structure of UGC 1347 as well as the strong recent star formation event at one tip of the bar may indicate that this object is well suited to study the corresponding physical processes of cluster galaxy evolution in great detail.
###### Acknowledgements.
We are grateful to the Calar Alto and La Silla staff for their excellent support and hospitality. We thank M. Lehnert, H.W. Rix, R. Genzel, L. Tacconi and T. Boller for helpful discussions. We are particularly grateful to P. Amram and M. Marcelin who kindly provided the H$`\alpha `$ continuum and line images to us. We furthermore thank the referee for very helpful and constructive suggestions.
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# Foreground and background dust in star cluster directions
## 1 Introduction
Full-sky surveys in the far infrared have been achieved by means of the IRAS and COBE satellite observations. Schlegel et al. (1998) built a reddening map from the 100 $`\mu `$m IRAS dust emission distribution considering temperature effects using 100/240 $`\mu `$m DIRBE data. The transformation to E(B-V) maps was obtained from dust columns calibrated via the (B-V)-Mg2 relation for early type galaxies. This far-infrared reddening (hereafter E(B-V)<sub>FIR</sub>) presents a good agreement at high galactic latitudes with that derived from H I and galaxy counts by Burstein & Heiles (1978, 1982) with an offset of 0.02 mag (lower values for the latter method). Recently, Hudson (1999) analysed E(B-V)<sub>FIR</sub> maps using 50 globular clusters with $`|b|>10^{}`$ and distance from the plane $`|Z|>3`$ kpc, as well as 86 RR Lyrae from Burstein & Heiles (1978). These two samples provided slightly lower values on the average as compared to Schlegel et al.’s reddening values ($`\mathrm{\Delta }`$E(B-V) = -0.008 and -0.016, respectively). The reddening comparisons above hardly exceed the limit E(B-V) $``$ 0.30, so that a more extended range should be explored.
Since the Galaxy is essentially transparent at 100 $`\mu `$m, the far-infrared reddening values should represent dust columns integrated throughout the whole Galaxy in a given direction. Star clusters probing distances as far as possible throughout the Galaxy should be useful to study the dust distribution in a given line of sight. Globular clusters and old open clusters are ideal objects for such purposes because they are in general distant enough to provide a significant probe the galactic interstellar medium and have a suitable sky coverage. Clearly, star clusters beyond the disk dust layer are expected to have reddening values essentially comparable to those of galaxies in the same direction. On the other hand, clusters within the dust layer should have contributions from clouds in background regions. Another issue is the thickness of the Milky Way dust lane and whether some dust clouds occur at higher distances from the Plane. Recently, several edge-on spiral galaxies have been studied in detail (Howk & Savage 1999) and a comparison of their dust distribution with that of the Milky Way is worthwhile.
The aim of the present study is to compare star cluster reddening values measured from direct methods, i.e. sampling the dust effects seen in the light emitted by the cluster members, with those derived from the 100 $`\mu `$m dust emission. We investigate the possibility of background and foreground dust contributions in star clusters directions. In Sect.2 we present an overview of Schlegel et al.’s (1998) reddening values predicted in different environments in the Galaxy. In Sect.3 we gather the necessary data for globular clusters and old open clusters and describe the sample properties. In Sect.4 we discuss the results, especially the star cluster lines of sight with evidence of background dust. Finally, the concluding remarks are given in Sect.5.
## 2 Overview of dust emission reddening values E(B-V)<sub>FIR</sub>
For a better understanding of the reddening distribution throughout the Galaxy we extracted E(B-V)<sub>FIR</sub> values from Schlegel et al.’s maps using the software dust-getval provided by them. We discuss (i) directions along the galactic plane which accumulate reddening from sources in different arms and other large structures, and (ii) galactic latitude profiles to see the effects of relatively isolated nearby (high latitude) dust clouds.
We show in Fig.1 the entire Galaxy longitude profile. The upper panel is in direction of the galactic centre and the lower one is in direction of the anticentre. Note the enormous reddening differences between the two panels: the lower panel has typical values of E(B-V)$`{}_{\mathrm{FIR}}{}^{}`$1.5, and the values in the upper panel are a factor $``$ 10 higher. We indicate a series of H I, CO and optical features which help interpret the reddening distribution: (i) tangent regions of the spiral arms Sagittarius-Carina, Scutum (5 kpc arm) and 4 kpc arm (Henderson 1977, Georgelin & Georgelin 1970a, Cohen et al. 1980); (ii) the extent of the 3 kpc arm (Kerr & Hindman 1970, Bania 1980); (iii) the extent of the far side of the Sagittarius-Carina arm (Grabelsky et al. 1988); (iv) the Molecular Ring (MR) and the Central Molecular Zone (CMZ), (Combes 1991, Morris & Serabyn 1996); and finally, (v) the extent of the Local (Orion) and Perseus arms (Georgelin & Georgelin 1970b).
The relatively low reddening in the anticentre panel can be basically explained by the cumulative effect of the three external arms: Orion, Perseus and Outer arm (Digel et al. 1990). It is worth noting that E(B-V)<sub>FIR</sub> on the average is higher in the second quadrant than in the third quadrant, probably by the interruption of the Perseus arm. The steady increase of E(B-V)<sub>FIR</sub> in the first and fourth quadrants towards the direction of the Galactic center can be explained by the cumulative effect of inner arms and especially their tangent zones. Owing to the 100 $`\mu `$m dust emission transparency the far side arms of the Galaxy will also contribute to E(B-V)<sub>FIR</sub> (see the extent of far side of the Sagittarius-Carina arm in the fourth quadrant). The Molecular Ring is also a major contributor, leading to a plateau level E(B-V)$`{}_{\mathrm{FIR}}{}^{}`$ 20. Finally, the Central Molecular Zone is responsible for the central cusp.
Figure 2 shows E(B-V)<sub>FIR</sub> profiles in the interval $`25^{}b25^{}`$ for selected galactic longitudes including well-known dark cloud centers. The individual clouds, especially their central parts can attain comparable (in some cases higher) E(B-V)<sub>FIR</sub> values to disk zones at lower latitudes. Individual dark clouds have a core-halo structure. In the $`\rho `$ Oph dark cloud the core FWHM is 35’ while at E(B-V)<sub>FIR</sub> = 0.5 the halo diameter is 4. For the Chamaleon I complex the core FWHM is 48’ (in the region of the reflection nebula IC 2631) while the halo diameter at E(B-V)<sub>FIR</sub> = 0.5 is 2.9. From these simple comparisons and corresponding solid angles the zones responsible for the accumulation of reddening throughout the arms appear to be the cloud halos rather than the cores, possibly combined with diffuse galactic dust.
### 2.1 Comparison with reddening values from JHK photometry for nearby dark clouds
We show in Table 1 embedded infrared star clusters and T Tauri groups which are related to nearby dust complexes. They can be useful to analyse the high extinction regime and to compare Schlegel et al.’s reddening values with those measured directly from Infrared (JHK) photometry of the stellar content. We give E(B-V)<sub>FIR</sub> values for the central positions of these stellar aggregates and indicate to which complexes they belong. The identifications and locations of these objects are from: (i) Gomez & Lada (1998) for the T Tauri groups related to the dark clouds Barnard 30 and 35; (ii) Lada et al. (1991) for the infrared clusters embedded in the nebulae NGC 2071, M 78, NGC 2023 and NGC 2024 in the Orion Complex LDN 1630 Molecular Cloud; (iii) Minchin et al. (1991) , Jones et al. (1994) and Reipurth et al. (1999 and references therein) respectively for the three deeply embedded clusters OMC-1, OMC-2 and OMC-3 in Orion Complex Molecular clouds; (iv) Strom et al. (1993) for the 7 objects in the Orion complex LDN 1641 Molecular Cloud; (v) Carpenter et al. (1997) for the Mon R2 IR cluster; (vi) Lawson et al. (1996) for two concentrations of T Tauri stars around the reflection nebulae IC 2631 and Cederblad 110/111 in the Chamaleon I dark cloud; (vii) Comerón et al. (1993 and references therein) for the $`\rho `$ Ophiuchi IR cluster. The larger E(B-V)<sub>FIR</sub> values occur for IR star clusters, while the T Tauri groups tend to be associated with lower reddening values. This must reflect the need of higher dust densities for the formation of star clusters and massive stars, conditions which occur in the cores of Giant Molecular Clouds in contrast to less massive dark clouds (e.g. Comerón et al. 1993, Carpenter et al. 1997).
For the infrared photometric reddening comparisons with E(B-V)<sub>FIR</sub> we adopt a total to selective extinction ratio $`R=\frac{A_\mathrm{V}}{E(BV)}=3.1`$. When the original studies do not express their results in terms of $`A_V`$, we adopt the ratios $`\frac{A_\mathrm{J}}{A_\mathrm{V}}=0.276`$, $`\frac{A_\mathrm{H}}{A_\mathrm{V}}=0.176`$ and $`\frac{A_\mathrm{K}}{A_\mathrm{V}}=0.112`$ from Schlegel et al. (1998).
JHK photometry of embedded stars in NGC 2024 (Comerón et al. 1996) indicates an average E(B-V) $`14.5`$, lower than that given by E(B-V)<sub>FIR</sub> (Table 1) which might be accounted for if the sources are not the deepest embedded ones and if there still is a considerable amount of dust in the back half of the cloud. Another possibility is that hot stars heat the cloud core beyond the upper limit (21K) available in Schlegel et al.’s temperature maps.
HK photometry of deeply embedded stars in OMC-2 (Johnson et al. 1990) provided an average $`E(BV)=7.7`$, so that the reddening through the whole cloud should be twice as much. These values are consistent with E(B-V)$`{}_{\mathrm{FIR}}{}^{}=11.4`$ (Table 1).
JHKL photometry of a deeply embedded source in OMC-1 (Minchin et al. 1991) indicated A$`{}_{\mathrm{L}}{}^{}=6`$ and A$`{}_{\mathrm{K}}{}^{}>9`$ which implied A$`{}_{\mathrm{V}}{}^{}200`$ or 90 depending on models with large and small dust grain respectively. The latter values convert to E(B-V) = 64.5 or 29.0 respectively, which bracket that from E(B-V)$`{}_{\mathrm{FIR}}{}^{}=50.9`$ (Table 1), in this case favouring the large grain model.
JHK measurements of stars in the Mon R2 IR cluster (Carpenter et al. 1997) indicated that most of the embedded stars are in the range $`3.3E(BV)4.5`$ somewhat higher than E(B-V)$`{}_{\mathrm{FIR}}{}^{}=2.5`$.
McGregor et al. (1994) obtained JHK photometry of a star projected close to the center of the reflection nebula Ced110 in Chamaleon I. They derived E(B-V) $`2.4`$ or $`5.6`$ depending on the assumed spectral type. In this direction E(B-V)<sub>FIR</sub> = 3.25, in reasonable agreement.
Comerón et al. (1993) estimated from JHK photometry that deeply embedded sources in the $`\rho `$ Oph dark cloud core have $`E(BV)>16`$ and that field stars in the background of the cloud are possibly affected by E(B-V) $`22.5`$. These values bracket that of E(B-V)<sub>FIR</sub> (Table 1) which is sensitive to reddening arising from the cloud as a whole.
Finally, we compare reddening values for the old/young components of the galactic nucleus (Catchpole et al. 1990, Krabbe et al. 1991), and for the two young star clusters projected close to the nucleus which contain WR stars (Quintuplet cluster = AFGL2004 and the Arches cluster = Object 17, respectively Glass et al. 1990, Nagata et al. 1995). The galactic center reddening is E(B-V) $`9.7`$, that for the Quintuplet cluster E(B-V) $`7.1`$ and that for the Arches cluster E(B-V)$`10.6`$. The E(B-V)<sub>FIR</sub> values in these directions are exceedingly high (E(B-V)$`{}_{\mathrm{FIR}}{}^{}98100`$). High values are expected since in this direction there is an integrated dust contribution behind the nucleus out to the disk edge. In the foreground of the nucleus we have contributions from the central molecular zone, the molecular ring and the four arms between the sun and the galactic nucleus. By symmetry arguments considering that there are three extra arms in the background outside the solar circle Schlegel et al.’s reddening values should not exceed E(B-V) $`{}_{\mathrm{FIR}}{}^{}25`$. A possible explanation is that the dust temperature in the nucleus and surroundings is significantly higher than those in Schlegel et al.’s temperature map. Indeed the possibility of non-thermal photon flux and existence of three massive young clusters may account for the dust heating.
We conclude that Schlegel et al.’s reddening values are in general consistent with values obtained from infrared photometry of embedded sources in dark clouds. Some discrepancies for infrared star clusters in the cores of molecular clouds might be explained by dust heated a few degrees above 21K. As pointed out by Schlegel et al. (1998) the same dust column density provides a factor of 5 in 100 $`\mu `$m flux when heated from 17K to 21K which are the temperature extremes considered in their study. For the Central Molecular Zone in the Galaxy E(B-V) <sub>FIR</sub> values appear to be exceedingly high so that the dust temperature near the galactic nucleus must be higher. Indeed, assuming that E(B-V)$`{}_{\mathrm{FIR}}{}^{}25`$ the dust temperature implied is $`T26`$K.
### 2.2 Dust layer height from nearby dark clouds
In the following we calculate the distance from the galactic plane (Z) of the nearby dust complexes to estimate the dust layer height to which significant reddening is expected, at least in the solar neighbourhood.
Considering that the sun appears to be located 15 pc above the galactic plane (Cohen 1995, Hammersley et al. 1995), the corrected distance from the galactic plane $`Z^{}`$ is
$$Z^{}=d_{\mathrm{sun}}\mathrm{sin}b+15,$$
(1)
where d<sub>sun</sub> is the object distance from the sun in pc.
Hipparcos distances of the Orion and $`\rho `$ Ophiuchi complexes are respectively d<sub>sun</sub> = 490 pc and d<sub>sun</sub> = 125 pc (de Zeeuw et al. 1999), which imply distances from the plane $`Z^{}`$ = -148 pc and $`Z^{}`$ = 59 pc respectively. The distance to Monoceros R2 is d<sub>sun</sub> = 830 pc (Carpenter et al. 1997) and $`Z^{}`$ = -166 pc. Finally, for Chamaleon I d<sub>sun</sub> = 140 pc (Lawson et al. 1996) and $`Z^{}`$ = -21 pc. Of these well-studied high latitude clouds Mon R2 and the Orion complex are intrinsically distant enough from the plane to estimate the dust layer height. The centers of the complexes imply $`|Z^{}|150`$ pc. Since dust in the Orion complex appears to absorb significantly (E(B-V)$`{}_{\mathrm{FIR}}{}^{}`$ 0.20) to at least b = -25 – see in Fig. 2 the b profiles for OMC-1 and Mon R2 (which includes Orion complex parts at higher latitudes), we adopt a dust layer height of 200 pc for the subsequent discussions.
## 3 Reddening Values in star cluster directions
Since star clusters span a wide range of galactic latitudes and distances from the sun, both within and outside the dust layer, they are ideal targets for comparison of Schlegel et al.’s reddening values with those derived from the stellar content methods. In the present section we compile reddening values of globular clusters and intermediate age open clusters. Owing to their higher ages they probe the interstellar medium without being physically related to the dust complexes, except for the possibility of interactions.
### 3.1 Globular clusters
Harris (1996) compiled parameters for the 147 Milky Way globular clusters, keeping an updated version in the Web interface
http://physun.physics.mcmaster.ca/Globular.html.
In previous compilations, e.g. Webbink (1985), many globular clusters had scanty or no information. Colour-magnitude diagrams (CMD) based on CCD observations are now almost complete for these objects as a consequence of recent efforts, especially for the reddened low latitude globular clusters in crowded fields (e.g. Ortolani et al. 1995a, Barbuy et al. 1998a and references therein). Just to illustrate the progress achieved we mention Terzan 3 for which Webbink (1985) provided $`E(BV)=0.32`$ based on the cosecant law, but the CMD showed a considerably higher reddening $`E(BV)=0.72`$ (Barbuy et al. 1998b).
Table 2 lists the galactic globular cluster, as follows: (1) name of object, (2) and (3) galactic coordinates, (4) distance from the sun, (5) reddening derived from dust 100 $`\mu `$m emission E(B-V)<sub>FIR</sub>, (6) E(B-V) derived from the light emitted by the cluster members, and (7) $`\beta `$E(B-V) which is the difference between E(B-V)<sub>FIR</sub> and E(B-V) (Sect.4.1). The E(B-V)<sub>FIR</sub> values were obtained from Schlegel et al.’s reddening maps using the cluster galactic coordinates. Reddening and distance values are from Harris’ (1996) compilation as updated to June 22 1999, except for low latitude globular clusters which come from the CMD studies indicated in the Table notes. Low latitude globular clusters have also been studied in detail via integrated spectral distribution in the near IR ($`7000<\lambda <10000`$ Å), which also is a direct estimator of the reddening affecting the stellar content (Bica et al. 1998). For these clusters, reddening values derived spectroscopically were also considered in Table 2 (see Table notes).
### 3.2 Old open clusters
The old open clusters (700 Myr or older), also usually referred to as Intermediate Age Clusters (IACs), are particularly suitable for studying the galactic reddening at low and moderately high galactic latitude directions because the old disk is relatively thick (Friel 1995). They are numerous for 90$`{}_{}{}^{}<`$ $`\mathrm{}`$ $`<`$ 270, thus complementary to the globular cluster sample which in turn probes numerous lines of sight towards the bulge. We looked for old open clusters in compilations (Janes & Phelps 1994, Friel 1995, Carraro et al. 1998), and individual clusters in the Open Cluster Database (Mermilliod 1996) as updated to November 1999 in the Web interface
http://obswww.unige.ch/webda.
We checked in the original references the CMD quality and the derived cluster parameters. In recent years the number of CCD photometric studies has been increasing steadily. They include clusters with CMD for the first time, CCD data on clusters previously observed photographically, and finally clusters newly observed in the infrared (J and K bands). Just to mention some recent studies: NGC 2204, NGC 2477, Berkeley 39 and Melotte 66 (Kassis et al. 1997), Trumpler 5 (Kaluzny 1998), Pişmiş 18, Pişmiş 19, NGC 6005 and NGC 6253 (Piatti et al. 1998a), Berkeley 17 and Berkeley 18 (Carraro et al. 1999), and ESO93-SC08 (Bica et al. 1999).
Janes & Phelps’ (1994) compilation included 72 IACs while the present sample has 103 entries. The Hyades were not included due to the proximity and large angular size. Table 3 lists the galactic old open clusters, as follows: (1) name of object, (2) and (3) galactic coordinates, (4) distance from the sun, (5) age, (6) E(B-V)<sub>FIR</sub>, (7) E(B-V), and (8) $`\beta `$E(B-V).
Figure 3 shows the sample properties. The histogram giving the distribution of old open clusters as a function of the distance from the sun. It shows that the sample is probing the interstellar medium quite far, mostly in the range 1-5 kpc and in some cases as far as 10-14 kpc. The age histogram shows a steady decrease for older ages probably related to the dissolution rate of IACs. There occurs a peak at t $``$ 5 Gyr which was also present in previous compilations, and the present increased sample supports its significance. A possible interpretation for this peak would be a burst of star formation in the old disk.
## 4 Discussion
The angular distribution in galactic coordinates of globular and old open clusters for $`|b|<`$ 20 is shown in Fig.4, centred on the galactic nucleus direction. The two samples are complementary, globular clusters probe mostly the galactic central regions while the old open clusters probe mostly the anticentre regions. The most frequent values for $`|b|<`$ 10 are intermediate ones ($`0.2E(BV)1.0`$). For higher latitudes smaller values dominate.
Reddening values for globular clusters in bulge regions often exceed E(B-V) = 1. Terzan 1, 4, 5, 6 and 10 exceed E(B-V) = 2, while Liller 1 and UKS 1 have E(B-V)$``$ 3 (Table 2). The dust emission reddening can be much larger, in some cases exceeding E(B-V)<sub>FIR</sub> = 4 which occurs for Terzan 5, 10, 4, 1 and UKS 1. The lowest galactic latitude globular cluster Liller 1 has the highest value (E(B-V)<sub>FIR</sub> = 11.57).
For the old open clusters the largest reddening derived from the CMD occurs for Pişmiş 2 (E(B-V) = 1.48), and seven other clusters exceed E(B-V) = 1 (Table 3). Three clusters have dust emission reddening exceeding E(B-V)<sub>FIR</sub> = 5 which are IC 4291 (Pişmiş 18), Pişmiş 19 and finally NGC 6134 with E(B-V)<sub>FIR</sub> = 25.66. They are at extremely low latitudes and in directions not far from the galactic centre (Table 3) which can accumulate reddening from dust in several arms and the Molecular Ring (Sect.2).
Figure 5 shows E(B-V)<sub>FIR</sub> as a function of the reddening derived from the stellar content. Panel (a) contains the clusters with $`|b|>20^{}`$, presenting a good agreement between the reddening values, except the clusters M 107 and NGC 1901 which are discussed in Sect.4.1. The highest values at such latitudes are E(B-V)$`{}_{\mathrm{FIR}}{}^{}E(BV)0.4`$. Panel (b) contains the clusters with $`|b|<20^{}`$ where most points follow the identity function up to E(B-V)$``$ 1.0. However, an important fraction of points in the range $`0<E(BV)<1`$ has large deviations in the sense of higher E(B-V)<sub>FIR</sub>. For $`E(BV)>1`$ the points deviate systematically from the identity function in the sense that E(B-V)<sub>FIR</sub> values are higher. A possible interpretation would be dust contributions for E(B-V)<sub>FIR</sub> arising from the cluster background.
### 4.1 $`\beta `$E(B-V): possibility of background reddening
In order to check the possibility of background reddening in the directions of globular and old open clusters we define the difference $`\beta `$E(B-V) = E(B-V)<sub>FIR</sub> \- E(B-V) (Tables 2 and 3, respectively). Figure 6 shows $`\beta `$E(B-V) histograms considering both samples together. For high latitude clusters ($`|b|>20^{}`$) we find a tight gaussian distribution suggesting an error distribution. We recall that E(B-V)<sub>FIR</sub> uncertainties typically amount to 16 % (Schlegel et al. 1998). The gaussian peak is in the bin 0-0.02, indicating a small offset between the two reddening types with higher values for E(B-V)<sub>FIR</sub>. This can also be seen in Panel (a) of Fig.5 as a small systematic shift of the points with respect to the identity function. There are two deviating objects in Panel (a) of Fig.6 which are the globular cluster M 107 (NGC 6171) and the open cluster NGC 1901. For M 107 E(B-V)<sub>FIR</sub> = 0.45 and E(B-V) = 0.33 (Table 3), thus $`\beta `$E(B-V) = 0.12. Recently Salaris & Weiss (1997) derived E(B-V) = 0.38 from isochrone fitting on CCD photometry, and they remarked that for this cluster values in the literature are in the range $`0.30<E(BV)<0.48`$. We suspect that the positions of M 107 in Figs. 5 and 6 reflect an uncertainty in the reddening derived from the stellar methods. On the other hand NGC 1901 with E(B-V)<sub>FIR</sub> = 0.33 and E(B-V) = 0.06 (Table 3) has the LMC disk as background (Sanduleak & Philip 1968), so that the high $`\beta `$E(B-V) = 0.27 must reflect dust emission from LMC complexes. Panels (b) and (c) deal with $`\beta `$E(B-V) histograms for low latitude clusters ($`|b|<20^{}`$), respectively for $`\beta E(BV)<1.5`$ and high $`\beta `$E(B-V) values. Similarly to panel (a) there occurs in panel (b) a peak near zero (bin 0-0.04), which indicates that most points indeed closely follow the identity function (Panel (b) of Fig.5). Thus, Schlegel et al.’s (1998) reddening values at low latitudes agree with those derived from stellar data for about two thirds of the clusters. Finally, the histogram for high $`\beta `$E(B-V) values (Panel (c)) shows 18 clusters with $`\beta E(BV)>1.0`$. These interesting objects together with the deviating clusters in Panel (b), typically with $`\beta E(BV)>0.30`$, are discussed in detail in Sect.4.4 for the possibility that their reddening values have an origin in the background dust.
Since mass loss is important in the last stages of red giant evolution dust accumulation in globular clusters is not unexpected. Cloudlets would contribute to differential reddening in CMDs as well as to 100 $`\mu `$m dust emission. Forte & Mendez (1988) found evidence for dust within globular clusters. They studied ten southern globular clusters, in particular NGC 362 and NGC 6624, and detected by means of CCD imaging regions with light deficiency which was attributed to dark clouds with intrinsic extinctions close to A<sub>V</sub> = 2.5. Their sizes are on the order of tenths of a parsec and they occur near the cluster nucleus. We checked whether Schlegel et al.’s reddening values are sensitive to internal dust contributions in the clusters NGC 362 and NGC 6624. We extracted E(B-V)<sub>FIR</sub> values for a cross with 17 pixels in Schlegel et al.’s maps (each pixel has 2.4 $`\times `$ 2.4). This cross samples the cluster main body and zones outside it, but still within the tidal radius (Trager et al. 1995). We noted fluctuations in E(B-V)<sub>FIR</sub> not exceeding 0.01 and 0.02, respectively. We conclude that Schlegel et al.’s reddening values are not particularly sensitive to the cloudlets owing to the large pixel size and the cloudlets’ small covering factor.
Star clusters within the dust layer as defined in Sect.2.2 ($`|Z^{}|200`$ pc) are expected to have significant differences between E(B-V)<sub>FIR</sub> and E(B-V) values. In order to study the behaviour of $`\beta `$E(B-V) we considered the cluster perpendicular distance to the Galactic plane ($`Z^{}`$) calculated with Eq.(1), using the data in Tables 2 and 3 for the globular and old open clusters. Figure 7 shows $`\beta `$E(B-V) as a function of $`Z^{}`$ for four $`Z^{}`$ ranges. Panel (a) shows objects up to 2 kpc from the Plane. Clearly, there is a large scatter of $`\beta `$E(B-V) values within the 200 pc dust layer, together with a significant wing which extends to $`400`$ pc. The scatter suggests that most of the differences between reddening values derived from dust emission and the stellar content are due to dust clouds in the disk background of the clusters. Panels (b), (c) and (d) show the behaviour of $`\beta `$E(B-V) in the regions increasingly away from the disk. In all three panels $`\beta `$E(B-V) values are small which is consistent with the fact that all these objects are halo globular clusters. The only strongly deviating object in Panel (b) is NGC6144 (Sect.4.4). In (b) and (c) the average $`\beta `$E(B-V) value is slightly positive corresponding to the small offset caused by higher E(B-V)<sub>FIR</sub> values. Unless an extremely thin diffuse distribution occurs throughout the halo caused by e.g. cooling flows and/or debris from dwarf galaxies accreted by the Milky Way, this offset observed for halo globular clusters implies that either Schlegel et al.’s zero point is slightly overestimated or that intrinsic reference colors, spectral distributions and isochrones were exceedingly red. Finally, in Panel (d) the offset is not present, but the sample is small.
Large differences $`\beta `$E(B-V) could be caused by the existence of dust clouds behind the clusters, primarily within the dust layer. In order to investigate this possibility we calculated cluster positions in the Galaxy and compared them to the assumed dust layer distribution. The latitude with respect to the true Galactic plane is given by
$$\mathrm{tan}b^{}=\frac{Z^{}}{d\mathrm{sin}b}\mathrm{tan}b.$$
(2)
Using this corrected latitude $`b^{}`$ it is possible to calculate the perpendicular distance of the cluster to the true Galactic plane $`d^{}`$, the distance along the cluster line of sight from the true Plane up to the dust layer edge d<sub>layer</sub>, and for the clusters within the dust layer the path behind them d<sub>bck</sub>:
$$d^{}=\frac{Z^{}}{\mathrm{sin}b^{}},d_{\mathrm{layer}}=\frac{200pc}{\mathrm{sin}b^{}},d_{\mathrm{bck}}=d_{\mathrm{layer}}d^{}$$
(3)
where 200 pc refers to the conservative dust layer height, assumption of Sect.2.2. We also assumed a galactic disk radius $`R=15`$ kpc.
Figure 8 shows in panel (a) and in the blowup (b) E(B-V)<sub>FIR</sub> as function of d<sub>layer</sub>. We note in panel (a) that objects with large E(B-V)<sub>FIR</sub> values (e.g. NGC 6134, Pişmiş 19, Liller 1, IC 4291) have a large path length within the layer. This is expected since the reddening derived from dust emission (E(B-V)<sub>FIR</sub>) should integrate dust contributions along the whole path throughout to the disk edge (d<sub>layer</sub>). Panel (b) suggests two correlations with different slopes, which might indicate differences in the cumulative effect of dust emission due to a discrete distribution of dust clouds. In panels (c) and its blowup (d) the behaviour of $`\beta `$E(B-V) is shown as a function of the path behind the cluster up to the disk edge (d<sub>bck</sub>). Panel (c) suggests that large $`\beta `$E(B-V) values for clusters like NGC 6134, Pişmiş 19, Liller 1, IC 4291 could arise from a dust cloud distribution behind the clusters since their directions and positions in the Galaxy imply a large background path up to the disk edge. Panel (d) shows considerable scatter which might be due to different origins: (i) inhomogeneous distribution of dust clouds, (ii) considerable uncertainties, (iii) the assumptions of a dusty disk are not satisfactory. Notice that for the assumed dust layer height of 200 pc the clusters outside the dust layer (d$`{}_{\mathrm{bck}}{}^{}=0`$) have an considerable range of values $`0<\beta E(BV)<1`$. This suggests that the Milky Way dust lane could be thicker (Sect.4.4).
### 4.2 Directions of some reddened young open clusters
Since the young disk is considerably thinner than the old disk (Janes & Phelps 1994, Friel 1995) it is worthwhile to study some interesting cases. We discuss some of the most reddened optical open clusters.
NGC 3603 and Westerlund 2 are clusters embedded in H II region complexes, where internal reddening is important. NGC 3603 ($`\mathrm{}`$ = 291.61, b = -0.52) is located at a distance d<sub>sun</sub> = 7 kpc and has E(B-V) = 1.44 from the CMD (Melnick et al. 1989), which comprises both the internal and foreground reddening. From the integrated spectrum Santos & Bica (1993) obtained a foreground reddening of E(B-V)<sub>f</sub> = 1.18, implying an internal reddening E(B-V)<sub>i</sub> = 0.26, by using a template spectrum which included internal absorption. Westerlund 2 ($`\mathrm{}`$ = 284.27, b = -0.33) at a distance d<sub>sun</sub> = 5.7 kpc has E(B-V) = 1.67 from the CMD (Moffat et al. 1991). Piatti et al. 1998b derived E(B-V)<sub>f</sub> = 1.40 and E(B-V)<sub>i</sub> = 0.27 by means of an integrated spectrum analysis.
Westerlund 1 ($`\mathrm{}`$ = 339.55, b = -0.40) is possibly the most reddened open cluster which can be optically observed. By means of CMDs and integrated spectrum Piatti et al. (1998b) derived E(B-V) = 4.3 and d<sub>sun</sub> = 1.0 kpc.
From Schlegel et al.’s (1998) reddening map we obtained very high reddening values derived for these young disk objects at very low galactic latitudes: E(B-V)<sub>FIR</sub> = 59.7, 65.7 and 12.3 respectively for NGC 3603, Westerlund 2 and Westerlund 1. The extremely high E(B-V)<sub>FIR</sub> for NGC 3603 and Westerlund 2 are probably related to lines of sight intercepting dust cloud cores (Sect.2), presumably the molecular clouds from which they were formed. Since Westerlund 1 is projected close to the Plane not far from the galactic center direction, its high E(B-V)<sub>FIR</sub> can be explained by the dust cumulative effect produced by a series of spiral arms and the Molecular Ring in that direction (Sect.2.1).
### 4.3 Reddening in the Sagittarius Dwarf direction
The globular clusters associated with the Sagittarius Dwarf are indicated in Table 1. Their E(B-V) values derived from the stellar content are comparable to those derived from the dust emission. Only a small systematic difference occurs, in the sense that values derived from dust emission are larger by $`\beta `$E(B-V) = 0.01-0.02. Assuming that M 54, Terzan 8, Arp 2 and Terzan 7 are slightly foreground or within Sagittarius itself, this sets a very low upper limit to the dust content in Sagittarius, in agreement with the fact that it is very depleted in H I (Koribalski et al. 1994).
### 4.4 Evidence for dark clouds with $`|Z^{}|>200`$ pc
Assuming the Milky Way dust layer height to be $`|Z^{}|=200`$ pc, Table 4 presents the 60 clusters with $`\beta E(BV)>0.30`$, showing their height from the Plane $`|Z^{}|`$, the distance in the cluster line of sight from the true Plane up to the dust layer edge d<sub>layer</sub>, and the path length behind the cluster up to the disk border d<sub>bck</sub>. They are separated in two groups, one formed by clusters within the dust layer (10 globular and 33 old open clusters) and the other by clusters outside it (15 globular and 2 old open clusters). As discussed in Sect.4.1 the large $`\beta `$E(B-V) values for clusters within the dust layer can be explained by dust clouds behind the clusters. The clusters outside the dust layer with large $`\beta `$E(B-V) values (see also their distribution in Fig.7 and the wing in the distribution above the dust layer) might be explained by higher $`|Z^{}|`$ dust clouds. A well-known example is the high latitude dust cloud Draco Nebula at a height from the Plane 300 to 400 pc (Gladders et al. 1998). In addition, star forming complexes as traced by means of Wolf-Rayet stars indicate that they are concentrated within a height from the Plane of 225 pc, but some attain 300 pc (Conti & Vacca 1990).
Assuming a dust layer height $`|Z^{}|=300`$ pc there would remain only 7 clusters outside the dust layer (Table 4). This corresponds to 3% of the total sample of 250 star clusters in the present study. They are all globular clusters: Lyngå 7, NGC 6144, NGC 6256, NGC 6355, NGC 6380, Tonantzintla 2 and NGC 6401. Among these clusters only NGC 6144 is very far from the plane at $`|Z^{}|2.8`$ kpc (Table 4). Although NGC 6144 has no CCD photometry yet, the photographic CMD and the integrated light reddening estimates (Harris 1996 and references therein) are consistent at about E(B-V) $``$ 0.32. We suspect that E(B-V)<sub>FIR</sub> = 0.71 for this cluster is overestimated, arising from foreground dust heated above the upper limit 21 K (Schlegel et al. 1998) by the hot star $`\sigma `$ Scorpii. Indeed, NGC 6144 is seen through the edge of the $`\rho `$ Ophiuchi dark cloud complex in the association Upper Scorpius at distance d<sub>sun</sub> = 125 pc (de Zeeuw et al. 1999). The cluster pathsight crosses the reflection nebula illuminated by the red supergiant Antares, also designated as IC 4606, Cederblad 132 or vdB-RN 107 in the catalogue of reflection nebulae by van den Bergh (1966). The neighbouring star $`\sigma `$ Scorpii is double (B2 III + O9.5 V), and ionizes the H II region Sh 2-9 (Sharpless 1959), or Gum 65 (Gum 1955). This hot double star also has its reflection nebula component Cederblad 130 (vdB-RN 104), which almost overlaps with the Antares reflection nebula. It is possible that dust grains in the direction of NGC 6144 are being heated by this particular configuration. At any rate CCD photometry of NGC 6144 is necessary to definitely establish the reddening E(B-V) affecting the cluster stars.
Howk & Savage (1999) detected high-Z dust structures in a sample of 7 edge-on spiral galaxies (NGC 891, NGC 3628, NGC 4013, NGC 4217, NGC 4302, NGC 4565 and NGC 4634). The thickness of the dust lanes in these galaxies is in the range $`500<2\times |Z|<900`$ pc. The high-Z dust features have typical dimensions of hundreds of parsecs and are located at heights in the range 500-1450 pc. The present study indicates the need of a Milky Way dust layer of thickness $`2\times |Z|600`$ pc in order to explain the lines of sight of 97 % of the known intermediate age and old clusters. The remaining 3 % would require some higher Z clouds.
## 5 Concluding remarks
Since the 100 $`\mu `$m dust emission reddening maps of Schlegel et al. (1998) provide whole-sky reddening estimates with relatively high angular resolution it is important to explore them in detail for a better understanding of the dust properties in different directions.
We provided an overview of the distribution along the galactic plane and in some interesting latitude directions. The accumulation of dust clouds in different arms and large structures such as the Molecular Ring can be distinguished. Individual dust complexes, including their cores, have E(B-V)<sub>FIR</sub> values compatible with those of embedded clusters derived from infrared photometry. An exception is the Nuclear Region where the temperature in the Central Molecular Zone appears to be underestimated in the Schlegel et al.’s temperature maps.
The 100 $`\mu `$m dust emission reddening maps provide E(B-V)<sub>FIR</sub> values compatible to E(B-V) derived from the stellar content of globular and old open cluster for $``$ 75 % of the 250 directions probed in the present study ($`\beta E(BV)<0.30`$). The values for most high-latitude clusters ($`|b|>20^{}`$) are in good agreement; these are objects in general outside the disk dust layer, so that all dust in the line of sight is sensitive to both methods. An interesting exception is NGC 1901 which is outside the dust layer and has a E(B-V)<sub>FIR</sub> much larger than E(B-V). The background dust source in this case is the LMC disk.
The differences between the dust emission and stellar content reddening values occur most frequently for low latitude clusters ($`|b|<20^{}`$). Brightness selection effects due to reddening and distance particularly affect the open cluster sample. In the existing catalogues many intermediate age open clusters are yet to be studied, while future infrared surveys should reveal many new open clusters. These distant objects are expected to have increasing reddening values, but the $`\beta `$E(B-V) values should decrease since the pathsight behind the cluster within the dust layer decreases. Most of the known globular clusters have CMDs, but infrared surveys should reveal some new ones in disk and bulge zones. Those far in the disk should behave like the heavily reddened far open clusters.
Bandwidth effects on the intrinsic reddening law have implications on the results for larger extinctions. The fact that most of the points follow a 1:1 FIR/Optical relation up to E(B-V) = 1.0 (Fig.5) suggests that this effect should become important for reddening values beyond this limit. An additional complication for objects so heavily reddened that the stellar content can be studied only in the infrared is that the transformations to optical reddening values depend on grain properties.
From the spatial distribution of available objects and their relative positions with respect to the dust layer we conclude that background dust clouds are probably responsible for these differences. A dust layer with thickness $`2\times |Z|`$ 600 pc is required to explain the distribution of $``$ 97 % of the sample. Some additional higher Z dust clouds, like the Draco Nebula, would also be required to explain the rest.
The present study of reddening in star cluster directions suggests that the Milky Way is similar in dust layer thickness and occurrence of some high-Z dust structures to edge-on spirals studied by Howk & Savage (1999). In particular the Milky Way dust layer may be thicker than previously thought.
###### Acknowledgements.
We acknowledge support from the Brazilian institution CNPq. We thank an anonymous referee for interesting remarks.
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# Structure of the Mg ii and damped Lyman-𝛼 systems along the line of sight to APM 08279+5255 Based on observations collected at the W.M. Keck Observatory, which is operated as a scientific partnership among the California Institute of Technology, the University of California and the National Aeronautics and Space Administration. The Observatory was made possible by the generous financial support of the W.M. Keck Foundation.
## 1 Introduction
The gravitationally lensed Broad Absorption Line (BAL) QSO APM 08279+5255 ($`z_{\mathrm{em}}`$ = 3.911) has been given much attention since its discovery by Irwin et al. (1998), as it is one of the most luminous objects in the universe even after correction for the gravitational lensing induced amplification. Adaptive-optics imaging has revealed two main components (Ledoux et al. 1998b), separated by 0.378$`\pm `$0.001 arcsec as measured on HST/NICMOS data (Ibata et al. 1999), and of relative brightness $`F_\mathrm{B}/F_\mathrm{A}`$ = 0.773$`\pm `$0.007. The HST images reveal also the presence of a third object C with $`F_\mathrm{C}/F_\mathrm{A}`$ = 0.175$`\pm `$0.008, located in between A and B and almost aligned with them. The point-spread-function model fits on the three objects are consistent with the three components being point-sources, and their colors are similar within the uncertainties. There is no trace of the lensing object up to magnitude $`V`$ = 23.
A high S/N ratio high-resolution spectrum of APM 08279+5255 was obtained at the Keck telescope (Ellison et al. 1999a,b), and made available to the astronomical community. This spectrum, though complicated by the combination of light traveling along three different sightlines, is a unique laboratory for studying the intervening and associated absorption systems.
In this paper we study the structure of six intervening Mg ii systems at 1.2 $`<`$ $`z`$ $`<`$ 2.07 and the physical characteristics of the gas in two very strong Mg ii systems detected at $`z_{\mathrm{abs}}`$ = 1.06 and 1.18, which, we argue, are damped Lyman-$`\alpha `$ systems and may well reveal the lensing galaxies. We also comment on a third probable damped Lyman-$`\alpha `$ system at $`z_{\mathrm{abs}}`$ = 2.974. This paper is organized as follows: the data are described in Section 2; the structure of the intervening Mg ii systems is investigated using the covering factor analysis in Section 3; we demonstrate that the Mg ii systems at $`z_{\mathrm{abs}}`$ = 1.06 and 1.18 are damped Lyman-$`\alpha `$ systems in Section 4; we discuss a probable damped Lyman-$`\alpha `$ system at $`z_{\mathrm{abs}}`$ = 2.974 in Section 5. We adopt $`H_\mathrm{o}`$ = 75 km s<sup>-1</sup> Mpc<sup>-1</sup> and $`q_\mathrm{o}`$ = 0.5 througout the paper.
## 2 Data
A high S/N ratio, high-resolution spectrum of the $`z_{\mathrm{em}}`$ = 3.911 quasar APM 08279+5255 was obtained with the HIRES echelle spectrograph at the 10m Keck-I telescope (Ellison et al. 1999a,b). This data was made public together with a low-resolution spectrum of the quasar and a high-resolution spectrum of a standard star. We have corrected the high-resolution spectrum of APM 08279+5255 for small discontinuities in the continuum, which are probably due to the inappropriate merging of different orders. These discontinuities have been recognized by comparing the high and low-resolution spectra. The latter is also used for normalization of the high-resolution data. Atmospheric absorption features were identified from the standard star spectrum. Voigt profile fitting of the absorption features have been performed using the context FITLYMAN (Fontana & Ballester 1995) of the European Southern Observatory data reduction package MIDAS and the code VPFIT (Carswell et al. 1987). We have measured the final spectral resolution by fitting the narrow atmospheric absorption lines which are free of blending. We find $`FWHM`$ $``$ 8 km s<sup>-1</sup> ($`b`$ $``$ 4.8 km s<sup>-1</sup>) at 6900 Å, $`R`$ = 37500, and use this value throughout the paper.
## 3 Structure of the intervening $`z`$ $`>`$ 1 Mg ii absorbers
Ellison et al. (1999b) have already noted that the two lines of some of the Mg ii$`\lambda \lambda `$2796,2803 doublets cannot be fitted with the same column density and Doppler parameter. As can be seen on Fig. 1, most of the systems have a doublet ratio close to unity inspite of having residual intensities in the normalized spectrum close to 0.5.
If the background source is a point source and if the Mg ii$`\lambda `$2796 absorption line is resolved, then the residual intensity of the normalized spectrum measured at any velocity $`v`$ with respect to the centroid of the line, is equal to $`e^{\tau (v)}`$, where $`\tau (v)`$ is the optical depth at $`v`$ and the residual intensity of the Mg ii$`\lambda `$2803 line is $`e^{\tau (v)/2}`$. It is apparent from Fig. 1 that the above condition is not fullfilled for most of the Mg ii doublets.
There are two possibilities to explain this. If the lines are not resolved, the measured residual intensity is affected by convolution of the true absorption profile with the instrumental profile. This can introduce artificial residuals at the bottom of the saturated absorption features (see e.g. Lespine & Petitjean 1997). Alternatively, as the observed light is a combination of light from different images, it is possible that the column densities are different along different lines of sight and, as a limiting case, that the absorbing cloud does not cover all the images. We discuss these different possibilities below.
### 3.1 Unresolved narrow-components
If the absorption profiles are made up of unresolved components, the observed residual intensity does not correspond to the real optical depth. Given the resolution of the spectrum ($`R`$ $``$ 37500), a saturated Mg ii$`\lambda `$2796 line can have a residual intensity in the normalized spectrum of 0.5 if its Doppler parameter is smaller than 1.5 km s<sup>-1</sup>. In that case, the residual intensity of the Mg ii$`\lambda `$2803 line is in the range 0.5–0.6, depending on the actual column density. The two equivalent widths differ by no more than 20% (see e.g. Lespine & Petitjean 1997).
The system at $`z_{\mathrm{abs}}`$ = 1.5497 can be indeed fitted this way using 10 components with $`b`$ values in the range 1.1–1.7 km s<sup>-1</sup>. In that case the well detached cloud in the red wing (see Fig. 1) is fitted with two adhoc nearly identical components though the Mg ii$`\lambda `$2796 profile is perfectly fitted with a single resolved component model. However the one-component model cannot explain the strengths of the two Mg ii lines without invoking partial coverage (see below).
We use the $`z_{\mathrm{abs}}`$ = 1.5523 system to illustrate the case. The final spectrum has a defect in the center of the Mg ii$`\lambda `$2796 line. We have therefore used three individual exposures of high S/N ratio (Ellison private communication) to correct for this. The final optical depth variations from one spectrum to the other is about 2%. Fig. 2 shows the best fit to the doublet considering full coverage ($`f`$ = 1). Five narrow components ($`b`$ $`<`$ 1.5 km s<sup>-1</sup>) are needed. Although good (reduced $`\chi ^2`$ of 1.5), the fit is not completely satisfactory. We have fitted consistently the Mg ii together with the Fe ii lines considering that only one of the brightest sources is covered. With five components, a good fit (reduced $`\chi ^2`$ of 0.9), shown on Fig. 3, is obtained if B is not covered ($`f`$ = 0.6). Details of the subcomponent parameters are given in Table 1. As the Mg ii doublet ratio is close to one, small $`b`$ values are needed even with the assumption of partial coverage.
It should be noted however that a Doppler parameter smaller than 1.5 km s<sup>-1</sup> corresponds to a temperature smaller than 3500 K, a surprizingly small temperature for this gas which is most probably ionized. Allowing for a smaller number of components, we can find a good fit with only three components, $`f`$ = 0.45 and reduced $`\chi ^2`$ = 0.6 (see Fig. 3 and Table 1). The $`b`$ values are larger than 2.5 km s<sup>-1</sup> relaxing the restriction on the temperature.
A statistically acceptable fit is difficult to find for the systems at $`z_{\mathrm{abs}}`$ = 1.221 and 1.5523. The red wing of the Mg ii$`\lambda `$2803 line at $`z_{\mathrm{abs}}`$ = 1.221 is blended with another absorption feature which we found difficult to identify. It could be Mg i$`\lambda `$2852 at $`z_{\mathrm{abs}}`$ = 1.1727. This system is possibly detected by Fe ii$`\lambda \lambda `$2344,2382, Mn ii$`\lambda `$2576 and Ca ii$`\lambda `$3934 (see Fig. 4). The other lines are either below the detection limit or blended. Confirmation of the Ca ii lines, which are free of blending, would be particularly important as, probably, the presence of this additional system, together with the presence of the three damped systems at $`z_{\mathrm{abs}}`$ = 1.062, 1.181 and 2.974 (see below) should be taken into account in any model of the lens.
### 3.2 Partial covering factor
#### 3.2.1 Computing a covering factor
We can interpret the observations in terms of a covering factor which is the fraction of the background source covered by the absorbing cloud. The relative brightness of the three sources are $`F_{\mathrm{A},\mathrm{B},\mathrm{C}}`$/$`F_{\mathrm{tot}}`$ = 0.513, 0.397 and 0.090 (Ibata et al. 1999). If one line-of-sight is completely absorbed (condition imposed by the fact that the doublets are saturated) and the other free of absorption, then the covering factor is 0.40, 0.49, 0.51 and 0.60 if, respectively, B only, B+C, A only and A+C are covered (the case C only is very unlikely as C is close to A and located in between A and B, Ibata et al. 1999; see however Srianand & Petitjean 2000). Given the uncertainties, if only one line-of-sight is completely absorbed, the covering factor should be in the range 0.4–0.6. Of course, it is larger if the second line-of-sight is not completely clear. We have computed the covering factor for the Mg ii systems using the method described by Srianand & Shankaranarayanan (1999). This assumes that the lines are resolved. It can be seen in Fig. 1 that for the three systems with $`z_{\mathrm{abs}}`$ $`<`$ 1.7, the covering factor ranges between 0.5 and 0.6 whereas for the systems with $`z`$ $`>`$ 1.7, the covering factor is larger (but always less than 0.8) with the possible exception of the $`z_{\mathrm{abs}}`$ = 2.0668 system. The latter system is quite weak however and uncertainties are large. The values of the covering factor for the three lower redshift systems suggest that the clouds cover one of the two brightest background sources only. This has to be investigated in more detail however.
#### 3.2.2 Optical depths along different sightlines
In this Section we investigate the effect of the optical depth being different along different sight lines. In order to make our analysis simpler, we consider only two images, A and B, with fractional flux contributions $`F_1`$ = 0.6 and $`F_2`$ = 0.4. Suppose the optical depth along the two sight lines are $`\tau _1`$ and $`\tau _2`$ then the measured residual intensities are,
$`R(2796)`$ $`=`$ $`F_1e^{\tau _1}+F_2e^{\tau _2}`$ (1)
$`R(2803)`$ $`=`$ $`F_1e^{\tau _1/2}+F_2e^{\tau _2/2}`$
The residual intensities can be written as (Srianand & Shankaranarayanan 1999),
$`R(2796)`$ $`=`$ $`1f+fe^\tau `$ (2)
$`R(2803)`$ $`=`$ $`1f+fe^{\tau /2}`$
where, $`f`$ and $`\tau `$ are the resulting covering factor and optical depth and therefore,
$$f=\frac{[1R(2803)]^2}{1+R(2796)2R(2803)}$$
(3)
From Eqs. (1) and (3) one can derive $`f`$ as a function of $`F_1`$, $`F_2`$, $`\tau _1`$ and $`\tau _2`$. This analysis assumes that the absorption profiles are resolved in the HIRES spectrum. As $`F_1`$ and $`F_2`$ are known from observation, the covering factor depends on $`\tau _1`$ and $`\tau _2`$ only.
In order to investigate the parameter space we have computed and plotted on Fig 5, the covering factor (panel d), residual intensities of the lines (panels a and b) and their ratio (panel c) as a function of $`\tau _2`$ for different values of $`\tau _1`$. In panels (a) and (b), the two horizontal dotted lines show the range of observed values for the Mg ii systems with $`z_{\mathrm{abs}}`$$`<`$1.7. In panel (c) the dotted lines give the measured values for the three low-redshift Mg ii systems.
As expected, when $`\tau _1=\tau _2`$, $`f=1`$ (albeit with various ratios of residual intensities) and for $`\tau _1\tau _2`$ the covering factors are less than 1.0. When $`\tau _1\tau _2`$ (respectively $`\tau _1\tau _2`$), the covering factor is in the range 0.6–1.0 (respectively 0.4–1.0). Conversely, the observed residual intensities, $`R`$(2796) and $`R`$(2803), together with the measured covering factor, $`f`$, can be used to constrain $`\tau _1`$ and $`\tau _2`$.
The covering factor estimates for $`z_{\mathrm{abs}}`$$`<`$ 1.7 systems are in the range 0.5 and 0.6 (see Fig. 1). This, together with the observed residual intensities, indicates that the absorbing gas is saturated along one sight line only with optical depth ratios as large as ten. For example, the well detached component in the red wing of the $`z_{\mathrm{abs}}`$ = 1.5497 system has $`f=0.6`$ (with a typical error of 0.02), a residual intensity ratio $``$0.84 (with a typical error of 0.02) and R(2796)$``$0.40 at the core of the line. This implies that the contribution to this absorption comes mainly from the line of sight toward A+C with the optical depth along B being more than an order of magnitude smaller.
### 3.3 Equivalent width ratio
The precise determination of the strength of the absorption lines along different lines of sight should await HST/STIS spectroscopic observations. Though we derived some information on this in the previous Section using the absorption line residual intensities, the results depend crucially on the assumption that the absorption lines are resolved. We can complement the previous analysis using the total equivalent widths and their ratios without making any assumption about the spectral resolution In addition, when considering total equivalent widths, the consequence of contamination by weak lines is small unlike in the case of the analysis of the residual intensities. The constrains are much weaker however.
Let us assume that the absorption is saturated along line of sight number one and optically thin along line of sight number two. Therefore $`W_1^{\mathrm{real}}`$(2796) = $`W_1^{\mathrm{real}}`$(2803) and $`W_2^{\mathrm{real}}`$(2796) = 2$`\times `$$`W_2^{\mathrm{real}}`$(2803). The combined observed equivalent width ratio is:
$$\frac{W^{\mathrm{obs}}(2803)}{W^{\mathrm{obs}}(2796)}=\frac{\frac{1}{2}+\frac{W_1^{\mathrm{real}}(2796)}{W_2^{\mathrm{real}}(2796)}\frac{F_1}{F_2}}{1+\frac{W_1^{\mathrm{real}}(2796)}{W_2^{\mathrm{real}}(2796)}\frac{F_1}{F_2}}$$
(4)
where $`F_1`$ and $`F_2`$ are the fractional flux contributions of the two distinct background sources. In Fig. 6, we have plotted $`W_1^{\mathrm{real}}`$(2796) /$`W_2^{\mathrm{real}}`$(2796) versus $`W^{\mathrm{obs}}`$(2803) /$`W^{\mathrm{obs}}`$(2796) for $`F_1`$/$`F_2`$ = 0.65 (1 is B; 2 is A+C), 1 (1 is A or B+C; 2 is B+C or A) and 1.5 (1 is A+C; 2 is B). The vertical dashed-dotted lines correspond to the observed doublet ratios of the systems at $`z_{\mathrm{abs}}`$ = 1.221, 1.5523 and of the reddest component of the $`z_{\mathrm{abs}}`$ = 1.5497 system.
From Fig. 6, it can be seen that the $`z_{\mathrm{abs}}`$ = 1.221 system (doublet ratio of 0.96) requires the ratio of the equivalent widths of Mg ii$`\lambda `$2796 along the two lines of sight to be larger than 8. and the column densities along the two sightlines differ by more than an order of magnitude.
### 3.4 Dimension of individual clouds
Rauch et al. (1999) have observed strong variations of C ii and Si ii absorptions at $`z`$ = 3.538 along two sightlines separated by only 13$`h^1`$ pc. However, as the velocity difference between the quasar and the system is only 6000 km s<sup>-1</sup>, it cannot be excluded that the latter system is somehow associated with the quasar. Variations of the strength of metal line systems have also been reported along adjacent lines of sight with larger separations (5–10 kpc) by Monier et al. (1998) and Lopez et al. (1999). Each time however a damped Lyman-$`\alpha `$ system is seen along one of the sightlines. Contrary to these previous studies, the Mg ii systems we examine here are most likely to be associated with halos of intervening galaxies.
From the detection of associated galaxies, radii of the order of 35$`h^1`$ kpc have been derived for Mg ii halos producing absorptions with equivalent widths $`W_\mathrm{r}`$ $`>`$ 0.3 Å at $`z`$ $`<`$ 1 (Bergeron & Boissé 1991, Steidel 1993). Dimensions of the same order have been derived from the study of Mg ii systems seen along two lines of sight separated by 3 arcsec (Smette et al. 1995). The latter authors find a lower limit of 22 $`h_{50}^1`$ kpc for the radius of Mg ii absorbers with $`W_\mathrm{r}`$ $`>`$ 0.3 Å at 0.5 $`<`$ $`z`$ $`<`$ 1.3.
If we assume that the lensing galaxy of APM 08279+5255 is at $`z_{\mathrm{lens}}1`$ (see next Section), the separation between the two lines of sight to A and B decreases from 1.7 to 0.7 $`h_{75}^1`$ kpc ($`q_\mathrm{o}`$ = 0.5) between $`z`$ = 1.22 and $`z`$ = 2.04 and is more than an order of magnitude smaller than the radius of Mg ii halos at intermediate redshift. Although evolution is probable, it would be really surprizing that the Mg ii systems studied here with $`W_\mathrm{r}`$ $`<`$ 0.5 Å at $`z_{\mathrm{abs}}`$ $`>`$ 1.2 have characteristic dimensions more than an order of magnitude smaller than what is derived at lower redshift. If true, this would suggest that the structure of the Mg ii halos at these redshifts differs substantially from that at lower redshifts.
A more likely explanation of these observations is that the halos are composed of a collection of clouds (see Petitjean & Bergeron 1990; Srianand & Khare 1994) and that individual clouds cover only one sight line. The number density of clouds is large enough so that the total covering factor of the halo is close to one, consistent with observations of associated galaxies. However, individual clouds, regularly spread over the velocity profile by kinematics, cover only one image of the lens. The number density of clouds is not large enough for the absorption material to cover the two lines of sight at all velocities. The distance over which the optical depth, and hence the column density of Mg ii , changes by at least one order of magnitude at $`z_{\mathrm{abs}}`$$`<`$1.7 is smaller than $`1h_{75}^1`$ kpc. In contrast, the two strong Mg ii systems at $`z_{\mathrm{abs}}`$ = 1.06 and 1.18 (see below) have covering factor equal to one (the lines are saturated and go to the zero level) over more than 200 km s<sup>-1</sup>. These latter systems are likely to arise due to absorption through central regions of galaxies where the number of clouds is so large that saturated absorption occurs along both lines of sight whatever the radius of the individual clouds might be.
## 4 The two Mg ii systems at $`z_{\mathrm{abs}}`$ =1.062 and 1.181
The presence of a strong Mg ii system ($`W_\mathrm{r}\lambda `$2803 $``$ 2.4 Å) at $`z_{\mathrm{abs}}`$ $``$ 1.18 was already mentioned by Irwin et al. (1998). There is an additional even stronger Mg ii system at $`z_{\mathrm{abs}}=1.062`$ ($`W_\mathrm{r}\lambda `$2803 $``$ 3.3 Å). Although the Mg ii lines are redshifted in the Lyman-$`\alpha `$ forest and may be blended with Lyman-$`\alpha `$ intervening absorptions, the existence of the system is confirmed by numerous lines redshifted redward of the quasar Lyman-$`\alpha `$ emission. As the two brightest images of the lensed quasar have similar magnitudes, it is expected that the lines of sight to both images pass through the core of the lensing object where strong Mg ii absorption is likely to occur. The gravitational lensing may thus result from the cumulative effect of the two galaxies associated with these two absorbing systems together with the objects responsible for the possible system at $`z_{\mathrm{abs}}`$ = 1.1727 (see Section 3.1) and the other damped Ly$`\alpha `$ system at $`z_{\mathrm{abs}}`$ = 2.974. Absorptions from Mg ii, Fe ii, Ca ii, Mn ii, Ti ii and Na i are seen in both systems (see Figs. 7,8).
### 4.1 Na i absorptions
In both Mg ii systems, the weak Na i$`\lambda \lambda `$3303,3303 doublet is detected (see Figs. 7 and 9). The fact that the two lines of the doublet are seen with consistent strengths gives confidence that the identification is correct. We have identified lines from other metal line systems in the vicinity of the doublet. The spectra of APM 08279+5255 and the standard star are compared in Fig. 9 to rule out the possibility that the absorption features are of atmospheric origin.
The column densities obtained by Voigt profile fitting are large, log $`N`$(Na i) = 12.9 and 13.5 at, respectively, $`z_{\mathrm{abs}}`$ = 1.0626 and 1.1801. The separation of the two principal lines of sight is $``$1.9$`h_{75}^1`$ kpc at $`z`$ $``$ 1. It is therefore possible that the clouds seen by their Na i absorptions do not cover the two brightest lines of sight. In that case, however, the column density could be even larger by a factor of two (see Section 4.4).
Using the measurements by Sembach et al. (1993) and Diplas & Savage (1994), Bowen et al. (1995) find that , in our Galaxy, log $`N`$(H i) = 0.688 log $`N`$(Na i) + 12.16. Note that this correlation holds up to column densities log $`N`$(H i$`>`$ 21 (see e.g. Ferlet et al. 1985). Applying this correlation for the $`z_{\mathrm{abs}}`$ = 1.0626 and 1.1801 absorbers gives neutral hydrogen column densities of the order of log $`N`$(H i$``$ 21.0 and 21.4 respectively, for gas-phase metallicities comparable to what is seen in the interstellar medium in our Galaxy (the neutral hydrogen column densities could be even larger if the metallicity in these intermediate-redshift systems is smaller than in our Galaxy). Such large values for $`N`$(H i) are supported by the column densities of other species that are found to be surprisingly close to what is observed in typical interstellar clouds (see below and Table 2). This leaves little doubt that the systems are indeed damped Lyman-$`\alpha `$ systems.
In our Galaxy, such high Na i column densities are seen only in dense and cool gas (see below). The typical $`b`$ values of the Na i diffuse components in both the local and low-halo gas is about 0.7 km s<sup>-1</sup> corresponding to $`T`$ $`<`$ 500 K (Welty et al. 1994). Although the conclusion is very uncertain given the resolution of the spectrum ($`R`$ $``$ 37500) and the double nature of the background source, the lines we observe are consistent with $`b`$ values as small as 1 km s<sup>-1</sup> (see below). From this, we derive an upper limit on the temperature of $`T`$ $`<`$ 2000 K.
### 4.2 Comments on each system
#### 4.2.1 $`z_{\mathrm{abs}}`$ = 1.062
A subset of the absorptions detected in this system is shown in Fig. 7. It can be seen that strong (but unsaturated) absorptions from Ti ii and Ca ii are detected. The profile of the Ti ii absorption is spread over about 250 km s<sup>-1</sup> but does not show any edge leading pattern (Prochaska & Wolfe 1998). We have selected to examine two components at $`z_{\mathrm{abs}}`$ = 1.0613 and 1.0631 because they show well detached absorptions in Ti ii, Mn ii, Ca ii and/or Fe ii plus the component at $`z_{\mathrm{abs}}`$ = 1.0626 in which we see Na i. Column densities are listed in Table 2. We have adjusted the best values for $`b`$ from the fit to the lines that are free of any blending, considering for simplicity that all species have the same Doppler parameter and assuming complete coverage. We find $`b`$ = 1.5, 1.5 and 1.9 for, respectively, $`z_{\mathrm{abs}}`$ = 1.0613, 1.0626 and 1.0631.
Ca i is not detected and the 3$`\sigma `$ upper limit on the column density in the three components we have selected is $`<`$ 10.43, 10.43 and 10.20. The Mg ii$`\lambda `$2796,2803 and Fe ii$`\lambda \lambda \lambda `$2382,2600,2586 lines are badly saturated. Moreover, they are redshifted in the Lyman$`\alpha `$ forest, which prevents any fit of the lines. However, the unsaturated Fe ii$`\lambda `$2260 features detected at $`z_{\mathrm{abs}}`$ = 1.0613 and 1.0626 (see Fig. 7) give a reliable estimate of the Fe ii column density, log $`N`$(Fe ii) = 13.90$`\pm `$0.80 and 14.10$`\pm `$0.80 for $`b`$ = 1.5 km s<sup>-1</sup>, the continuum being adjusted locally. The non-detection of Fe ii$`\lambda `$2367 at $`z_{\mathrm{abs}}`$ = 1.0631 gives an upper limit log $`N`$(Fe ii$`<`$ 14.60.
#### 4.2.2 $`z_{\mathrm{abs}}`$ = 1.181
A subset of the absorptions detected in this system is shown in Fig. 8. The profile of the Mg i absorption is spread over more than 200 km s<sup>-1</sup> but, as for the previous system, it does not show any edge leading pattern. A number of absorption lines are optically thin or moderately saturated and reliable column densities can be derived even though difficulties arise from most of the components being badly defined (see Fig. 8). We have selected for study two subcomponents which are clearly seen in all absorption profiles at $`z_{\mathrm{abs}}`$ = 1.1799 and 1.1801. They are indicated on Fig. 8 by vertical dashed lines, and the column densities obtained from Voigt-profile fitting are given in Table 2. Doppler parameters have been considered to be identical for all species. For $`z_{\mathrm{abs}}`$ = 1.1799 we find $`b`$ = 2.5 km s<sup>-1</sup>.
For the component at $`z_{\mathrm{abs}}`$ = 1.1801 in which Na i absorption is detected, the Doppler parameter is estimated by fitting the well-defined lines of the sodium doublet after having taken into account the effect of Na i$`\lambda `$3303.3 being partially blended with an atmospheric feature (see Fig. 9). We obtain a best value $`b`$ = 1.1$`{}_{0.5}{}^{}{}_{}{}^{+1.0}`$ km s<sup>-1</sup>. The column densities derived using the two values $`b\pm `$ 1$`\sigma `$ differ by large factors. We have therefore refined the determination of $`b`$ and $`N`$ using the following indirect argument.
The ratio of the Mg i to the Na i column densities in neutral gas can be written
$$\frac{N(\mathrm{Mg}\mathrm{i})}{N(\mathrm{Na}\mathrm{i})}=\frac{\mathrm{Mg}\mathrm{i}}{\mathrm{Mg}}\times \frac{\mathrm{Na}}{\mathrm{Na}\mathrm{i}}\times \frac{\delta _{\mathrm{Mg}}}{\delta _{\mathrm{Na}}}\times \frac{Z_{\mathrm{Mg}}}{Z_{\mathrm{Na}}},$$
(5)
where $`\delta `$ is the depletion of the element due to the presence of dust and $`Z`$ the abundance. Assuming that (i) the relative abundance of Na to Mg is solar, $`Z_{\mathrm{Mg}}`$/$`Z_{\mathrm{Na}}`$ = 19, (ii) the relative depletion into dust-grains is $`\delta _{\mathrm{Mg}}`$/$`\delta _{\mathrm{Na}}`$ $`>`$ 0.3 (Savage & Sembach 1996) and (iii) (Mg i/Mg ii)$`\times `$(Na ii/Na i$`>`$ 0.15 in cold and neutral gas (Péquignot & Aldrovandi 1986), we derive $`N`$(Mg i)/$`N`$(Na i$``$ 0.8. This estimation is certainly very approximate. However, this simple argument shows that the latter ratio cannot be much smaller than 0.5. If we assume $`b`$ = 1.5 km s<sup>-1</sup>, then we find $`N`$(Mg i)/$`N`$(Na i) = 0.045 which is definitively too small. We therefore have fitted the absorption lines decreasing $`b`$ from 1.5 to 0.5 km s<sup>-1</sup> to find the largest $`N`$(Mg i)/$`N`$(Na i) ratio. We find a maximum $`N`$(Mg i)/$`N`$(Na i) = 0.3 for $`b`$ = 0.8 km s<sup>-1</sup>, log $`N`$(Mg i) = 13.0 and log $`N`$(Na i) = 13.5.
### 4.3 Physical state of the gas
Table 2 contains the column densities measured in the five subcomponents defined above. The penultimate column gives for comparison the column densities measured by Welty et al. (1999) in the neutral gas toward 23 Ori. This gas is found to have temperature $`T`$ $``$ 100 K, hydrogen density $`n_\mathrm{H}`$ $``$ 10–15 cm<sup>-3</sup> (and therefore total thickness of 12–16 pc) and electronic density $`n_\mathrm{e}`$ $``$ 0.15$`\pm `$0.05 cm<sup>-3</sup>. The last column of Table 2 gives the measurements for the strongest component of the warm neutral gas toward $`\mu `$Col (Howk & Savage 1999). This gas is found to have $`T`$ $``$ 6000–7000 K and $`n_\mathrm{e}`$ $``$ 0.3 cm<sup>-3</sup>. The two sets of column densities are similar except for $`N`$(Na i) which is much smaller toward $`\mu `$Col. This illustrates directly that at least in the component at $`z`$ $``$ 1.1802 toward APM 08579+5255 where Na i is detected, the gas is most likely to be neutral and cold. Moreover, note also that in our Galaxy, a ratio $`N`$(Na i)/$`N`$(Ca ii$`>`$ 1, as observed at $`z_{\mathrm{abs}}`$ = 1.1802 toward APM 08579+5255, is characteristic of cold gas in the disk. Indeed, along the line of sight to the LMC, such large ratios are observed only at the systemic velocities of the LMC and the Galaxy; gas in between has $`N`$(Na i)/$`N`$(Ca ii$`<`$ 1 (Vidal-Madjar et al. 1987, Vladilo et al. 1993).
From the upper limit on Ca i and Fe i we can derive, in the components where $`N`$(Ca ii) and $`N`$(Fe ii) are measured, a lower limit for the electronic density for a given ionizing field. Indeed, $`n_\mathrm{e}`$ = ($`X^\mathrm{o}`$/$`X^+`$)$`\times `$($`\mathrm{\Gamma }`$/$`\alpha `$), where $`\mathrm{\Gamma }`$ is the photoionization rate and $`\alpha `$ the recombination coefficient. The corresponding electronic densities for a Galactic ionizing field (Péquignot & Aldrovandi 1986), are $`n_\mathrm{e}`$ $`<`$ 0.13 and 3 cm<sup>-3</sup> for Fe and Ca at $`z_{\mathrm{abs}}`$ = 1.0626 and 1.1801 respectively. Note that the determination of $`n_\mathrm{e}`$ in the interstellar medium from the ratio of singly ionized to neutral species is highly uncertain probably because of contamination of the singly ionized column density determination by adjacent components (Welty et al. 1999). Writing the same relation for sodium and equating the expression of $`n_\mathrm{e}`$ obtained for sodium and iron or calcium leads to upper limits on the Na i/Na ii ratio which depends only on the shape of the ionizing spectrum and not on its absolute value. With the only assumption that the ionizing spectrum has the same shape as in our Galaxy and using the coefficients derived from Péquignot & Aldrovandi (1986; see Welty et al. 1999), we find log Na i/Na ii $`<`$ $``$1.3 and 0.1 from the constraints obtained on Fe and Ca in the 1.0626 and 1.1801 systems respectively.
From $`N`$(H i) = $`N`$(Na i)$`\times `$(Na/Na i)/$`Z`$(Na) and using the two upper limits on the Na i/Na ii ratios derived above, we can write $`N`$(H i$`>`$ 19.9 and 19.5/($`Z`$(Na)/$`Z_{}`$(Na))/$`\delta `$(Na) at $`z_{\mathrm{abs}}`$ = 1.0626 and 1.1801 respectively, where $`\delta `$(Na) is the fraction of sodium remaining in the gas phase after depletion into dust-grains. This factor is equal to about 0.1 in the ISM (Savage & Sembach 1996). This adds support to arguments presented previously that these systems are damped. To illustrate the discussion, simple photo-ionization models using the code Cloudy (Ferland 1996) have been constructed. The absorbing cloud is modelled as a plane parallel slab with uniform density, solar chemical composition and neutral hydrogen column density $`5\times 10^{20}`$ cm<sup>-2</sup>. The elements are considered to be depleted into dust-grains as in the cool cloud observed toward $`\zeta `$Oph (Savage & Sembach 1996). The shape of the UV flux is taken to be a power-law $`F_\nu `$ $``$ $`\nu ^{1.0}`$. The resulting column densities of various species along a line-of-sight perpendicular to the slab are given versus the ionizing parameter in Fig. 10. As discussed above, the $`N`$(Mg i)/$`N`$(Na i) ratio can be smaller than one only if magnesium is more depleted into dust grains than sodium. Note that every model that produces enough Na i has temperature less than 100 K.
Finally, we do not detect any CH$`\lambda `$4300 and CH$`{}_{}{}^{+}\lambda `$4232 absorption. The limits on the column densities are log $`N`$(CH) $`<`$ 13.5 and log $`N`$(CH<sup>+</sup>) $`<`$ 13.0 at both $`z_{\mathrm{abs}}`$ = 1.06 and 1.18 (see Table 2). This is just what would be expected in our Galaxy along an otherwise similar line of sight. Indeed, along the line of sight to 23 Ori, log $`N`$(CH<sup>+</sup>) = 13.06 and log $`N`$(CH) = 12.69 (Welty et al. 1999; see Table 2). More generally, the column density of CH is observed to increase from 1.5 to 7.5$`\times `$10<sup>13</sup> cm<sup>-2</sup> for lines of sight with $`E_{\mathrm{B}\mathrm{V}}`$ increasing from 0.5 to 1.5 (Gredel et al. 1993). It would be of prime interest to obtain better data in this wavelength range to better constrain the molecule column densities.
### 4.4 Consequence of partial covering factor
From the detection of the Na i$`\lambda \lambda `$3303.3,3303.9 doublet, we can derive that the two components at $`z_{\mathrm{abs}}`$ = 1.0626 and 1.1801 arise in cold, dense and neutral gas (see Section 4.3). It is therefore possible that the dimension of the cloud is less than $``$1.9 kpc which is the separation of the lines of sight to the two brightest images at the redshift of the absorber assuming that the lensing object is at the same redshift. Indeed, large variations of Na i column density have been reported in the nearby interstellar medium on very small scales (Meyer & Lauroesch 1999). We have therefore investigated the impact on the column density measurements of the assumption that the cloud covers the three images. If it is the case that the cloud covers only one of the images the column density is larger. The discrepancy cannot be very large, however, as most of the lines used for column density determination are weak. The strong lines are completely saturated and blended, which prevents in any case any determination of the line parameters.
We have considered the Mg i and Na i lines in the $`z_{\mathrm{abs}}`$ = 1.1801 component. It can be seen on Fig. 8 that if only one image is covered, it cannot be the brightest as the residual normalized flux in the Mg i absorption is smaller than 0.5. As the flux ratio of the two brightest components is 1.2, we artificially placed the zero level at 0.4 on the scale of Fig. 8. Voigt profile fitting of the Na i doublet gives log $`N`$(Na i) = 13.7 and $`b`$ = 1 km s<sup>-1</sup> which is within a factor of two of what has been derived previously (see Table 2). Four components have been used to fit the Mg i blend at $`\mathrm{\Delta }v`$ = $``$60 km s<sup>-1</sup> (see Fig. 8). For the component at $`z_{\mathrm{abs}}`$ = 1.1801, we obtain log $`N`$(Mg i) = 13.8.
Note that in this case, the Na i and Mg i column densities are nearly identical to what is observed toward 23 Ori (see Table 2). We therefore conclude that $`N`$(Mg i)/$`N`$(Na i$``$ 1 is a robust measurement in this system.
Mg i is the line with the largest saturation among those used to derive column densities quoted in Table 2. Therefore, the column densities derived from weaker lines should not differ from what is quoted in Table 2 by more than a factor of two.
### 4.5 Metallicity and dust content
In the following we use the conventional definition \[X/H\] = log($`Z/Z_{}`$), with $`Z`$(X) the metallicity of species X relative to hydrogen. Ca ii and Ti ii are both detected at $`z_{\mathrm{abs}}`$ = 1.0613 and 1.0631. The column densities are consistent with what is seen in the interstellar medium of our Galaxy (Stokes 1978; see Table 2). However, log $`N`$(Ca ii)/$`N`$(Ti ii$``$ $``$0.5 and $``$0.2 at, respectively, $`z_{\mathrm{abs}}`$ = 1.0613 and 1.0631 when the relative solar metallicity is log $`Z_{}`$(Ca) $``$ log $`Z_{}`$(Ti) = 1.38. Various explanations for this discrepancy can be invoked, amongst them the most likely are: (i) Calcium is mostly in the form of Ca iii, (ii) Calcium is more depleted into dust-grains than Titanium. Note that the relative metallicities \[Ca/Fe\] and \[Ti/Fe\] are both observed to be $``$+0.3 for \[Fe/H\] $`<`$ $``$1 in late-type stars (Thévenin 1998). Note also that the Ca iii/Ca ii ratio derived in our Galaxy is in the range 5–10 which is much smaller than the discrepancy mentioned above. This favors the explanation that Calcium is heavily depleted into dust-grains (see below).
In the $`z_{\mathrm{abs}}`$ = 1.0631 component, Mn ii is also seen and log $`N`$(Mn ii)/$`N`$(Ti ii$``$ $`+`$0.8. As the solar metallicity of Mn is $``$6.47, the relative solar metallicity is log $`Z_{}`$(Mn) $``$ log $`Z_{}`$(Ti) = = $`+`$0.6. This is consistent with similar depletion of Mn and Ti as observed in warm gas (Savage & Sembach 1996). Therefore, in this component, we cannot rule out that the low Ca ii column density is due to ionization.
In the $`z_{\mathrm{abs}}`$ = 1.0613 component, Fe ii is also seen with log $`N`$(Fe ii)/$`N`$(Ti ii$``$ $`+`$2.1. The solar metallicity of iron and titanium are, respectively, $``$4.49 and $``$7.07 and the relative solar metallicity is log $`Z_{}`$(Fe) $``$ log $`Z_{}`$(Ti) = $`+`$2.58. There is no differential ionization correction for these two elements. The discrepancy, $``$+0.5 dex, between the two ratios can only be explained by a larger depletion of titanium compared to iron into dust-grains by $``$0.5 dex as in the cool gas of the ISM (Savage & Sembach 1996). Indeed, from nucleosynthesis alone, we would expect titanium to be enhanced compared to iron (Thévenin 1998) contrary to what is observed. We therefore conclude that depletion into dust-grains is present in this system. The low Ca ii column density can indeed be explained by a large depletion of calcium into dust-grains as is observed in the ISM.
The fact that Ti ii is not detected in the $`z_{\mathrm{abs}}`$ = 1.18 system is surprising, although the limit on the column density is not stringent and only a factor of four smaller than what is seen in the $`z_{\mathrm{abs}}`$ = 1.06 components.
The presence of dust is supported by the analysis of the column density ratios in the $`z_{\mathrm{abs}}`$ = 1.0626 and 1.1801 components where Na i absorption is detected. Indeed, we can compute for different elements the quantity $`N_{}`$ = 10<sup>\[X/H\]</sup>$`\times `$$`N`$(H i) which is the H i column density of a cloud with solar metallicity that would have the same column density of element $`X`$ as the observed one. We compute,
$$\mathrm{log}N_{}=[\mathrm{X}/\mathrm{H}]+\mathrm{log}N(\mathrm{HI})=\mathrm{log}\left(N(\mathrm{X}^i)\frac{\mathrm{X}}{\mathrm{X}^i}\frac{1}{\delta }\frac{1}{Z_{}(\mathrm{X})}\right)$$
(6)
where $`Z_{}`$ is the solar metallicity and 1-$`\delta `$ the fraction of the element tied up into dust-grains. We assume solar metallicity and depletion pattern given in Table 5 of Savage & Sembach (1996). Solar metallicities relative to hydrogen for Na, Fe, Ca and Ti are, respectively, $``$5.69, $``$4.49, $``$5.66 and $``$7.07. If we assume that the absorption arises in cool gas similar to the gas seen in front of $`\zeta `$Oph, we find that at $`z_{\mathrm{abs}}`$ = 1.1801, log $`N_{}`$ $`>`$20.44, $`<`$21.80, $`<`$21.18, $`<`$21.30 from Na, Fe, Ca and Ti and at $`z_{\mathrm{abs}}`$ = 1.0626, log $`N_{}`$ $`>`$20.85, = 20.86, $`<`$22.10 from Na, Fe and Ti respectively.
These consistent results suggest that, in the $`z_{\mathrm{abs}}`$ = 1.0626 and 1.1801 components, log $`N`$(H i) + \[X/H\] $``$ 21 and the depletion into dust-grains is similar to what is seen in cool Galactic interstellar clouds. If the relation of Bohlin et al. (1978) holds, this implies a color excess $`E_{\mathrm{B}\mathrm{V}}`$ $``$ 0.2 with the only assumption that the overall dust-to-metal ratio does not depend on metallicity. At the wavelength of the Mg ii absorption, ($`\lambda _{\mathrm{z}=1}`$ $``$ 3000 Å; $`\lambda _{\mathrm{rest}}^{\mathrm{QSO}}`$ $``$ 1200 Å; $`\nu _{\mathrm{rest}}^{\mathrm{QSO}}`$ $``$ 2$`\times `$10<sup>15</sup> Hz) this would induce an extinction of about 1 mag (Seaton 1979). At $`\nu _{\mathrm{rest}}^{\mathrm{QSO}}`$ $``$ 10<sup>14</sup> Hz, the extinction would be negligible. Note that there is some evidence that $`\nu F_\nu `$ decreases by a factor of two from 10<sup>15</sup> to 10<sup>14</sup> Hz in the APM 08279+5255 SED (Lewis et al. 1998).
The amount of dust suggested by the previous discussion is significant. We have therefore searched the spectrum of the quasar for some signature of this amount of dust. For this, we have compared the spectrum of APM 08279+5255 with the composite QSO spectrum obtained with the FOS-HST (Zheng et al. 1997) attenuated by dust with optical depth $`\tau _{\mathrm{dust}}`$($`\lambda `$) at the observed wavelength $`\lambda `$,
$$\tau _{\mathrm{dust}}(\lambda )=k\left[\frac{N}{10^{21}\mathrm{cm}^2}\right]\xi \left(\frac{\lambda }{1+z}\right)$$
(7)
where $`\xi (\lambda )`$ is the ratio of the extinction at the wavelength $`\lambda `$ to that in the B-band as observed in our Galaxy, $`k`$ = 10<sup>21</sup> cm<sup>-2</sup> $`\tau _\mathrm{B}`$/$`N`$(H i) is the dimensionless dust-to-gas ratio and $`N`$ the H i column density (e.g. Pei et al. 1991, Srianand & Kembhavi 1997). We assume here log $`N`$(H i) = 21. The redshift of the QSO is taken to be $`z_{\mathrm{QSO}}`$ = 3.91 and that of the absorber $`z_{\mathrm{abs}}`$ = 1.1. The results are shown in Fig. 11 where the spectrum of APM 08279+5255 as observed over part of the $`R`$-band (solid line) is plotted together with the composite HST-FOS QSO spectrum attenuated by dust with optical depth in the B-band $`\tau _{\mathrm{dust}}`$(B) = 0.1, 0.2 and 0.3. It is apparent that the best fit is obtained with $`\tau _{\mathrm{dust}}`$(B) $``$ 0.3. In our Galaxy, this would correspond to about $`E_{\mathrm{B}\mathrm{V}}`$ $``$ 0.1. This is two times smaller than what has been derived above. This suggests that the dust to metal ratio in the redshifted gas is about half that in our Galaxy.
Altogether we find that the H i column density at $`z`$ = 1 is of the order of $`1\times 10^{21}`$ cm<sup>-2</sup> to $`5\times 10^{21}`$ cm<sup>-2</sup>, the corresponding metallicity is in the range 1–0.3 $`Z_{}`$, the dust-to-metal ratio is about half that in our Galaxy and the relative depletion of species into dust-grains is similar to what is observed in cool gas in the disk of our Galaxy.
The colors of the images derived from HST imaging are nearly identical (Ibata et al. 1999). The differential reddening over kpc scales is smaller than 10%. This indicates that, although the medium is highly inhomogeneous, the extinction is fairly uniform over distances of the order of the separation between the lines of sight. It must be noted that the number of components in the two strong Mg ii systems must be quite large. Indeed, the Mg ii absorptions reach the zero level over $``$200 and 350 km s<sup>-1</sup> at, respectively, $`z_{\mathrm{abs}}`$ = 1.181 and 1.062. One single component cannot be much larger than about $``$10 km s<sup>-1</sup> and therefore the number of components is larger than 20 and 40 respectively along all three lines of sight. This is probably larger than what is seen in the disk of our Galaxy (Sembach & Danks 1994). Therefore, it may well be possible that the total reddening we see is not due to one strong component only but rather is due to the accumulated effect of a large number of diffuse clouds with small and similar extinctions. As the number of clouds is large and similar along the different lines of sight, the differential extinction is small. To probe this, higher S/N ratio data should be obtained in the region of the Na i absorption to investigate what is the velocity spread of this absorption.
### 4.6 Nature of the systems
Churchill (1999) has shown that Mg ii systems at intermediate redshift can be classified in five categories: DLA, Double, Classic, C iv deficient and Weak. The first class is characterized by strong Mg ii absorption saturated over $``$150 km s<sup>-1</sup> and, when observed, the Lyman-$`\alpha `$ line is damped. The Double systems are characterized by kinematic velocity spreads up to 400 km s<sup>-1</sup>. Other classes correspond to much weaker systems.
The system at $`z_{\mathrm{abs}}`$ = 1.181 would be classified as Double (see Fig. 8): it has Mg ii$`\lambda `$2796 absorption saturated over more than 150 km s<sup>-1</sup> in total. The characteristic double profile is very similar to that of the system at $`z_{\mathrm{abs}}`$ = 1.17 toward Q 0450–132 (see Petitjean et al. 1994). The system at $`z_{\mathrm{abs}}`$ = 1.062 would be classified as DLA, as the Mg ii absorption is continuously saturated over more than 300 km s<sup>-1</sup> (see Fig. 7). To our knowledge, this system has one of the strongest Mg ii absorption features known ($`W_\mathrm{r}`$ $``$ 3.3 Å). As seen from Fig. 7, the Ti ii absorption spans $``$ 300 km s<sup>-1</sup> and coincides exactly with the saturated part of the Mg ii absorption.
At high-z, Prochaska & Wolfe (1998) have shown that most of the low-ionization absorptions associated with DLAs, have an edge-leading profile, possibly revealing large-scale rotational motions. They conclude that DLA systems arise in large rotating disks. Haehnelt et al. (1998) have claimed that the observed profiles can be reproduced as well if the line of sight passes through several interacting blobs. Indeed, Ledoux et al. (1998a) have shown that the observed profiles are consistent with rotation when they span less than $`\mathrm{\Delta }V`$ $``$ 150 km s<sup>-1</sup>. For larger velocity spreads, several sub-systems are usually seen.
The characteristic edge-leading profile is seen in the $`z_{\mathrm{abs}}`$ = 2.974 system (see below, Fig. 12) but not in the $`z_{\mathrm{abs}}`$ = 1.062 and 1.181 Mg ii systems. The large spread of the profiles ($`>`$ 200 and 300 km s<sup>-1</sup>) is more reminiscent of the profile seen toward the supernova SN 1993J in the large nearby spiral galaxy M 81. The latter is part of a complex interacting group together with M 82, NGC 3077 and NGC 2976 with tidally stripped H i linking individual galaxies over an area of $``$50$`\times `$100 kpc<sup>2</sup> (Yun et al. 1994). The Mg ii absorption is characteristic of the Double class as defined by Churchill (1999). It is spread over $``$400 km s<sup>-1</sup> with two strong saturated absorptions of width, respectively, $`\mathrm{\Delta }V`$ $``$ 150 and 90 km s<sup>-1</sup>, and separated by $``$180 km s<sup>-1</sup> (Bowen et al. 1995). Most of the absorption is due to tidal debris expelled outside the disks of the interacting galaxies.
In the case of the $`z_{\mathrm{abs}}`$ = 1.062 system, the Mg ii absorption does not show any sub-structure. Such strong absorption is expected to occur in the central part of galaxies. Although the statistics are very poor, it seems that when strong absorption occurs, the equivalent width of the absorption is anti-correlated with the impact parameter between the line of sight and the center of the galaxy (Bowen et al. 1995). The impact parameter could be as low as 1 kpc for a system with $`W_\mathrm{r}`$ $``$3 Å. The separation between the two bright images of APM 08279+5255 is $``$1.9$`h_{75}^1`$ kpc at $`z`$ $``$ 1. This suggests that the object giving rise to the $`z_{\mathrm{abs}}`$ = 1.06 Mg ii system could be nearly exactly aligned with the quasar.
## 5 The system at $`z_{\mathrm{abs}}`$ = 2.974
There is a strong absorption feature with $`W_{\mathrm{obs}}`$ $`>`$ 19 Å at $`\lambda `$ = 4831 Å. It corresponds to H i Lyman-$`\alpha `$ at $`z_{\mathrm{abs}}`$ = 2.974. The only possibility of coincidence with a BAL transition could be Lyman$`\beta `$ at $`z_{\mathrm{abs}}`$ = 3.71 but there is no corresponding Lyman$`\alpha `$ transition (see Srianand & Petitjean 2000). Although the red-wing of the absorption has the characteristic shape of a damped transition, uncertainties in the continuum determination prevent an accurate determination of the column density. Associated absorptions from Al ii, Fe ii, Si ii, C ii and O i are detected in four components spanning $``$100 km s<sup>-1</sup> (see Fig. 12). By fitting the Lyman-$`\alpha `$ line, we estimate that the total H i column density in the four components is in the range 19.8 $`<`$ log $`N`$(H i$`<`$ 20.3.
It is important to note that there is no evidence that the cloud does not cover the three lines of sight which are separated by $``$ 200$`h_{75}^1`$ pc at the redshift of the absorber. Indeed, the core of the H i absorption is black over $``$ 15 Å. Column densities integrated over the absorption profiles have been obtained for all species and summarized in Table 3. Column #3 of Table 3 gives the abundance of the element assuming that log $`N`$(H i) = 20.3 in the cloud and that the observed ion is the dominant species. Given the uncertainty in the neutral hydrogen column density, the absolute values are unreliable and could be 0.5 dex higher. Column #4 and #5 of Table 3 give the solar metallicity and the metallicity relative to solar respectively. It is remarkable how consistent the measurements are, which point toward metallicities less than 10<sup>-1.5</sup> $`Z_{}`$. The low metallicities are not due to depletion into dust-grains. Indeed, the relatively small neutral hydrogen column density implies that column densities of relatively abundant elements can be measured. It is apparent that iron is not depleted compared to carbon or oxygen and that the amount of dust in this cloud must be very small.
Metallicity in damped Lyman-$`\alpha `$ systems is usually measured using zinc, an element that is not very much depleted into dust-grains in our Galaxy and, because it has relatively low metallicity compared to other elements, induces non-saturated absorptions even for clouds of high hydrogen column density. There is barely no evolution in the measured Zinc metallicity from $`z`$ $``$ 1 to $`z`$ $``$ 3 (Pettini et al. 1997, 1999). At $`z`$ $`>`$ 3, in most damped Lyman-$`\alpha `$ studied up to now, zinc is not detected. This is most probably a consequence of limited S/N ratio of the data however. Indeed, the detection limit of most of the spectra is log $`N`$(Zn ii$``$ 11.5 which means \[Zn/H\] $`<`$ $``$1.4 (see e.g. Prochaska & Wolfe 1997). It can be noted that in the system at $`z_{\mathrm{abs}}`$ = 2.974 toward APM 08279+5255, the limit on zinc is of this order. However, as we can measure metallicities for more abundant elements, we know that metallicities are less than $``$1.5. For such abundances, zinc would have been detectable in the spectrum of APM 08279+5255 only for H i column densities larger than 10<sup>21</sup> cm<sup>-2</sup>. It is therefore possible that the upper limit found for the zinc metallicity at $`z`$ $`>`$ 3 in previous surveys is indicative of a true evolution of the metallicity in individual systems (see also Prochaska & Wolfe 2000; Savaglio et al. 1999). This should be checked by measuring in the same systems abundances of species like carbon, aluminium, silicon and iron.
## 6 Conclusion
The doublet ratio of several intervening Mg ii$`\lambda \lambda `$2796,2803 systems along the line of sight to APM 08279+5255 is observed close to unity, indicating saturation of the lines, whereas the depth of the lines is close to 0.5 in the normalized spectrum (see Fig. 1). This can be understood if the absorption profile is made of components with arbitrarily small Doppler parameters ($`b`$ $``$ 1–1.5 km s<sup>-1</sup>). This would imply however a surprisingly low temperature (1500–3000 K) when the gas is expected to be heated by photo-ionization to temperatures larger than 10<sup>4</sup> K (e.g. Petitjean et al. 1992). A more likely explanation of these observations is that Mg ii galactic halos are composed of a collection of clouds each of them having dimensions less than $``$1 kpc. Individual clouds, regularly spread over the velocity profile by kinematics, cover only one of the two brighest image of the lensed quasar. The number density of clouds is not large enough for the absorption material to cover the two lines of sight at all velocity positions. The total covering factor of the halo however is close to one, consistent with observations of associated galaxies at intermediate redshift. In contrast, the two strong Mg ii systems at $`z_{\mathrm{abs}}`$ = 1.06 and 1.18 have covering factor equal to one (the lines are saturated and go to the zero level) over more than 200 km s<sup>-1</sup>. These latter systems are likely to arise due to absorption through the central part of galaxies where the number of clouds is so large that saturated absorption occurs along both lines of sight whatever the radius of the individual clouds might be.
Models by Mo & Miralda-Escudé (1996) have shown that halos with small rotation velocity ($`<`$ 100 km s<sup>-1</sup>) should contribute little to the total cross-section of Mg ii systems as they have dimensions as small as 5 kpc. This conclusion is probably true at $`z`$ $`<`$ 1 (see also Churchill et al. 1996). It is however interesting to note that most of the $`z`$ $`>`$ 1 Mg ii systems studied here (see Fig. 1) are spread over much less than 100 km s<sup>-1</sup>. Moreover, they have equivalent widths $`W_\mathrm{r}`$ $``$ 0.43, 0.37, 0.12, 0.99, 0.28 and 0.17 Å at $`z_{\mathrm{abs}}`$ = 1.211, 1.5497, 1.5523, 1.813, 2.0418 and 2.0668 respectively. Therefore the systems we see have similar strengths as the systems which, at lower redshift, are associated with large halos of galaxies. In particular, they are generally stronger than the weak Mg ii systems studied by Churchill et al. (1999). This means that, contrary to what is seen at lower redshift, these systems could be associated with halos of low rotation velocity ($`<`$ 100 km s<sup>-1</sup>) and thus small radii.
The two strong Mg ii systems at $`z_{\mathrm{abs}}`$ = 1.06 and 1.18 are studied in detail. Absorption from Ca ii, Mg i, Ti ii, Mn ii and Fe ii have been observed in several damped Lyman-$`\alpha `$ systems over a large range of redshift (Meyer et al. 1995, Lu et al. 1996, Vladilo et al. 1997, Proschaska & Wolfe 1997, Churchill et al. 2000). This is, however, the first time that Na i is also detected at such redshift, thanks to the combination of high S/N ratio and high spectral resolution. This additional strong constraint shows that the gas in these systems is cool and neutral. Indeed, similar column densities are observed in our Galaxy for Ca ii, Mg i, Ti ii, Mn ii and Fe ii in warm and cool gas clouds toward, respectively, $`\mu `$Col and 23 Ori. Only the Na i column density differs; it is more than an order of magnitude larger through the cool cloud. Doppler parameters as low as $`b`$ $``$ 1 km s<sup>-1</sup> are derived from Voigt-profile fitting of isolated subcomponents. We find that the H i column density at $`z`$ = 1 is of the order of $`1\times 10^{21}`$ cm<sup>-2</sup> to $`5\times 10^{21}`$ cm<sup>-2</sup>, the corresponding metallicity is in the range 1–0.3 $`Z_{}`$, the dust-to-metal ratio is about a third that in our Galaxy and the relative depletions of iron, titanium, manganese and calcium are similar to those in cool gas in the disk of our Galaxy. The dust-to-metal ratio measured here is similar to what is derived in most of the damped Lyman-$`\alpha `$ systems (Vladilo 1998, Savaglio et al. 1999). The presence of dust is supported by the reddening of the QSO spectrum over the $`R`$-band. These are probably amongst the most metal and dust-rich damped Lyman-$`\alpha `$ systems at $`z`$ $``$ 1. The dust depletion pattern is similar to that observed in cool gas in the Galaxy. All this is consistent with the finding by Petitjean et al. (1992) that although most of the damped Lyman-$`\alpha `$ systems arise in warm gas, the highest column densities are due to a collection of clumps that condense out of the warm phase due to thermal instability (see also Lane et al. 2000).
Another damped Lyman-$`\alpha `$ system is seen at $`z_{\mathrm{abs}}`$ = 2.974 with 19.8 $`<`$ log $`N`$(H i$`<`$ 20.3. As the Lyman-$`\alpha `$ line is black over about 15 Å, the cloud must cover the three QSO images. The transverse dimension of the absorber is therefore larger than 200 $`h_{75}^1`$ pc. Column densities of Al ii, Fe ii, Si ii, C ii and O i indicate abundances relative to solar of $``$2.31, $``$2.26, $``$2.10, $``$2.35 and $``$2.37 for, respectively, Fe, Al, Si, C and O (for log $`N`$(H i) = 20.3). Metallicities are therefore less than 10<sup>-1.5</sup> $`Z_{}`$ and, if any, the amount of dust in the cloud is very small, as are any deviations from relative solar abundances. It seems likely that the difficulty to detect Zinc in several damped Lyman-$`\alpha `$ systems at $`z`$ $`>`$ 3 in previous surveys is indicative of a true cosmological evolution of the metallicity in individual systems (see also Prochaska & Wolfe 2000, Savaglio et al. 1999). This should be checked by measuring in the same systems abundances of species like carbon, aluminium, silicon and iron.
###### Acknowledgements.
We thank the team headed by Sara L. Ellison to have made this beautiful data available for general public use and particularly Sara L. Ellison for an access to individual spectra. We gratefully acknowledge support from the Indo-French Centre for the Promotion of Advanced Research (Centre Franco-Indien pour la Promotion de la Recherche Avancée) under contract No. 1710-1. PPJ thanks Elisabeth Flam for useful discussions and Patrick Boissé for a critical reading of the manuscript.
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# Untitled Document
An Accelerated Expansion Model in the Absence of the Cosmological Constant
Yi-Ping Qin<sup>1,2,3</sup>
<sup>1</sup> Yunnan Observatory, Chinese Academy of Sciences, Kunming, Yunnan 650011, P. R. China; E-mail: ypqin@public.km.yn.cn
<sup>2</sup> National Astronomical Observatories, Chinese Academy of Sciences
<sup>3</sup> Chinese Academy of Science-Peking University joint Beijing Astrophysical Center
Summary
> Based on some observations, the apparent energy, associated with gravity, of vacuums is defined, with that of normal vacuums to be zero and that of the vacuums losing some energy to be negative. An important application of the energy is its contribution to Einstein’s equation. A cosmological model, accounting for recent observations of the accelerated expansion of the universe, in the absence of the cosmological constant, can be well constructed. In a certain case, the expansion of the universe would be decelerated at its early epoch and accelerated at its late epoch. The curvature of the universe would depend on the ratio of matter energy to total energy. The missing mass problem does no longer exist in this model. Most negative apparent energy vacuums might be contained in voids, then the spacetime of galaxy clusters or that of the solar system would not be significantly affected by this kind of energy.
> PACS number: 98.80.Bp, 98.80.Dr, 98.80.Es, 98.80.Ft
Recent observations showed that the expansion of the universe is accelerated rather than decelerated \[1–3\]. An economic approach to this phenomenon is to adopt the cosmological constant which is often referred to the vacuum energy produced by the phase transitions the universe undergoes as it cools. However, there are some reasons against this scenario . For example, the amount of vacuum energy produced by all the phase transitions can be about $`10^{120}`$ times greater than the density of all the matter in the universe ; in particular, the constant corresponds to a universal energy density but its influence on the nearby spacetime has never been observed. These facts suggest that vacuum energy acts as something like potential energy and the apparent energy associated with gravity, of the vacuums in the nearby spacetime might be zero.
Recently, the Casimir force was detected in laboratory . It is reasonable that a work done by the force would extract more or less energy from vacuums. We call a vacuum not losing any energy a normal vacuum and that losing some a deficit vacuum. According to the above comprehension, we define the apparent energy associated with gravity, of normal vacuums to be zero and that of deficit vacuums to be negative. The energy is assumed to contribute to Einstein’s equation the way matter energy does.
Now we consider a cosmological model of the Robertson-Walker metric following Einstein’s equation and the conservation equation of the energy-momentum tensor. The difference is that, we take
$$\rho =\rho _m+\rho _v$$
(1)
and
$$p=p_m+p_v$$
(2)
with that $`\rho _v`$ can be negative, where $`m`$ denotes matter and $`v`$ represents vacuums.
For the Robertson-Walker metric, Einstein’s equation gives
$$3\stackrel{}{R}=4\pi G(\rho +3p)R,$$
(3)
$$R\stackrel{}{R}+2\stackrel{2}{\stackrel{}{R}}+2k=4\pi G(\rho p)R^2,$$
(4)
and the conservation equation of the energy-momentum tensor yields
$$\stackrel{}{p}R^3=\frac{d}{dt}[R^3(\rho +p)],$$
(5)
where $`\stackrel{}{R}=dR/dt`$. From (3) and (4) one can obtain
$$\stackrel{2}{\stackrel{}{R}}+k=\frac{8\pi G}{3}\rho R^2$$
(6)
and
$$2R\stackrel{}{R}+\stackrel{2}{\stackrel{}{R}}+k=8\pi GpR^2.$$
(7)
Let us define
$$H\frac{\stackrel{}{R}}{R},$$
(8)
$$q\frac{R\stackrel{}{R}}{\stackrel{2}{\stackrel{}{R}}},$$
(9)
$$\rho _c\frac{3H^2}{8\pi G},$$
(10)
and
$$\mathrm{\Omega }\frac{\rho }{\rho _c}.$$
(11)
Then equations (6) and (7) can be written as
$$\frac{k}{R^2}=H^2(\mathrm{\Omega }1)$$
(12)
and
$$\frac{k}{R^2}=H^2(2q1\frac{3p}{\rho }\mathrm{\Omega }),$$
(13)
respectively.
The observation of $`q<0`$ suggests that, at least at the present time, the following condition must be satisfied (see equation (3)):
$$\rho +3p<0.$$
(14)
One means to meet condition (14) is to consider a universe containing both deficit and normal vacuums (then on the average, $`\rho _v<0`$), and to assume that deficit vacuums act as negative energy photons (then $`p_v=\rho _v/3<0`$). (According to quantum electrodynamics, normal vacuums are full of all kinds of electromagnetic modes.)
At the late epoch of the universe, the pressure of matter particles is negligible. Then $`p=p_v`$, $`3p=3p_v=\rho _v`$. Condition (14) leads to $`\rho <\rho _v`$, which allows $`\rho >0`$ (note $`\rho _v<0`$). Hence, it is possible that a positive energy density (where $`\rho _m>\rho _v`$) may lead to an accelerated expansion of the universe at its late epoch so long as $`\rho <\rho _v`$ or $`\rho _m<2\rho _v`$. For $`\rho >0`$, we find $`\rho _v<\rho _m<2\rho _v`$. Equation (13) leads to $`k/R^2=H^2(2q1\alpha \mathrm{\Omega })`$, where
$$\alpha \frac{\rho _v}{\rho }.$$
(15)
This together with (12) yield $`\mathrm{\Omega }=2q/(1+\alpha )`$. As $`\rho <\rho _v`$, for $`\rho >0`$, we find $`\alpha <1`$. For $`q<0`$, this indicates that $`\mathrm{\Omega }>0`$. The curvature of the universe depends on $`\alpha `$: when $`\alpha >(12q)/\mathrm{\Omega }`$, then $`k=1`$ and $`0<\mathrm{\Omega }<1`$; when $`\alpha =(12q)/\mathrm{\Omega }`$, then $`k=0`$ and $`\mathrm{\Omega }=1`$; when $`\alpha <(12q)/\mathrm{\Omega }`$, then $`k=+1`$ and $`\mathrm{\Omega }>1`$. Equation (15) shows, if $`\alpha `$ is a constant, the sign of $`\rho `$ will remain unchanged. For a flat universe, $`\alpha =(12q)`$ at the late epoch. When $`\alpha `$ is known (e.g., determined by the late epoch acceleration), $`\rho `$ and $`\rho _v`$ may be known since $`\rho _m`$ is measurable.
At the early epoch of the universe, $`p_m=\rho _m/3`$. Then $`p=\rho /3`$ (as $`p_v=\rho _v/3`$). For $`\rho >0`$, we find from (3) that $`\stackrel{}{R}<0`$ (then $`q>0`$), indicating that the expansion of the universe is decelerated. Equation (13) leads to $`k/R^2=H^2(2q1\mathrm{\Omega })`$. This together with (12) yield $`\mathrm{\Omega }=q`$. Since $`q>0`$, then $`\mathrm{\Omega }>0`$. When $`k=1`$, then $`2q1<\mathrm{\Omega }<1`$, $`0<q<1`$; when $`k=0`$, then $`\mathrm{\Omega }=1`$, $`q=1`$; when $`k=+1`$, then $`1<\mathrm{\Omega }<2q1`$, $`q>1`$.
We find in this model that the universe can possess a positive energy density ($`\rho >0`$, where $`\rho _m>\rho _v`$) and a positive energy parameter ($`\mathrm{\Omega }>0`$). The curvature of the universe depends on the ratio of vacuum energy (or matter energy) to total energy. As the pressure of matter particles becomes less important as time goes on, the expansion of the universe will change. When $`p_m=((\alpha +1)/3(\alpha 1))\rho _m`$, we find from (3) that $`\stackrel{}{R}=0`$. It is at this moment the expansion changing from deceleration to acceleration.
It is obvious that deficit vacuums, as they possess negative energy, will be expelled by gravitation. Then there will seldom be deficit vacuums remained in galaxy clusters, and therefore the spacetime of galaxy clusters and that of the solar system will not be significantly affected. We suspect that the place in the universe containing most deficit vacuums might be voids. In the above model, deficit vacuums are assumed to act as negative energy photons. These photons will be deflected towards the center of voids and then the voids might be crowded with them. Within the voids, if the amount of vacuum apparent energy (negative) is over that of matter energy, the spacetime might be somewhat like that of the Schwarzschild solution with (more or less) negative mass, and then matter objects would experience (strong or weak) anti-gravitation. Many matter objects must have been driven to the sheets around the voids, and the number of those remained will be small. In addition, any matter objects would undergo the negative pressure ($`p(void)p_v(void)=\rho _v(void)/3<0`$ ) in the voids. Those survived would be that with their components being firmly connected. The cloud structure of matter must finally be disintegrated and some other matter structures must be reduced. As a result, the amount of matter in the sheets would grow and that in the voids would reduce (the disintegrated matter would be easier to be expelled to the sheets by anti-gravitation). In this model, we can assign the observed matter density to be $`\rho _m`$ (e.g., at the present time, taking $`\rho _m=0.3\rho _c`$ ), then the missing mass problem will no longer exist. (We suggest that many conventional problems should be reexamined in this new model.)
It might be possible that deficit vacuums act as negative mass particles, or some act as negative mass particles while others act as negative energy photons. Situations in these cases will be different.
ACKNOWLEDGEMENTS
The author is grateful to Professors G. Z. Xie, Xue-Tang Zheng and Shi-Min Wu for their guide and help. This work was supported by the United Laboratory of Optical Astronomy, CAS, the Natural Science Foundation of China, and the Natural Science Foundation of Yunnan.
> REFERENCES
>
> 1. Garnavich, P. M. et al. (1998). Astrophys. J. 493, L53.
>
> 2. Perlmutter, S. et al. (1998). Nature 391, 51.
>
> 3. Riess, A. G. et al. (1998). Astron. J. 116, 1009.
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> 4. Coles, P. (1998). Nature 393, 741.
>
> 5. Weinberg, S. (1989). Rev. Mod. Phys. 6, 1.
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> 6. Lamoreaux, S. K. (1997). Phys. Rev. Lett. 78, 5.
>
> 7. Weinberg, S. (1972). In Gravitation and Cosmology, (John Wiley, New York).
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> 8. Trimble, V., and Aschwanden, M. (1999). Publ. Astron. Soc. Pac. 111, 385.
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# Hidden Order in the Cuprates
## I Introduction
In this paper we argue that much of the strange phenomenology of the cuprate superconductors may be simply explained as the disorder-frustrated development of a new order parameter. There are a number of potential candidates for this order, but the one we favor on phenomenological grounds is orbital antiferromagnetism or $`d`$-density wave (DDW) order , which is characterized by a local order parameter which distills the universal physics underlying the staggered flux state divorced from the uncontrolled approximations associated with the gauge theory formalism. The essence of our idea is that the pseudogap observed in underdoped cuprates is an actual gap in the one-particle excitation spectrum at the wavevector $`(\pi ,0)`$ and symmetry-related points of the Brillouin zone associated with the development of this new order. It is “pseudo” in experiment only because of extreme sensitivity to sample imperfection caused by proximity to the phase transition. Moreover, the DDW couples weakly to common experimental probes, and is thus difficult to detect.
Our proposal has much in common with theoretical ideas already in the literature , and borrows heavily from them. For example, Wen and Lee have proposed staggered currents that fluctuate but do not order . Varma has proposed currents which alternate in the unit cell but do not break translational symmetry . Emery and Kivelson and Caprara et al. have proposed states with broken symmetries of different kinds. Our strategy for constructing a theory and confronting experiments differs from most others in deemphasizing modeling of the “strange metal” behavior and focusing on order, low-temperature phenomenology, and material imperfection - all issues with sharp experimental dichotomies amenable to falsification. DDW order can be detected if it is present. If it is not present, the proposal is disproved.
## II Competing Order
Order-parameter competition has always been a natural candidate for explaining why the superconducting transition temperature $`T_c`$ first grows and then retreats as doping is reduced. Let us consider the generic zero-temperature Ginzburg-Landau free energy
$$F=\lambda (y^2+|z|^2)^2+\gamma y^2|z|^2\alpha y^2\alpha ^{}|z|^2,$$
(1)
describing the development of order parameters $`y`$ and $`z`$ in the case that low-order mixing is forbidden by symmetry. In Fig. 1b we plot the values of $`y`$ and $`z`$ that minimize $`F`$ for the case of $`\lambda =1`$ and $`\gamma =0.8`$ as a function of the abstract tuning parameter $`p`$. The variables $`\alpha `$ and $`\alpha ^{}`$ are the simple linear functions of $`p`$ shown in Fig. 1c. One sees that $`z`$ develops at $`p=0.3`$, $`y`$ develops at $`p=0.2`$, and that $`0.1<p<0.2`$ is a coexistence region in which the growth of $`y`$ suppresses and eventually eliminates $`z`$. Thus if we imagine $`z`$ to be the magnitude of the order parameter for $`d`$-wave superconductivity and $`p`$ to be doping, then we can understand the onset, growth, saturation, and eventual destruction of superconductivity with reduced doping as an effect of a monotonically strengthening $`d`$-wave pairing interaction, as opposed to one that first strengthens and then weakens. The underdoped side of the superconducting dome is then fundamentally different from the overdoped side in that the superfluid density is suppressed there by the development of a second order parameter $`y`$.
## III D-Density Wave
Let us now consider the order parameter
$$y=i\underset{𝐤,s}{}f(𝐤)<c_{𝐤+𝐐,s}^{}c_{𝐤,s}>,$$
(2)
where $`f(𝐤)=\mathrm{cos}(k_x)\mathrm{cos}(k_y)`$. If $`f(𝐤)`$ were replaced by a function with $`s`$-wave symmetry, $`y`$ would simply be the order parameter of a charge-density wave (CDW) - hence, we call this state a $`d_{x^2y^2}`$ density wave state (DDW) . For the particular case of $`𝐐=(\pi ,\pi )`$, which we think
most relevant to the cuprates, the equivalence of $`𝐐`$ and $`𝐐`$ enforced by the underlying band structure requires the sum to be imaginary. Thus this state necessarily breaks parity and time-reversal symmetry (i.e. exhibits magnetism), as well as translation by one lattice spacing and rotation by $`\pi /2`$. It is, however, symmetric under the combination of any two of these operations. The order parameter is equivalent to the array of bond currents illustrated in Fig. 2.
The excitation spectrum of the DDW at very low energies is generic and consists of conventional fermionic particles and holes in a band structure like that of the $`d`$-wave superconductor with which it competes. Introducing a mean-field ansatz (cf. Eq. 2) we obtain the 1-body Hamiltonian
$$=\underset{𝐤,\sigma }{}ϵ(𝐤)c_{𝐤\sigma }^{}c_{𝐤\sigma }+\mathrm{\Delta }(𝐤)c_{𝐤\sigma }^{}c_{𝐤+𝐐\sigma }$$
(3)
where $`ϵ_𝐤=2t[\mathrm{cos}(k_x)+\mathrm{cos}(k_y)]`$ and $`\mathrm{\Delta }_𝐤=yV[\mathrm{cos}(k_x)\mathrm{cos}(k_y)]`$, and $`V`$ is a coupling constant in the microscopic Hamiltonian. Microscopic Hamiltonians with short-range repulsion and superexchange are favorable for such order but are even more favorable for an antiferromagnetic state. However, correlated hopping terms tend to tip the balance in favor of DDW order . Since the ordering occurs at $`𝐐=(\pi ,\pi )`$, it is most favorable at half-filling or low doping. The corresponding band structure is
$$E_𝐤=\pm \sqrt{ϵ_𝐤^2+|\mathrm{\Delta }_𝐤^{\mathrm{DDW}}|^2}.$$
(4)
At half-filling there are gapless quasiparticles only at the nodal points $`𝐤=(\pm \pi /2,\pm \pi /2)`$. At finite doping, Fermi pockets are opened, as shown in Fig. 2. While the DDW state is semimetallic at half-filling, it is a conventional metal (with 2-d localization prevented by interlayer tunneling) with a disconnected Fermi surface at dopings other than half-filling. It is possible for the DDW to discommensurate, thereby opening a full gap, as occurs with a traditional spin density wave, but this is not automatic because the remaining Fermi surface is not nested. Some related density-wave states are discussed in the Appendix.
The excitation spectrum at high energies is not generic. There is no reason for the quasiparticle at $`(\pi ,0)`$ to have integrity, particularly if the system is near the continuous quantum phase transition at $`p=0.2`$ in Fig. 1. This is a Fermi-surface reconnection, at which the Hall conductance jumps, a van Hove singularity develops at $`(\pi ,0)`$, and quasiparticles scatter violently even at low energy scales .
## IV D-wave Superconductivity
The Heisenberg exchange nominally responsible for DDW order also tends to favor $`d`$-wave superconductivity. This is the underlying reason the band structures of the two are so similar, and why the competition of these two kinds of order is natural. If we allow the superconducting bond expectation value $`<c_jc_k>=\pm z`$ to develop, where the sign is positive on $`x`$ bonds and negative on $`y`$ bonds, the Hartree-Fock Hamiltonian becomes
$$_{\mathrm{HF}}^{}=_{\mathrm{HF}}+J\underset{<jk>}{}\pm (zc_k^{}c_j^{}+z^{}c_jc_k),$$
(5)
and the corresponding superconducting quasiparticle dispersion relation becomes
$$E_𝐤=\pm \sqrt{[(ϵ_𝐤^2+|\mathrm{\Delta }_𝐤^{\mathrm{DDW}}|^2)^{1/2}\pm \mu ]^2+|\mathrm{\Delta }_𝐤^{\mathrm{DSC}}|^2},$$
(6)
where $`\mathrm{\Delta }_𝐤^{\mathrm{DSC}}=zJ[\mathrm{cos}(k_x)\mathrm{cos}(k_y)]`$ and $`\mu `$ is the chemical potential. Thus not only does this kind of interaction stabilize both kinds of order, it allows the two order parameters to evolve continuously into each other without collapsing the quasiparticle gap at the zone face. This allows us to use the ground state expectation value of $``$ and similar Hamiltonians as a sensible model for the energy functional $`F`$, i.e. one that does not throw away important low-energy excitations.
This calculation illustrates an important feature of the mixed state that the superfluid density is not fixed by sum rules on the underlying Fermi surface but is rather determined by the balance between the DDW and DSC order parameters. This is because the superfluid is primarily a condensate of Cooper pairs drawn from the gapped region near $`(\pi ,0)`$ rather than the residual Fermi surface near $`(\pi /2,\pi /2)`$. This effect is not difficult to understand if the DDW order parameter is small, for then the semimetallic state with Fermi points - or, away from half-filling, the conventional metallic state with a small, disconnected Fermi surface - is not significantly different from the parent metal with a full Fermi surface at the energy scales relevant to superconductivity. As the DDW order parameter becomes large, however, we have more and more the case of a powerful attractive force exciting electrons and holes virtually into the insulating part of the band structure and then binding these into superfluid. The result is a condensate fraction that falls precipitously as the DDW order parameter grows. An insulating ground state (or, in this case, nearly insulating, since there is a small disconnected Fermi surface) that becomes a superfluid without first becoming a metal is unusual in solids, but perhaps not in nature, for this is the central idea behind Higgs condensation in electroweak theory.
## V S-wave Competition
The competition between DDW and DSC has a simple analogue in the $`s`$-wave case that is particularly instructive because it is exact . The Hubbard model
$$=t\underset{<jk>}{}\underset{\sigma }{}c_{j\sigma }^{}c_{k\sigma }+U\underset{j}{}c_j^{}c_j^{}c_jc_j$$
(7)
has the special property at half-filling that replacing the fermion operator on a lattice site $`j`$, $`c_j`$, by $`(1)^jc_j^{}`$ (a unitary transformation at half-filling) reverses the sign of $`U`$. When $`U>0`$ this model has a ground state which is an ordered antiferromagnet characterized by the expectation values
$$\left[\begin{array}{c}<S_j^x>\\ <S_j^y>\\ <S_j^z>\end{array}\right]=\frac{1}{2}\left[\begin{array}{c}<c_j^{}c_j+c_j^{}c_j>\\ i<c_j^{}c_jc_j^{}c_j>\\ <c_j^{}c_jc_j^{}c_j>\end{array}\right].$$
(8)
When $`U<0`$ the ground state is thus a degenerate mixture of $`s`$-wave superconductivity and checkerboard charge order characterized by the expectation values
$$\left[\begin{array}{c}<\mathrm{Re}(\mathrm{\Delta }_j)>\\ <\mathrm{Im}(\mathrm{\Delta }_j)>\\ <n_j>\end{array}\right]=\left[\begin{array}{c}<c_j^{}c_j^{}+c_jc_j>\\ i<c_j^{}c_j^{}c_jc_j>\\ (1)^j<c_j^{}c_j+c_j^{}c_j>\end{array}\right]$$
(9)
Both kinds of order occur simultaneously, are equivalent energetically, and may be rotated into each other by analogy with spin rotation of an antiferromagnet. More precisely, this system lies at a quantum phase transition between the two kinds of order and can be made to acquire one, the other, or a mixture of the two by means of an arbitrarily small perturbation, exactly the way the parameters $`\gamma `$ and $`\alpha \alpha ^{}`$ in Eq. (1) break the rotational invariance of $`F`$.
The Hartree-Fock solution, which is only approximate, also has this symmetry. Allowing the expectation values $`y=(1)^j<c_j^{}c_jc_j^{}c_j>/2`$ and $`z=<c_jc_j>`$ we obtain for the Hartree-Fock Hamiltonian
$$_{\mathrm{HF}}=t\underset{<jk>}{}\underset{j\sigma }{}c_{k\sigma }^{}c_{k\sigma }+U\underset{j}{}[(1)^jy(c_j^{}c_j$$
$$c_j^{}c_j)+(zc_j^{}c_j^{}+z^{}c_jc_j)],$$
(10)
and for the corresponding quasiparticle dispersion relation
$$E_𝐤=\pm \sqrt{[(ϵ_𝐤^2+|\mathrm{\Delta }^{\mathrm{CDW}}|^2)^{1/2}\pm \mu ]^2+|\mathrm{\Delta }^{\mathrm{SSC}}|^2},$$
(11)
where $`\mathrm{\Delta }^{\mathrm{CDW}}=yU`$ and $`\mathrm{\Delta }^{\mathrm{SSC}}=zU`$, are the charge density wave and $`s`$-wave superconducting gaps, respectively. This is exactly the same as Eq. (6) with $`s`$-wave quantities substituted for $`d`$-wave ones. Thus, as in the $`d`$-wave case, the superconducting order parameter may, at half-filling ($`\mu =0`$) be rotated continuously from pure superconductivity to pure checkerboard charge order without closing the quasiparticle gap. In this case, however, the rotation also leaves the ground state energy invariant, and is an exact symmetry .
This calculation illustrates the important feature of charge order that it competes easily and naturally with $`s`$-wave superconductivity but not with $`d`$-wave. This is because it is an $`s`$-wave condensate, per Eq. (2).
## VI Pseudogap
A large number of experimental properties of the cuprates are consistent with the presence of DDW order in underdoped samples.
### A Gap Evolution
The $`d`$-wave superconducting gap in the electron spectral function evolves continuously with underdoping into the $`d`$-wave-like pseudogap without collapsing. In the top of Fig. 3 we reproduce point-contact tunneling measurements on underdoped YBCO of Renner et al. showing the excessive size of the tunneling gap and its persistence above the superconducting $`T_c`$, both of which are characteristic of underdoped cuprates. Identification of this feature with the $`d`$-wave gap follows from its evolution out of the simpler BCS-like gap found in overdoped materials and its rough compatibility with the magnitude of $`T_c`$. However, its persistence above $`T_c`$ is not consistent with a traditional BCS gap, for this should disappear at $`T_c`$, as occurs in overdoped samples, on quite general grounds. That this gap has the correct angular dependence is shown in the middle of Fig. 3, where we reproduce angle-resolved photoemission spectra from underdoped Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> at two different points in the Brillouin zone reported by Norman et al. . The upper trace, taken from near the zone face at $`(\pi ,0)`$, shows a large gap that persists well above $`T_c`$, whereas the lower trace, taken from near the node at $`(\pi /2,\pi /2)`$, shows a smaller gap that is destroyed at $`T_c`$. This angular dependence is also seen in the bottom of Fig. 3, where we reproduce the retreat of the photoemission “leading edge” as a function of position on the weak-coupling Fermi surface reported by Harris et al. . The $`d`$-wave-like character of the gap is clear, as is its persistence at the zone face above $`T_c`$ for even slightly underdoped samples. Thus it appears that the pseudogap and the superconducting gap have identical functional forms and evolve continuously into each other as the doping is reduced, just as expected from order-parameter rotation.
Energetic competition as the cause of this rotation is suggested by the similarity between the superconducting and pseudogap energy scales. It may be seen in the bottom of Fig. 3 that the maximum “leading-edge” gapis 30 meV while the retreat caused by heating above $`T_c`$ is
between 5 meV and 10 meV, depending on doping. The pseudogap scale k<sub>B</sub>T$`{}_{}{}^{}`$ 30 meV is also identified in a number of other measurements , notably neutron
scattering , NMR , electronic Raman scattering , and optical reflectivity .
### B Superfluid Density
Rapid collapse of superfluid density below optimal doping is seen in many experiments . The zero-temperature penetration depth, for example, grows rapidly in the pseudogap regime and correlates with the suppression of $`T_c`$ with underdoping, yet saturates at overdoping in a way reminiscent of a traditional BCS superconductor . In Fig. 4 we reproduce the heat capacity measurements on Bi<sub>2.15</sub>Sr<sub>1.85</sub>CaCu<sub>2</sub>O<sub>8+δ</sub> recently reported by Tallon and Loram . Above ahole concentration of about $`p=0.19`$ per Cu the specific heat jump at the superconducting transition varies weakly with p, as one would expect if the material were an ordinary metal undergoing a transition to BCS superconductivity. At $`p=0.19`$, however, there is an abrupt transition and a rapid decrease of this height with underdoping, as though all or part of the Fermi surface were being destroyed by the removal of holes. As a result of this, there are fewer low-energy excitations remaining to
be affected by the superconducting transition. Hence, the specific heat jump at the transition, $`\mathrm{\Delta }\gamma `$, is reduced. All of this behavior is compatible with Fig. 1 if the phase transition is at $`p=0.2`$, where the order parameter $`y`$ begins to develop, is associated with the onset of DDW order and the consequent continuous opening of a gap at $`(\pi ,0)`$ in the quasiparticle spectrum.
### C Spin Susceptibility
There is evidence that spin ordering - and thus presumably stripe ordering \- has not taken place at optimal doping in YBCO, but only occurs at much lower doping levels. In Fig. 5 we reproduce the inelastic neutron measurements for optimally-doped and underdoped
YBCO at a momentum transfer of $`(\pi ,\pi ,\pi )`$. These experiments show that the 41 meV resonance, which disappears above $`T_c`$ and is presumably associated with the superconductivity, evolves continuously with underdoping into the magnetic fluctuation spectrum of the ordered antiferromagnet. Thus, we interpret the piling up of low-frequency spectral weight in the experiment at low doping as signaling the approach of magnetic order, and conversely of showing that magnetic order is neither present nor imminent at the onset of DDW order. The spin-fluctuation spectrum in the superconducting region remains fully gapped and has no low-energy structure of any kind. The resonance continues to be destroyed by elevated temperature, but the requisite temperature grows with underdoping even as $`T_c`$ is evolving to zero. In this way an excitation manifestly associated with the superconductivity at optimal doping transforms into an excitation irrelevant to superconductivity.
This effect is simply understood as a triplet exciton that vanishes at elevated temperature because the quasiparticle gap required for it to be well-defined vanishes. This is quantified in Fig. 6, where we plot the imaginary part of
$$\chi _𝐪(\omega )=\frac{\chi _𝐪^0(\omega )}{1+U\chi _𝐪^0(\omega )},$$
(12)
where
$$\chi _𝐪^0(\omega )=\frac{1}{2\pi ^2}_\pi ^\pi _\pi ^\pi 𝑑k_x𝑑k_y\frac{E_𝐤+E_{𝐤+𝐪}}{(\omega +i\eta )^2(E_k+E_{𝐤+𝐪})^2}$$
$$\times \left(1\frac{\epsilon _𝐤\epsilon _{𝐤+𝐪}+\mathrm{\Delta }_𝐤^{DDW}\mathrm{\Delta }_{𝐤+𝐪}^{DDW}+\mathrm{\Delta }_𝐤^{DSC}\mathrm{\Delta }_{𝐤+𝐪}^{DSC}}{E_𝐤E_{𝐤+𝐪}}\right)$$
(13)
at $`𝐪=(\pi ,\pi )`$ for various values of $`U`$. This is a crude ladder sum in which $`\chi _𝐪^0(\omega )`$ represents the susceptibility of the ideal BCS superconductor characterized by $`E_𝐤`$, $`\mathrm{\Delta }_𝐤^{\mathrm{DDW}}`$, and $`\mathrm{\Delta }_𝐤^{\mathrm{DSC}}`$ per Eq. (6), while $`U`$ represents a coulomb interaction added to push this system toward spin antiferromagnetism. One sees that as $`U`$ is increased the sharp resonance in the spectrum decreases in energy and broadens, just as occurs with decreased doping in Fig. 5. This width is due to efficient decay of the exciton into nodal quasiparticle pairs. At a slightly higher value of $`U`$ the continuum evolves into a divergence at $`\omega =0`$ associated with onset of spin order. Note that the DDW and DSC order parameters in this calculation are effectively interchangeable. Since $`\mathrm{\Delta }_𝐤^{\mathrm{DDW}}=\mathrm{\Delta }_{𝐤+𝐪}^{\mathrm{DDW}}`$ and $`\mathrm{\Delta }_𝐤^{\mathrm{DSC}}=\mathrm{\Delta }_{𝐤+𝐪}^{\mathrm{DSC}}`$ for $`𝐪=(\pi ,\pi )`$, the coherence factor is unity and unchanged close to the Fermi energy whether or not both gaps, or only one of them, are present. For an $`s`$-wave gap the corresponding coherence factor would have been zero.
### D High-Field Transport
Stripes and antiferromagnetic order are naturally associated with the insulating behavior of the cuprates seen near half-filling . In a conventional doped band insulator, insulation is caused by impurities, which trap carriers and prevent them from moving. The system becomes a metal when it is doped sufficiently that the impurity orbitals touch. One of the most significant characteristics of the cuprates is that they continue to insulate to phenomenally high dopings, typically 5% or 1 hole for every 20 Cu atoms. It is very difficult to understand how an insulator with an energy gap less than that of the common semiconductor GaAs should still insulate at these high dopings through impurity trapping solely. But development of antiferromagnetic order with antiphase domain walls, which then trap carriers and pin, is easy to understand, physically sensible, and supported experimentally by the simultaneous occurrence in these materials of discommensurated magnetic Bragg peaks and X-ray satellites at exactly half their momentum displacements . Thus our view is that charge ordering (which would have an order parameter of the form (2), but with an $`f(𝐤)`$ which has $`s`$-wave symmetry and, in all likelihood, incommensurate $`𝐐`$) impedes conduction, rather than facilitating it , and moreover is characteristic of the insulating state.
The issue of coexistence of superconductivity with stripes and antiferromagnetism, and potential causative relations among them, is still highly controversial and a matter of experimental study . There is, however, increasing evidence that the coexistence found in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub>:Nd is anomalous and that the cuprates with the highest values of $`T_c`$ have charge ordering only at the low-doping edge of the superconducting dome. The recent neutron scattering from YBa<sub>2</sub>CuO<sub>3</sub>O<sub>7-x</sub> reported by Mook et al. find the charge-ordering line shown in Fig. 1 and no static antiferromagnetism anywhere in the superconducting region. This is consistent with the the high-field transport experiment on Bi<sub>2</sub>Sr<sub>2-x</sub>La<sub>x</sub>CuO<sub>6+δ</sub> recently reported by Ono
et al. reproduced in Fig. 7, which finds a metal-insulator transition at essentially the same doping as the charge ordering line of when the superconductivity is crushed by a large magnetic field. The phenomenology of this transition is qualitatively similar to that observed previously in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> except that it occurs near the edge of the dome rather than near optimal doping. This is important, for Castellani, Di Castro, and Grilli were led by this observation to propose that the strange-metal behavior of the cuprates might be quantum criticality associated with the charge-ordering transition. These more recent experiments suggest that it is instead quantum criticality associated with the development of DDW order. LSCO is unique among the high-T<sub>c</sub> cuprates in having a low transition temperature, a strong tendency to stripe-order near 1/8 doping, and an extreme sensitivity to Nd doping , all of which suggest mechanical weakness of the crystal structure.
The large-field experiment also reveals another important aspect of the cuprates, namely the lack of evidence for strange-metal behavior in the zero-temperature normal state. It may be seen in the top of Fig. 7 that the resistivities on the metallic side of the transition become constant at low temperatures and that they evolve
continuously across the transition into the linear-$`T`$ resistivity characteristic of the high-temperature normal state of the cuprates. The resistivity at the transition is also about 200 $`\mu \mathrm{\Omega }`$-cm, a typical saturation resistivity in strong-scattering metals. Both of these properties are consistent with the zero-temperature normal state being a conventional metal. They do not prove this, but they make the argument for a non-Fermi-liquid phase more difficult, as linear-$`T`$ resistivity is one of its key signatures.
Thus on the basis of these experiments we predict that in large magnetic fields there should be a second zero-temperature phase transition near $`p=0.19`$ associated with the onset of DDW order. At this transition the system should remain a conventional metal but violently change the topology of its Fermi surface. This transition should be plainly visible in all transport measurements and should be characterized by powerful critical scattering.
### E $`c`$-Axis Conductivity
DDW formation provides a simple explanation for the perplexing semiconducting $`c`$-axis resistivity in many cuprates. In Fig. 8 we reproduce the optical conductivity measurements on YBa<sub>2</sub>Cu<sub>3</sub>O<sub>6.7</sub> of Homes et al. showing the steady reduction of the oscillator strength below 40 meV beginning at a temperature far above the superconducting T<sub>c</sub>. That this reduction does not conserve the f-sum rule locally - which any mean-field theory, including that of the DDW, does - is interesting but not necessarily significant, as a mean-field description is obligated to be quantitative only at arbitrarily small energy scales. Band structure studies of these materials have shown that the $`c`$-axis tunneling matrix element is largest at $`(\pi ,0)`$ and symmetry-related points - precisely at the points where the DDW gap is large. (The functional form is roughly $`t_{}(\mathrm{cos}k_x\mathrm{cos}k_y)^2`$ .) Thus the opening of the DDW gap suppresses the $`c`$-axis transport because the remaining Fermi surface does not conduct efficiently in the c-direction due to small tunneling matrix element. The above matrix element holds for simple tetragonal materials (Hg1201, Tl1201, etc.). For body centered tetragonal materials (LSCO, Tl2201, Bi2212, etc.), the maxima of $`t_{}`$ are shifted towards the zone center, and the effect of opening the DDW gap at $`(\pi ,0)`$ and symmetry-related points is weaker.
## VII Orbital Magnetism
The distinguishing characteristic of DDW order is the magnetic field it makes. Since the possibility of spontaneous breaking of time-reversal and parity in the cuprates was first proposed in the late 1980s there have been a number of attempts to detect such fields, most of which have reported null results . However there has always been confusion about the size of the effects one would expect, and there have always been mysterious magnetic signals in the cuprates, including a recent report of spin antiferromagnetism coexisting with superconductivity in a sample of superoxygenated La<sub>2</sub>CuO<sub>4+y</sub> with y=0.12 and T<sub>c</sub> = 42 K . This fundamentally conflicts with a previous report of no magnetism in La<sub>1.85</sub>Sr<sub>0.15</sub>CuO<sub>4</sub> . We feel that the magentic experiments are so contradictory that they can at present neither rule out nor confirm the presence of DDW order.
We estimate the magnetic field at the center of a plaquette associated with DDW order to be between 1 and 30 gauss . The bond currents of Fig. 2 are roughly $`e\mathrm{\Delta }^{\mathrm{DDW}}/\mathrm{}`$, where $`\mathrm{\Delta }^{\mathrm{DDW}}`$ is the maximum DDW gap. If we take this to be 30 meV, we find bond currents of about $`7\mu A`$. The large uncertainty in the corresponding field strength is due mainly to uncertainty in the current path. One can reasonably consider models ranging from Cu sites connected by 1 Å “wires” to split current carried between adjacent O atoms.
Let us now briefly review the current experimental situation relevant to direct detection of DDW magnetism.
### A Neutron Scattering
DDW order is, in principle, visible in magnetic neutron scattering. Unfortunately the signals are quite small compared with those from ordered spins and easily overwhelmed by them. The ratio of the staggered magnetic field associated with DDW fields to that nominally produced by an ordered array of spins is
$$\frac{B_{\mathrm{DDW}}}{B_{\mathrm{AFM}}}=\left(\frac{e\mathrm{\Delta }^{\mathrm{DDW}}}{\mathrm{}cr}\right)\left(\frac{mcr^3}{e\mathrm{}}\right)=\frac{mr^2}{\mathrm{}^2}\mathrm{\Delta }^{\mathrm{DDW}},$$
(14)
or about 0.06, with $`r=4`$ Å taken for the bond length. Effective magnetic moments of this size are just barely detectable in the cuprates .
It is also unfortunate that the doping levels at which DDW order should be well developed lie close to the spin-glass regime where the system crosses over between Néel and superconducting order. The spin glass is characterized by slightly incommensurated short-range antiferromagnetism with strongly suppressed scattering intensities along one orthorhombic axis - behavior consistent with unpaired Cu spins pointing in the plane and inconsistent with DDW magnetism. However, numerous incursions of this magnetism into the superconducting phase have been reported, in one case deeply , and this has always been difficult to understand from the point of view of traditional magnetic models. It implicitly raises the question of whether there might be two kinds of antiferromagnetism in the cuprates - one, due to spins, which is incompatible with superconductivity and one, due to DDW, which is fully compatible with it and associated with pseudogap formation. Spin-orbit coupling would then mix these and conceivably make them evolve into each other with increased doping.
### B X-ray Scattering
DDW order cannot be seen in X-ray scattering. The DDW order parameter is odd under time-reversal while atomic displacements are even, so there is no first-order coupling between them, and Bragg scattering through circular birefringence from the valence electrons is too weak. For a 10 KeV X-ray of frequency $`\omega `$ the Bragg intensity is down by the factor ($`\mu _BB^{\mathrm{DDW}}/\mathrm{}\omega )^210^{16}`$ from the Bragg intensity of valence electrons - already small compared with the signal from the core electrons. The absence of an X-ray signal is a key characteristic DDW order distinguishing it experimentally from stripes.
### C Magnetic Resonance
The static magnetic field of ideal DDW order cannot be seen directly through NMR of Cu or O nuclei in ideal CuO planes, as these lie at centers of symmetry where the DDW magnetic field is zero. However, magnetic fluctuations associated with the onset of DDW order or a glassy state of a disorder-frustrated DDW could be seen by NMR, although it would be difficult to distinguish from antiferromagnetic spin fluctuations for the reasons stated above. Also, DDW order can in principle be seen in NMR of ions out of the Cu-O planes, such as Y, Ba, La, or Sr. It has long been established that there are unusual magnetic signals below T<sub>c</sub> in all the cuprates, but attempts to quantify these have been plagued by the inherent model-sensitivity of NMR analysis. Tallon and Loram have recently argued using the ratio of <sup>63</sup>Cu and <sup>17</sup>O spin-lattice relaxation rates analyzed with
the model of Millis, Monien, and Pines that short-range antiferromagnetic fluctuations develop in the pseudogap regime with a functional dependence on $`p`$ tracking roughly the value of $`y`$ in Fig. 1. This analysis is not persuasive evidence for DDW order.
### D Muon Spin Resonance
Muon spin resonance has consistently found evidence for magnetism in the superconducting state of the cuprates for dopings less than $`p=0.1`$. In Fig. 9 we reproduce the phase diagram of Niedermeyer et al. showing boundaries of distinct magnetic behaviors observed in powders of both La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> and Y<sub>1-x</sub>Ca<sub>x</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>6.02</sub>. They report a “spin freezing” transition ($`T_f`$) below the Néel transition, and a spin-glass transition ($`T_g`$) that cuts in to the superconducting dome. Below this transition, and at doping levels as high as $`p=0.09`$, the muons depolarize in about 0.1 $`\mu `$S. In the case of LSCO, the measurements extended into the range of the 1/8 anomaly at p = 0.12, beyond which the spin-glass neutron signal tends to disappear and where no magnetism was found in previous $`\mu `$SR measurements . However the fact that both cuprates behave similarly, and that the spin-glass temperature in Y<sub>1-x</sub>Ca<sub>x</sub>Ba<sub>2</sub>Cu<sub>3</sub>O<sub>6</sub> is substantially higher, suggests that this behavior is characteristic of the cuprates as a class. Also, the way the spin-glass line ends has always been confusing.
In a recent paper Panagopoulos et al. have reported anomalous long-time magnetic fluctuations at temperatures just above the glass transition in La<sub>2-x</sub>Sr<sub>x</sub>CuO<sub>4</sub> powders. The scale of these is comparable to T<sub>c</sub> and has a functional dependence with doping identical to that of the parameter $`y`$ in Fig. 1 - i.e. decreasing with doping and vanishing at $`p=0.17`$. Thus they argue that the spin-glass line actually ends here, not at 1/8. The observation of the same effect in a different cuprate, which seems likely in light of Fig. 9, would suggest an intrinsic magnetic signal developing at the onset of DDW order.
## VIII Disorder and Crossover
The muon phenomenology suggests an answer to a question plaguing the idea of competing order in the cuprates, namely why there is no evidence for a genuine phase transition at the pseudogap temperature $`T^{}`$, the alleged phase boundary for onset of DDW, and also why previous searches for magnetism at optimal doping have found sample-dependent or null results. It is simply that the DDW order is corrupted by disorder and transformed into the spin-glass transition at the lower temperature $`T_g`$. In very dirty samples it is lowered so much as to be effectively destroyed.
There has been controversy over how much intrinsic disorder cuprates possess since they were discovered. The essence of the problem is that the most sensitive tests of disorder - transport and the degradation of superconductivity - are corrupted by the non-Fermi-liquid behavior of the normal state evidenced by resistivities which exceed the Ioffe-Regel limit of 100 $`\mu \mathrm{\Omega }`$-cm at T<sub>c</sub> (cf. Fig. 7) and increase with temperature from there. However, using criteria less dependent on the theory of metals, the case for chronic disorder is easier to make: All cuprates lose oxygen easily in arbitrary amounts. All of them have spin-glass phases at low dopings . All of them have magnetic scattering in the superconducting regime that is sample-dependent, difficult to reproduce, and difficult to quantify . All of them have anomalous widths in Cu and O NMR and NQR . Thus our view is that all cuprates made thus far have been significantly disordered, even ones showing evidence to the contrary such as narrow superconducting transition widths. This view is supported by new scanning tunneling microscope experiments on atomically perfect cleaves of optimally-doped BSCCO that find inhomogeneities in the tunneling density of states on the scale of 20 Å.
We note that the disorder need not occur within the CuO<sub>2</sub> planes to have a strong effect on the electronic properties. The non-superconducting cuprates which were studied in the late 1980s differ from the superconducting ones only in the elements that sit between planes. Substituting Hg between the planes raises T<sub>c</sub> substantially. Substituting Nd causes stripes .
The DDW transition is in the same universality class as the random-bond Ising model. The DDW order parameter breaks translational and rotational symmetries, and thus couples to disorder as in a model with a random uniaxial anisotropy. From the Imry-Ma argument as adapted to the random anisotropy case , we know that this symmetry-breaking transition will be spoiled by the random distribution of impurities. Thus, in the presence of disorder, time-reversal is the only true symmetry that can be spontaneously broken by the DDW state. The universality class is then that of a $`Z_2`$ symmetry preserved by the impurities.
The phase diagram of the random-bond Ising model depends critically on the disorder strength. Weak disorder is an irrelevant perturbation and can be ignored at the finite-temperature transition to a state with broken time reversal symmetry. Such a state has a non-vanishing expectation value for the staggered orbital magnetization (thereby breaking the disorder-averaged translational symmetry). If the disorder is strong, on the other hand, and the interlayer coupling is finite, then there can be a finite-temperature transition in the same universality class as the three-dimensional ($`3D`$) Ising spin-glass transition. Due to the weakness of the magnetic coupling between the planes, the spin-glass transition temperature estimated from the two-dimensional spin-glass susceptibility, $`\chi _{sg}^{2D}T^\gamma `$, $`\gamma 5.3`$, is small. Such a transition would not be possible if we could neglect the coupling between the planes, as the lower critical dimension of the Ising spin glass is known to be greater than two.
In contrast to this, the finite temperature transition to DSC remains sharp in the presence of disorder - although $`T_c`$ may be degraded. This is because disorder does not couple linearly to the order parameter as a random field, and because a superconducting transition in two dimensions is possible. The ultimate 3-dimensional transition driven by the coupling between the layers is robust. This can be further understood by invoking the Harris criterion assuming that the transition is in the $`3D`$-XY universality class. The criterion states that weak disorder is an irrelevant perturbation to the pure system if the specific heat exponent is negative, which is indeed the case for the $`3D`$-XY model. When the disorder is so strong that $`k_Fl1`$, where $`k_F`$ is the Fermi wave vector and $`l`$ is the mean free path, a superconductor-insulator transition will take place, and the Harris criterion will no longer apply.
The potential presence in this system of a disorder-sensitive, purely electronic, phase transition involving a Fermi-surface reconnection raises the disturbing possibility that many experiments in this field may be measuring corrupted critical properties of the DDW transition rather than the properties of new states of matter. The notorious non-Fermi-liquid behavior of the normal state, for example, appears to evolve at the lowest temperatures and in a strong magnetic field into behavior of a traditional metal. A possible explanation of this is that the high-temperature behavior is characteristic of a quantum critical region associated with a nearby critical point.
## IX Summary
In summary we find that most of the strange behavior of the cuprate superconductors is consistent evidence for the simultaneous occurrence of $`d`$-wave superconductivity and bond antiferromagnetism. On the basis of this we predict that the spin-glass transition temperature observed in muon spin resonance will climb to the pseudogap temperature $`T^{}`$ as the sample quality improves, that the onset of this effect with doping coincides perfectly with the loss of superfluid density at $`p=0.19`$, and that this transition will be found to be a metal-metal transition involving a Fermi surface reconnection, not a transition to stripe order , when magnetic fields sufficiently intense are available to crush the superconductivity at optimal doping.
## Acknowledgements
RBL wishes to thank Z.-X. Shen, D. Pines, S.-C. Zhang, and G. Aeppli for numerous helpful discussions. RBL, SC, and DM wish to thank the Institute for Complex Adaptive Matter at Los Alamos, where key ideas for this work were conceived. RBL was supported by the National Science Foundation under Grant No. DMR-9813899 and by NEDO. SC was supported by the National Science Foundation under Grant No. DMR-9971138 and, in part, by funds provided by the University of California for the conduct of discretionary research by Los Alamos Natonal Laboratory, under the auspices of the Department of Energy. CN was supported by the National Science Foundation under Grant No. DMR-9983544 and the Alfred P. Sloan Foundation. DM was supported by the Department of Energy at Los Alamos National Laboratory.
## Appendix: Related Density Waves
There are two related unconventional density wave order parameters potentially relevant to the cuprates. The first is that the frustration of the singlet DDW order parameter can lead to incommensurate ordering, in analogy with the Ferrell-Fulde-Larkin-Ovchinnikov state in superconductivity as nesting is destroyed. As in the superconducting case, this will take place for sufficiently strong frustration and at sufficiently low temperatures. Note that in this case the order parameter is allowed to couple with lattice displacements and can therefore be seen in X-ray scattering. When the order parameter is incommensurate, it will no longer have pure $`d_{x^2y^2}`$ symmetry, but will mix in $`p`$-wave terms. For $`𝐐=(\pi /a,\pi /a)+𝐪`$ with $`|𝐪|`$ small, the order parameter will take the form of
equation (2), with $`f(𝐤)=\mathrm{cos}(k_x)\mathrm{cos}(k_y)`$ replaced by
$$f(𝐤)=(1+\frac{i}{2}q_xa)\left[\mathrm{cos}(k_xa)\mathrm{cos}(k_ya)\right]$$
$$\frac{1}{2}q_xa\mathrm{sin}(k_xa)\frac{1}{2}q_ya\mathrm{sin}(k_ya).$$
(15)
The second interesting order parameter is the triplet version of the DDW. This is defined by
$$\stackrel{}{y}=\underset{k}{}f(𝐤)\underset{ss^{}}{}\stackrel{}{\sigma }_{ss^{}}<c_{𝐤+𝐐,s}^{}c_{𝐤,s^{}}>,$$
(16)
where $`\stackrel{}{\sigma }`$ is a Pauli spin matrix. If $`f(𝐤)`$ were chosen to be a function of $`s`$-wave symmetry, this would be a conventional spin-density wave. The order in this case is chracterized by broken time-reversal, translational, and rotational invariances. The combination of any two of time-reversal, a translation by one lattice spacing, or a rotation by $`\pi /2`$ is preserved, however. In addition, spin-rotational symmetry is also broken, which leads to gapless spin-1 excitations. The triplet DDW corresponds to an alternating pattern of spin currents analogous to charge currents of Fig. 2. Presently, the phenomenology of high temperature does not seem to be consistent with the choice of the triplet DDW as the competing order parameter.
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# Quadrupole shape invariants in the interacting boson model
## Abstract
In terms of the Interacting Boson Model, shape invariants for the ground state, formed by quadrupole moments up to sixth order, are studied in the dynamical symmetry limits and over the whole structural range of the IBM-1. The results are related to the effective deformation parameters and their fluctuations in the geometrical model. New signatures that can distinguish vibrator and $`\gamma `$-soft rotor structures, and one that is related to shape coexistence, are identified.
Nuclei are often regarded as drops of nuclear matter as in the geometrical model of Bohr and Mottelson. Having a view of nuclei as such geometrical objects leads directly to the importance of possible deformations of nuclei. The most important deformation of nuclei at low energies is the quadrupole deformation to which we restrict our discussion. These quadrupole deformations are of special interest as they enable us to make predictions of nuclear properties such as energies or $`E2`$ transition strengths of the lowest excited states.
Conversely one can deduce information about nuclear deformations by observing $`E2`$ transition matrix elements. Indeed from a complete set of $`E2`$ matrix elements one can calculate model independent moments and higher order moments of the quadrupole operator, tensorially coupled to a scalar – the shape invariants. Shape invariants were first introduced by Kumar and Cline in the discussion of a large set of $`E2`$ matrix elements obtained in Coulomb excitation experiments. Calculating shape invariants in the geometrical model shows their connection to the deformation parameters $`\beta `$ and $`\gamma `$ used by Bohr and Mottelson or, to be more precise, to effective values $`\beta _{\mathrm{eff}}`$ and $`\gamma _{\mathrm{eff}}`$ and the fluctuations of those. Recently Jolos et al. have introduced approximation formulae to the lowest shape invariants in the framework of the newly developed $`Q`$-phonon scheme . These approximations now make it possible to determine approximate values of the shape invariants from data by using only a few absolute $`B(E2)`$ values.
This is a substantial result since the advent of radioactive beams opens up entirely new nuclear regions for study but, at the same time, the very low intensities of such beams means that data will be sparse and that nuclear structure information must be obtained from fewer and simpler-to-obtain data. Hence the importance of the approximations to the Q-invariants which allow estimates not only of basic deformation parameters such as $`\beta `$ and $`\gamma `$, but of higher moments related to the stiffness of the potential in $`\beta `$ and $`\gamma `$ and to the amount of zero point motion. Such information has seldom if ever been available from any nuclear data. With these approximation formulae, they are now accessible from simple data.
It is therefore important to develop a global view of how these shape invariants behave as a function of structure so that they can be effectively used as signatures of structure. It is the purpose of this Rapid Communication to map out for the first time the behaviour of the five essential invariants, as well as several related quantities, over the full range of nuclear structure. To do so we will use the algebraic Interacting Boson Model (IBM) to study the behaviour of shape invariants in and between the dynamical symmetry limits of the IBM. Formulae will be given to transform the shape invariants into effective deformation parameters $`\beta `$ and $`\gamma `$. The values derived from the algebraic model will be compared to values in the appropriate limiting cases of the geometrical model.
Shape invariants are formed by the isoscalar electric quadrupole operator, which is also the $`E2`$ transition operator in the Consistent Q Formalism (CQF) ,
$$T(E2)=Q=e_BQ^{IBM},$$
(1)
where $`Q^{IBM}`$ is the quadrupole operator in the IBM
$$Q^{IBM}=Q^\chi =s^+\stackrel{~}{d}+d^+s+\chi [d^+\stackrel{~}{d}]^{(2)}$$
(2)
and $`e_B`$ is the effective boson charge which is fixed for a given nucleus. We define moments up to sixth order of the quadrupole operator in the ground state as
$`q_2=`$ $`0_1^+|(QQ)|0_1^+`$ (3)
$`q_3=`$ $`\sqrt{{\displaystyle \frac{35}{2}}}`$ $`|0_1^+|[QQQ]^{(0)}|0_1^+|`$ (4)
$`q_4=`$ $`0_1^+|(QQ)(QQ)|0_1^+`$ (5)
$`q_5=`$ $`\sqrt{{\displaystyle \frac{35}{2}}}`$ $`|0_1^+|(QQ)[QQQ]^{(0)}|0_1^+|`$ (6)
$`q_6=`$ $`{\displaystyle \frac{35}{2}}`$ $`0_1^+|[QQQ]^{(0)}[QQQ]^{(0)}|0_1^+,`$ (7)
where a dot denotes a scalar product and $`[QQQ]^{(0)}`$ abbreviates the tensor coupling $`[Q\times [Q\times Q]^{(2)}]^{(0)}`$. We should note that $`q_2`$ is equal to the total absolute $`E2`$ excitation strength from the ground state
$$q_2=\underset{j}{}B(E2;0_1^+2_j^+).$$
(8)
$`q_2`$ will be the only quantity in our discussion where an absolute value, namely the effective boson charge $`e_B`$, appears.
With the moments (37) we define the relative dimensionless shape invariants by normalizing to an appropriate power of $`q_2`$
$$K_n=\frac{q_n}{q_{2}^{}{}_{}{}^{n/2}}\text{for }n\{3,4,5,6\}.$$
(9)
The quantities $`K_n`$ do not depend on the effective boson charge $`e_B`$. The shape invariants $`K_n`$ differ from earlier definitions of Jolos et al. by normalization constants or tensor coupling. In the present definitions no value may become infinite and all shape invariants are exactly equal to unity in the limit of the rigid symmetric rotor or, in terms of the IBM, the $`SU(3)`$ limit for any boson number $`N`$.
For the calculation of the shape invariants it is convenient to write the expressions for the quadrupole moments (37) as sums over $`E2`$ matrix elements. Therefore the tensor properties of the quadrupole operator are taken into account and the unity operator $`\mathrm{𝟏}=_{J,i,M}|J_iMJ_iM|`$ is inserted between every pair of quadrupole operators. Using the Wigner-Eckert theorem and the unitarity relation of Clebsch Gordan coefficients it is possible to write the moments $`q_n`$ as
$`q_2=`$ $`{\displaystyle \underset{i}{}}0_1^+Q2_i^+2_i^+Q0_1^+`$ (10)
$`q_3=`$ $`\sqrt{{\displaystyle \frac{7}{10}}}`$ $`|{\displaystyle \underset{i,j}{}}0_1^+Q2_i^+2_i^+Q2_j^+2_j^+Q0_1^+|`$ (11)
$`q_4=`$ $`{\displaystyle \underset{i,j,k}{}}0_1^+Q2_i^+2_i^+Q0_j^+0_j^+Q2_k^+`$ (13)
$`2_k^+Q0_1^+`$
$`q_5=`$ $`\sqrt{{\displaystyle \frac{7}{10}}}`$ $`|{\displaystyle \underset{i,j,k,l}{}}0_1^+Q2_i^+2_i^+Q2_j^+2_j^+Q0_k^+`$ (15)
$`0_k^+Q2_l^+2_l^+Q0_1^+|`$
$`q_6=`$ $`{\displaystyle \frac{7}{10}}`$ $`{\displaystyle \underset{i,j,k,l,m}{}}0_1^+Q2_i^+2_i^+Q2_j^+2_j^+Q0_k^+`$ (17)
$`0_k^+Q2_l^+2_l^+Q2_m^+2_m^+Q0_1^+,`$
involving reduced matrix elements between $`0^+`$ and $`2^+`$ states only. In general only the lowest states contribute to the sums because convergence of the $`Q`$-phonon expansion of nuclear states is fast . Matrix elements between nuclear states that differ by several $`Q`$-phonons are usually small .
In the model of a quadrupole deformed rotor analytical expressions for $`E2`$ matrix elements and thus for shape invariants can be obtained. In the rigid rotor the shape invariants are functions of the fixed deformation parameters $`\beta `$ and $`\gamma `$. If we assume a non-rigid deformation, we can give expressions for the shape invariants as
$`q_2`$ $`=`$ $`\left({\displaystyle \frac{3ZeR^2}{4\pi }}\right)^2\beta ^2\left({\displaystyle \frac{3ZeR^2}{4\pi }}\right)^2\beta _{\mathrm{eff}}^{}{}_{}{}^{2}`$ (18)
$`K_3`$ $`=`$ $`{\displaystyle \frac{\beta ^3\mathrm{cos}3\gamma }{\beta ^2^{3/2}}}\mathrm{cos}3\gamma _{\mathrm{eff}}`$ (19)
$`K_4`$ $`=`$ $`{\displaystyle \frac{\beta ^4}{\beta ^2^2}}`$ (20)
$`K_5`$ $`=`$ $`{\displaystyle \frac{\beta ^5\mathrm{cos}3\gamma }{\beta ^2^{5/2}}}`$ (21)
$`K_6`$ $`=`$ $`{\displaystyle \frac{\beta ^6\mathrm{cos}^23\gamma }{\beta ^2^3}},`$ (22)
explicitly using expectation values of $`\beta `$ and $`\mathrm{cos}3\gamma `$. We can define effective values of the deformation parameters $`\beta _{\mathrm{eff}}`$ and $`\gamma _{\mathrm{eff}}`$ by Eqs. (18,19). The parameter $`\gamma _{\mathrm{eff}}=\frac{1}{3}\mathrm{arccos}K_3`$ is given in Table I for the appropriate dynamical symmetry limits of the IBM, where $`SU(3)`$ corresponds to a symmetric rigid rotor, $`O(6)`$ to a $`\gamma `$-soft nucleus with maximal triaxiality and $`U(5)`$ to a vibrator.
The shape invariants are measures of effective deformation parameters and their fluctuations. This is made more explicit by defining the following quantities as measures of the fluctuations of $`\beta `$ and $`\mathrm{cos}3\gamma `$:
$`\sigma _\beta `$ $`=`$ $`{\displaystyle \frac{\beta ^4\beta ^2^2}{\beta ^2^2}}=K_41`$ (23)
$`\sigma _\gamma `$ $`=`$ $`{\displaystyle \frac{\beta ^6\mathrm{cos}^23\gamma \beta ^3\mathrm{cos}3\gamma ^2}{\beta ^2^3}}=K_6K_{3}^{}{}_{}{}^{2}`$ (24)
Using expressions (1017,23, 24) one can analytically calculate the shape invariants in the dynamical symmetry limits of the IBM-1. In this paper we will employ the Extended Consistent Q Formalism (ECQF) of the IBM-1, using the IBM-1 Hamiltonian
$$H_{ECQF}=a\left[(1\zeta )n_d\frac{\zeta }{4N}Q^\chi Q^\chi \right],$$
(25)
with $`Q^\chi `$ taken from Eq. (2). This simple Hamiltonian contains three parameters ($`a`$,$`\zeta `$,$`\chi `$). While one parameter ($`a`$) sets the absolute energy scale, the wave functions depend only on two structural constants ($`\zeta `$,$`\chi `$). For a given nucleus the boson number $`N`$ is fixed. The ECQF-Hamiltonian covers the three dynamical symmetry limits as indicated in Fig. 1. We note that the structural parameter $`\chi `$ appearing in the shape invariants through the $`E2`$ transition operator is fixed in the $`SU(3)`$ limit ($`\chi `$=$`\sqrt{7}/2`$) and in the $`O(6)`$ limit ($`\chi `$=$`0`$) while it is unspecified by the Hamiltonian in the $`U(5)`$ limit. Therefore shape invariants in the $`U(5)`$ limit are functions of the structural parameter $`\chi `$. The analytical expressions for the shape invariants and the fluctuations in the $`U(5)`$ limit as functions of the boson number $`N`$ and the structural parameter $`\chi `$ are
$`K_3^{U(5)}`$ $`=`$ $`\sqrt{{\displaystyle \frac{7}{10}}}{\displaystyle \frac{1}{\sqrt{N}}}|\chi |`$ (26)
$`K_4^{U(5)}`$ $`=`$ $`{\displaystyle \frac{7}{5}}\left(1{\displaystyle \frac{2}{7N}}\right)`$ (27)
$`K_5^{U(5)}`$ $`=`$ $`\sqrt{{\displaystyle \frac{7}{2}}}{\displaystyle \frac{(11N6)}{(5N)^{3/2}}}|\chi |`$ (28)
$`K_6^{U(5)}`$ $`=`$ $`{\displaystyle \frac{21}{25}}{\displaystyle \frac{(N1)}{N^2}}\left(3\chi ^2+N2+{\displaystyle \frac{5N\chi ^2}{6(N1)}}\right)`$ (29)
$`\sigma _\beta ^{U(5)}`$ $`=`$ $`{\displaystyle \frac{2}{5}}\left(1{\displaystyle \frac{1}{N}}\right)`$ (30)
$`\sigma _\gamma ^{U(5)}`$ $`=`$ $`{\displaystyle \frac{21}{25}}{\displaystyle \frac{(N1)}{N^2}}\left(3\chi ^2+N2\right).`$ (31)
For completeness we also give $`K_6`$ as a function of the boson number $`N`$ in the $`O(6)`$ limit
$$K_6^{O(6)}=\frac{1}{3}\frac{(N2)(N1)(N+5)(N+6)}{[N(N+4)]^2}.$$
(32)
These results enlarge the well known symmetry triangle for wave functions of the IBM-1 to a structural ECQF-square for shape invariants and thus for the interpretation of nuclear shapes. This fact is illustrated in Fig. 1. The use of a similar rectangular representation of the parameter space has also been suggested by Bucurescu et al. . Known typical examples for particular points of the ECQF-square are, e.g., <sup>172</sup>Yb for $`SU(3)`$, <sup>196</sup>Pt for $`O(6)`$, <sup>116</sup>Cd for $`U(5)`$ with $`\chi `$=$`0`$ and <sup>152</sup>Sm for large $`\chi `$ and moderate values of $`\zeta `$ (see and discussion below).
From Eqs. (26,28,29,31) we note that the necessary extension of the IBM-1 symmetry triangle to the ECQF-square is a finite-N-effect, because the shape invariants of $`U(5)`$ wave functions converge in the limit $`N\mathrm{}`$ for any value of $`\chi `$. Table I shows the values of the shape invariants and their fluctuations in the dynamical symmetry limits of the IBM-1 for an infinite boson number $`N=\mathrm{}`$. Only the quantities given in Eqs. (2632) depend on the boson number $`N`$.
As we would expect the values of $`\gamma _{\mathrm{eff}}`$ in the symmetry limits are $`0^{}`$ and $`30^{}`$ while the effective triaxiality fluctuates in the $`O(6)`$ and the $`U(5)`$ limits. In the $`SU(3)`$ rigid rotor and $`O(6)`$ $`\gamma `$-soft limits the $`\beta `$ deformation is rigid while it fluctuates in the $`U(5)`$ vibrator limit.
We have discussed the shape invariants and fluctuations in the dynamical symmetry limits of the IBM using analytical expressions. These values provide useful benchmarks for the geometrical interpretation of IBM ground state wave functions. However, the dynamical symmetry limits of the IBM and the corresponding geometrical models are idealised, analytically solvable limits. More accurate descriptions of the low energy structure of collective nuclei can usually be obtained by IBM Hamiltonians outside the dynamical symmetry limits. To gain insight in the structure of actual nuclei the quantities of interest have been calculated between the symmetry limits, using the ECQF-Hamiltonian (25). The shape invariants and their fluctuations have been calculated gridwise over the whole IBM parameter space for $`N=10`$ bosons as functions of the structural parameters $`\zeta `$ and $`\chi `$. All calculations have been performed by diagonalizing the Hamiltonian numerically using the computer code PHINT . Calculations of the shape invariants have been done by a FORTRAN code (QINVAR) which evaluates the PHINT output.
All quantities behave smoothly and one obtains an impression of how the quantities vary outside of the dynamical symmetry limits. Fig. 2 represents the numerical results of this work presenting the variation of the most important quadrupole invariants over all ranges of structure. The behaviour of the invariants, $`q_2`$, $`K_3`$-$`K_6`$, between the symmetries is interesting. Strong variations towards and for deformed nuclei are typical. The invariant $`K_4`$, which is related to fluctuations in $`\beta `$ via Eq. (23), is one of the few observables that can distinguish $`U(5)`$ from $`O(6)`$. This can be useful in newly accessible exotic nuclei since $`K_4`$ can be approximately obtained from the simple expression
$$K_4\frac{7}{10}\frac{B(E2;4_1^+2_1^+)}{B(E2;2_1^+0_1^+)}K_4^{\mathrm{appr}.},$$
(33)
which involves two observables, easily measured, e.g., by Coulomb excitation experiments. The approximation (33) is valid within about $`10\%`$ for the ECQF-square and boson numbers $`N5`$. This was numerically checked for the whole ECQF-square and for boson numbers $`N=5,7,10,16`$. For a detailed analysis an experimental value of $`K_4^{\mathrm{appr}.}`$ can serve as a benchmark for starting points of numerical IBM calculations, which can be optimized to reproduce the measured transition strengths. The actual value of $`K_4`$ can then be determined from the complete set of calculated $`E2`$ transition matrix elements.
For large $`N`$, $`K_6`$ and $`\sigma _\gamma `$ are quite different in $`U(5)`$ and $`O(6)`$ which is evident from Table I. The bottom right panel of Fig. 2 shows $`\sigma _\gamma `$, which gives the fluctuations in $`\gamma `$, gridwise over the full structural range. Note, however, that Fig. 2 is calculated for a finite boson number which lowers the value of $`K_6`$ and $`\sigma _\gamma `$ in the $`U(5)`$ and $`O(6)`$ limits as seen from Eqs. (29,31,32). In the $`SU(3)`$ limit $`\sigma _\gamma `$ vanishes, which characterizes the $`SU(3)`$ limit as a model for a rigid rotor, also in the $`\gamma `$ degree of freedom. In contrast non-vanishing triaxiality fluctuations occur in the $`U(5)`$ and $`O(6)`$ limits, indicating that these limits and the whole transitional region between them model $`\gamma `$-soft nuclei.
Finally, we note that all shape invariants, especially $`\sigma _\gamma `$, change strongly between $`SU(3)`$ and $`U(5)`$-like values in an unusual region of the IBM-1 parameter space, namely for moderate values of $`\zeta `$ and $`\chi `$=$`\sqrt{7}/2`$. Interestingly, this is just the region appropriate to the nucleus <sup>152</sup>Sm ($`\zeta `$=$`0.57`$,$`\chi `$=$`\sqrt{7}/2`$) . The case of <sup>152</sup>Sm is currently under active discussion and it seems that it shows a certain degree of shape coexistence between spherical and deformed shapes with large effective triaxiality .
Above, we discussed the numerical calculation of the exact shape invariants within the $`sd`$-IBM-1 parameter space, using the ECQF-Hamiltonian (25). One aspect of this work is to establish the shape invariants as a convenient link between the geometrical model and any other nuclear structure model which is able to calculate $`E2`$ transition matrix elements. Here we have chosen the algebraical IBM. Our ansatz is alternative to the intrinsic state formalism by Ginocchio and Kirson which was used much earlier to link the IBM Hamiltonian to the geometrical Bohr Hamiltonian.
We note that the effective values of the shape parameters $`\beta _{\mathrm{eff}}`$ and $`\gamma _{\mathrm{eff}}`$ do in general not exactly coincide with the minima of a corresponding energy surface for the ground state in the deformation parameter plane. However, the shape invariants can easily be used to compare the predictions from different nuclear models in a geometrically transparent way.
In principle, the shape invariants can also be measured directly from extensive nuclear structure data, providing a direct test of nuclear structure models. Much more intriguing is the common case when only a few key observables, like $`E2`$ branching ratios from low-lying $`0^+`$ states and $`2^+`$ states, are known experimentally and when a phenomenological nuclear structure model, like the IBM, can be used to extrapolate the data to a complete set of $`E2`$ transition matrix elements.
To summarize, we have presented analytic expressions for moments up to sixth order of the quadrupole operator in the ground state and we have given definitions for the lowest shape invariants up to $`K_6`$. The shape invariants were calculated analytically in the dynamical symmetry limits of the IBM-1. Formulae were given to derive effective deformation parameters and their fluctuations from shape invariants, and thus from IBM-1 calculations. A study, using the ECQF-Hamiltonian (25), of the behaviour of the shape invariants over a full range of structures has been performed for the first time. It shows the smooth but yet widely varying behaviour of the invariants. Thus they can be used to determine the properties of nuclei by comparing the calculated invariants to experimentally obtained values or to results of fits. Moreover, approximate values of these invariants can be obtained experimentally simply from $`B(E2)`$ values involving just the $`2_1^+`$, $`2_2^+`$ and $`4_1^+`$ states, and, for $`K_5`$ and $`K_6`$, a $`B(E2)`$ branching ratio from the appropriate excited $`0^+`$ states.
The invariant $`K_4`$, as well as the fluctuation $`\sigma _\gamma =K_6K_3^2`$, are of special interest as they allow to distinguish between $`O(6)`$ and $`U(5)`$ symmetries which can be difficult otherwise . Finally, the values of $`\sigma _\gamma `$ change most rapidly for IBM-1 Hamiltonians that show shape coexistence.
For fruitful discussions the authors thank A. Gelberg, T. Otsuka and N.V. Zamfir. This work has been partly supported by the Deutsche Forschungsgemeinschaft under Contract Nos. Br 799/9-1 and Pi 393/1-1, and by the U.S. DOE under Grant No. DE-FG02-91ER40609.
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# Gaussian superpositions in scalar-tensor quantum cosmological models
## 1 Introduction
The existence of an initial singularity is one of the major drawbacks of classical cosmology. In spite of the fact that the standard cosmological model, based in the classical general relativity theory, has been successfully tested until the nucleosynthesis era (around $`t1s`$), the extrapolation of this model to higher energies leads to a breakdown of the geometry in a finite cosmic time. This breakdown of the geometry may indicate that the classical theory must be replaced by a quantum theory of gravitation: quantum effects may avoid the presence of the singularity, leading to a complete regular cosmological model.
The quantization of gravity is plagued with many conceptual and technical problems, and when it is applied to the whole universe new issues appear. In the Dirac quantization approach, a functional equation for the wave function of the Universe is obtained, the Wheeler-DeWitt equation , which is the basic equation of quantum cosmology. It is formulated in the so-called superspace, the space of all possible three-dimensional spatial geometries. It is very hard to find exact solutions of the full Wheeler-DeWitt equation, but solutions may be found in minisuperspaces where all but a finite number of degrees of freedom are frozen.
Among the fundamental questions that come from the quantization of the universe as a whole, one of the most important concerns the interpretation of the wave function coming from the Wheeler-DeWitt equation. In order to extract predictions from the wave function of the Universe, the Bohm-de Broglie ontological interpretation of quantum mechanics has been proposed , since it avoids many conceptual difficulties that follow from the application of the standard Copenhagen interpretation to an unique system that contains everything. In opposition to the latter one, the ontological interpretation does not need a classical domain outside the quantized system to generate the physical facts out of potentialities (the facts are there ab initio), and hence it can be applied to the universe as a whole<sup>1</sup><sup>1</sup>1Other alternative interpretations can be used in quantum cosmology like the many worlds interpretation of quantum mechanics .With this interpretation in hands, one can ask if the quantum scenario predicted by the Wheeler-DeWitt equation is free of singularities and which type of classical universe emerges from the quantum phase.
In a preceding work , we have applied this proposal to a free scalar-tensor model with minimal coupling in Friedmann-Robertson-Walker geometry, which can be obtained from a non-minimal scalar-tensor theory through a conformal transformation. Free scalar fields are good candidates to describe the material content of the early Universe because of their simplicity and because they represent stiff matter, the type of fluid advocated by Zel’dovich to be relevant at early stages of cosmic evolution. Only positive curvature spatial sections have been studied. The bohmian trajectories in configuration space revealed an unexpected scenario: they behaved as the classical solutions for small values of the scale factor, but display quantum behaviour when the scale factor is big. As a consequence, the initial singularity is still present in this quantum model.
The Wheeler-DeWitt solutions for this scalar-tensor model contain positive and negative frequency modes, the first leading to an expanding universe, and the second to a contracting universe, near the singularity. Inspired by this observation, we constructed in some particular superpositions mixing negative and positive models. In this way, we found non-singular quantum solutions which were, however, of planckian size and hence they could not be a model for our real Universe.
The aim of the present work is to explore further the possibilities of the minisuperspace model of Reference . First, we will not restrict ourselves to positive curvature spatial sections and second, we will explore more suitable superpositions of negative and positive modes, namely, the gaussian superposition. For the case the spatial section is flat, it is possible to solve analytically the expressions for the phase of the wave function, and to reduce the equations for the bohmian trajectories to a dynamical system. The critical points are calculated, and they are identified as center or nodes points. This leads to the existence of three kind of scenarios: periodic solutions representing oscillating universes; bouncing universes; models with a big-bang followed by a big-crunch. The bouncing universes contract classically from infinity until a minimum size, where quantum effects become important acting as repulsive forces avoiding the singularity, expanding afterwards to an infinite size, approaching the classical expansion as long as the scale factor increases. These are non-singular solutions which are viable models to describe the Universe we live in. For closed and open spatial sections, all calculations must be performed numerically, and the trajectories obtained in the configuration space reveal again the presence of oscillating universes besides those with a big-bang followed by a big-crunch. In all three cases, the oscillating universes have a characteristic scale of the order of the Planck length, except for very special gaussians in the case of zero saptial curvature. Hence, the most interesting scenarios emerge from the flat case, where we have succeeded to obtain a viable non-singular model.
The article is organized as follows. In section 2, we describe the classical model and the corresponding Wheeler-DeWitt equation in the minisuperspace. Section 3 is devoted to the study of the gaussian superposition of the quantum solutions found before, and their corresponding analysis. In section 4 we present our conclusions.
## 2 The classical and quantum minisuperspace models
Let us take the lagrangian
$$L=\sqrt{g}e^\varphi (Rw\varphi _{;\rho }\varphi ^{;\rho }).$$
(1)
For $`w=1`$ we have effective string theory without the Kalb-Ramond field. For $`w=3/2`$ we have a conformally coupled scalar field. Performing the conformal transformation $`g_{\mu \nu }=e^\varphi \overline{g}_{\mu \nu }`$ we obtain the following lagrangian:
$$L=\sqrt{g}[R(\omega +\frac{3}{2})\varphi _{;\rho }\varphi ^{;\rho }],$$
(2)
where the bars have been omitted. We will define $`C_w(\omega +\frac{3}{2})`$, which we will consider, from now on, to be strictly positive in order not to violate any of the energy conditions, at least classically.
We will consider the Robertson-Walker metric
$$ds^2=N^2\mathrm{d}t^2+\frac{a(t)^2}{1+\frac{ϵ}{4}r^2}[\mathrm{d}r^2+r^2(\mathrm{d}\theta ^2+\mathrm{sin}^2(\theta )\mathrm{d}\phi ^2)],$$
(3)
where the spatial curvature $`ϵ`$ takes the values $`0`$, $`1`$,$`1`$. Inserting this line element into the lagrangian (2), and using the units where $`\mathrm{}=c=1`$, we obtain the following action:
$$S=\frac{3V}{4\pi l_p^2}\frac{Na^3}{2}(\frac{\dot{a}^2}{N^2a^2}+C_w\frac{\dot{\varphi }^2}{6N^2}+\frac{ϵ}{a^2})\mathrm{d}t,$$
(4)
where $`V`$ is the total volume divided by $`a^3`$ of the spacelike hypersurfaces, which are supposed to be closed, and $`l_p`$ is the Planck length. $`V`$ depends on the value of $`ϵ`$ and on the topology of the hypersurfaces. For $`ϵ=0`$, $`V`$ can have any value because the fundamental polyhedra of $`ϵ=0`$ hypersurfaces can have arbitrary size (see Ref. ). In the case of $`ϵ=1`$ and topology $`S^3`$, $`V=2\pi ^2`$. Defining $`\beta ^2=\frac{4\pi l_p^2}{3V}`$, $`\overline{\varphi }\sqrt{\frac{C_w}{6}}\varphi `$, and omitting again the bars, the hamiltonian reads
$$H=N(\beta ^2\frac{p_a^2}{2a}+\beta ^2\frac{p_\varphi ^2}{2a^3}ϵ\frac{a}{2\beta ^2}).$$
(5)
where
$`p_a`$ $`=`$ $`{\displaystyle \frac{a\dot{a}}{\beta ^2N}},`$ (6)
$`p_\varphi `$ $`=`$ $`{\displaystyle \frac{a^3\dot{\varphi }}{\beta ^2N}}.`$ (7)
Usually, the scale factor has dimensions of length because we use angular coordinates in closed spaces. Hence we will define a dimensionless scale factor $`\stackrel{~}{a}a/\beta `$. In that case the hamiltonian becomes, omitting the tilde:
$$H=\frac{N}{2\beta }(\frac{p_a^2}{a}+\frac{p_\varphi ^2}{a^3}ϵa).$$
(8)
As $`\beta `$ appears as an overall multiplicative constant in the hamiltonian, we can set it equal to one without any loss of generality, keeping in mind that the scale factor which appears in the metric is $`\beta a`$, not $`a`$. We can further simplify the hamiltonian by defining $`\alpha \mathrm{ln}(a)`$ obtaining
$$H=\frac{N}{2\mathrm{exp}(3\alpha )}[p_\alpha ^2+p_\varphi ^2ϵ\mathrm{exp}(4\alpha )],$$
(9)
where
$`p_\alpha `$ $`=`$ $`{\displaystyle \frac{e^{3\alpha }\dot{\alpha }}{N}},`$ (10)
$`p_\varphi `$ $`=`$ $`{\displaystyle \frac{e^{3\alpha }\dot{\varphi }}{N}}.`$ (11)
The momentum $`p_\varphi `$ is a constant of motion which we will call $`\overline{k}`$. The classical solutions are, in the gauge $`N=1`$:
1) $`ϵ=0`$
$$\varphi =\pm \alpha +c_1,$$
(12)
where $`c_1`$ is an integration constant. In term of cosmic time they read:
$`a`$ $`=`$ $`e^\alpha =3\overline{k}t^{1/3},`$ (13)
$`\varphi `$ $`=`$ $`{\displaystyle \frac{\mathrm{ln}(t)}{3}}+c_2.`$ (14)
The solutions contract or expand forever from a singularity, depending on the sign of $`\overline{k}`$, without any inflationary epoch.
2) $`ϵ=1`$
$$a=e^\alpha =\frac{\overline{k}}{\mathrm{cosh}(2\varphi c_1)},$$
(15)
where $`c_1`$ is an integration constant, and from the conservation of $`p_\varphi `$ we get
$$\overline{k}=e^{3\alpha }\dot{\varphi }.$$
(16)
The cosmic time dependence is complicated and we will not write it here. These solutions describe universes expanding from a singularity till a maximum size and contracting again to a big crunch. Near the singularity, these solutions behave as in the flat case. There is no inflation.
3) $`ϵ=1`$
$$a=e^\alpha =\frac{\overline{k}}{\mathrm{sinh}(2\varphi c_1)},$$
(17)
where $`c_1`$ is an integration constant, and again, from the conservation of $`p_\varphi `$ we get
$$\overline{k}=e^{3\alpha }\dot{\varphi }.$$
(18)
As before, the cosmic time dependence is complicated and we will not write it here. These solutions describe universes contracting forever to or expanding forever from a singularity. Near the singularity, these solutions behave as in the flat case. There is no inflation<sup>2</sup><sup>2</sup>2In the case $`ϵ=1`$ there are classical solutions with $`C_w<0`$. Qualitatively, they represent universes contracting from an infinite to a minimum size and then expanding again to infinity.
Let us now quantize the model. The Wheeler-DeWitt equation is obtained through the Dirac quantization procedure where the wave function must be annihilated by the operator version of the constraint in Eq. (9). With the choice of factor ordering which makes it covariant through field redefinitions, it reads
$$\frac{^2\mathrm{\Psi }}{\alpha ^2}+\frac{^2\mathrm{\Psi }}{\varphi ^2}+ϵe^{4\alpha }\mathrm{\Psi }=0.$$
(19)
Employing the separation of variables method, we obtain the general solution
$$\mathrm{\Psi }(\alpha ,\varphi )=F(k)A_k(\alpha )B_k(\varphi )𝑑k,$$
(20)
where $`k`$ is a separation constant,
$$B_k(\varphi )=b_1\mathrm{exp}(ik\varphi )+b_2\mathrm{exp}(ik\varphi ),$$
(21)
and for $`ϵ=0`$
$$A_k(\alpha )=a_1\mathrm{exp}(ik\alpha )+a_2\mathrm{exp}(ik\alpha ),$$
(22)
for $`ϵ=1`$
$$A_k(\alpha )=a_1I_{ik/2}(e^{2\alpha }/2)+a_2K_{ik/2}(e^{2\alpha }/2),$$
(23)
and for $`ϵ=1`$
$$A_k(\alpha )=a_1J_{ik/2}(e^{2\alpha }/2)+a_2N_{ik/2}(e^{2\alpha }/2).$$
(24)
The functions $`J,N,I,K`$ are Bessel and modified Bessel functions of first and second kind.
The Bohm-de Broglie interpretation of homogeneous minisuperspace models goes as follows: in general, the minisuperspace Wheeler-De Witt equation is
$$(\widehat{p}^\alpha (t),\widehat{q}_\alpha (t))\mathrm{\Psi }(q)=0.$$
(25)
Writing $`\mathrm{\Psi }=R\mathrm{exp}(iS/\mathrm{})`$, and substituting it into (25), we obtain the following equation:
$$\frac{1}{2}f_{\alpha \beta }(q_\mu )\frac{S}{q_\alpha }\frac{S}{q_\beta }+U(q_\mu )+Q(q_\mu )=0,$$
(26)
where the quantum potential is
$$Q(q_\mu )=\frac{1}{2R}f_{\alpha \beta }\frac{^2R}{q_\alpha q_\beta }.$$
(27)
The Bohm-de Broglie interpretation applied to quantum cosmology states that the trajectories $`q_\alpha (t)`$ are real, independently of any observations. Eq. (26) is the Hamilton-Jacobi equation for them, which is the classical one amended with a quantum potential term (27), responsible for the quantum effects. This suggests to define
$$p^\alpha =\frac{S}{q_\alpha },$$
(28)
where the momenta are related to the velocities in the usual way:
$$p^\alpha =f^{\alpha \beta }\frac{1}{N}\frac{q_\beta }{t}.$$
(29)
To obtain the quantum trajectories we have to solve the following system of first order differential equations, called the guidance relations:
$$\frac{S(q_\alpha )}{q_\alpha }=f^{\alpha \beta }\frac{1}{N}\dot{q}_\beta .$$
(30)
In the present case of the hamiltonian (9), the quantum potential (27) becomes
$$Q(\alpha ,\varphi )=\frac{e^{3\alpha }}{2R}[\frac{^2R}{\alpha ^2}\frac{^2R}{\varphi ^2}],$$
(31)
and the guidance relations (30) read
$$\frac{S}{\alpha }=\frac{e^{3\alpha }\dot{\alpha }}{N},$$
(32)
$$\frac{S}{\varphi }=\frac{e^{3\alpha }\dot{\varphi }}{N}.$$
(33)
Eqs. (30) are invariant under time reparametrization. Hence, even at the quantum level, different choices of $`N(t)`$ yield the same spacetime geometry for a given non-classical solution $`q_\alpha (t)`$. There is no problem of time in the Bohm-de Broglie interpretation of minisuperspace quantum cosmology<sup>3</sup><sup>3</sup>3This is not the case, however, for the full superspace (see Reference ).. Let us then apply this interpretation to our minisuperspace models and choose the gauge $`N=1`$.
## 3 Bohm interpretation of gaussian superpositions
We will now make gaussian superpositions of these solutions and interpret the results using the Bohm-de Broglie interpretation of quantum mechanics. We will begin by the case $`ϵ=0`$, which is simpler, and it is the one to which the others reduce when $`\alpha \mathrm{}`$.
### 3.1 Hypersurfaces with $`ϵ=0`$
This case can be solved analytically. The function $`F(k)`$ is
$$F(k)=\mathrm{exp}[\frac{(kd)^2}{\sigma ^2}].$$
(34)
We can study two types of wave function:
$$\mathrm{\Psi }_1(\alpha ,\varphi )=F(k)B_k(\varphi )[A_k(\alpha )+A_k(\alpha )]𝑑k,$$
(35)
and
$$\mathrm{\Psi }_2(\alpha ,\varphi )=F(k)A_k(\alpha )[B_k(\varphi )+B_k(\varphi )]𝑑k,$$
(36)
both with $`a_2=b_2=0`$. We will restrict ourselves to $`\mathrm{\Psi }_1`$ because it yields the most interesting results. The results coming from $`\mathrm{\Psi }_2`$ can be obtained from the first by changing $`\alpha `$ with $`\varphi `$.
Performing the integration in $`k`$ we obtain for $`\mathrm{\Psi }_1`$:
$$\mathrm{\Psi }_1=\sigma \sqrt{\pi }\{\mathrm{exp}[\frac{(\alpha +\varphi )^2\sigma ^2}{4}]\mathrm{exp}[id(\alpha +\varphi )]+\mathrm{exp}[\frac{(\alpha \varphi )^2\sigma ^2}{4}]\mathrm{exp}[id(\alpha \varphi )]\}.$$
(37)
In order to obtain the bohmian trajectories, we have to calculate the phase $`S`$ of the above wave function and substitute it into the guidance formula (3233), working in the gauge $`N=1`$. These equations constitute a planar system which can be easily studied:
$`\dot{\alpha }`$ $`=`$ $`{\displaystyle \frac{\left[\varphi \sigma ^2\mathrm{sin}(2d\alpha )+2d\mathrm{sinh}(\sigma ^2\alpha \varphi )\right]}{\mathrm{exp}(3\alpha )\left\{2[\mathrm{cos}(2d\alpha )+\mathrm{cosh}(\sigma ^2\alpha \varphi )]\right\}}},`$ (38)
$`\dot{\varphi }`$ $`=`$ $`{\displaystyle \frac{\left[\alpha \sigma ^2\mathrm{sin}(2d\alpha )+2d\mathrm{cos}(2d\alpha )+2d\mathrm{cosh}(\sigma ^2\alpha \varphi )\right]}{\mathrm{exp}(3\alpha )\left\{2[\mathrm{cos}(2d\alpha )+\mathrm{cosh}(\sigma ^2\alpha \varphi )]\right\}}}.`$ (39)
The line $`\alpha =0`$ divides configuration space in two symmetric regions. The line $`\varphi =0`$ contains all singular points of this system, which are nodes and centers. The nodes appear when the denominator of the above equations, which is proportional to the norm of the wave function, is zero. No trajectory can pass through these points. They happen when $`\varphi =0`$ and $`\mathrm{cos}(d\alpha )=0`$, or $`\alpha =(2n+1)\pi /2d`$, $`n`$ an integer, with separation $`\pi /d`$. The center points appear when the numerators are zero. They are given by $`\varphi =0`$ and $`\alpha =2d[\mathrm{cotan}(d\alpha )]/\sigma ^2`$. They are intercalated with the node points, and their separations cannot exceed $`\pi /d`$. As $`\alpha \mathrm{}`$ these points tend to $`n\pi /d`$. As one can see from the above system, the classical solutions ($`a(t)t^{1/3}`$) are recovered when $`\alpha \mathrm{}`$ or $`\varphi \mathrm{}`$, the other being different from zero.
A field plot of this planar system is shown in Figure 1, for $`\sigma =d=1`$. We can see plenty of different possibilities, depending on the initial conditions. Near the center points we can have oscillating universes without singularities and with amplitude of oscillation of order 1<sup>4</sup><sup>4</sup>4As discussed above, these amplitudes can be very large as long as $`d`$ becomes very small because the separation of the center points are of the order of $`1/d`$. For negative values of $`\alpha `$, the universe arises classically from a singularity but quantum effects become important forcing it to recollapse to another singularity, recovering classical behaviour near it. For positive values of $`\alpha `$, the universe contracts classically but when $`\varphi `$ and $`\alpha `$ are small enough, quantum effects become important creating an inflationary phase which avoids the singularity. The universe contracts to a minimum size and after reaching this point it expands forever, recovering the classical limit when $`\alpha `$ becomes sufficiently large. These are models which can represent the early Universe. We can see that for $`\alpha `$ negative we have classical limit for small scale factor while for $`\alpha `$ positive we have classical limit for big scale factor.
For the wave function $`\mathrm{\Psi }_2`$, the analysis goes in the same way but we have to interchange $`\alpha `$ with $`\varphi `$. In this case we also have periodic solutions but the others are universes arising classically from a singularity, experiencing quantum effects in the middle of their expansion, and recovering their classical behaviour for large values of $`\alpha `$. There are no further possibilities.
We will now pass to the cases with curved spatial sections. One can immediately notice an important difference. The case $`ϵ=0`$ has a symmetry $`\alpha \alpha `$ which is present not only in the Wheeler-DeWitt equation (19) but also in the solution (37). The cases $`ϵ0`$ do not possess this symmetry (the potential $`ϵe^{4\alpha }`$ in the Wheeler-DeWitt equation breaks it), and one should not expect to obtain the $`\alpha >0`$ part of the field plots in these cases from the $`\alpha <0`$ part through a reflection, as in the case $`ϵ=0`$ (see Figure 1).
### 3.2 Hypersurfaces with $`ϵ=1`$
The Wheeler-DeWitt equation (19) for $`ϵ=1`$, in the case we neglect the $`_{\varphi \varphi }\mathrm{\Psi }`$, is analogous to a stationary Schroedinger equation with $`E=0`$ and $`V=e^{4\alpha }`$. Hence, one should make superpositions involving only the parts of $`A_k(\alpha )`$ which goes to zero as $`\alpha `$ goes to infinity, which are the Bessel functions $`K_{ik/2}(e^{2\alpha /2})`$. Consequently, we will take the following superposition:
$$\mathrm{\Psi }_3(\alpha ,\varphi )=\mathrm{exp}[\frac{(kd)^2}{\sigma ^2}]K_{ik/2}(e^{2\alpha }/2)B_k(\varphi )dk.$$
(40)
The limit $`\alpha \mathrm{}`$ does not give the preceding results for $`ϵ=0`$ because the Bessel function $`K`$ reduces in this limit to
$$K_{ik/2}(e^{2\alpha }/2)\frac{i}{k}[\mathrm{exp}[ik(\alpha \mathrm{ln}(2)]\mathrm{\Gamma }(1\frac{ik}{2})\mathrm{exp}[ik(\alpha \mathrm{ln}(2)]\mathrm{\Gamma }(1+\frac{ik}{2})],$$
(41)
and the presence of the Gamma functions spoils their similarity.
This case must be studied numerically, and the transformation
$`\alpha ^{}`$ $`=`$ $`\alpha \varphi ,`$ (42)
$`\varphi ^{}`$ $`=`$ $`\varphi +\alpha ,`$ (43)
eases this task. The guidance relations (3233) become
$`\dot{\alpha }^{}`$ $`=`$ $`2\mathrm{exp}\left[{\displaystyle \frac{3(\alpha ^{}+\varphi ^{})}{2}}\right]{\displaystyle \frac{S}{\varphi ^{}}},`$ (44)
$`\dot{\varphi }^{}`$ $`=`$ $`2\mathrm{exp}\left[{\displaystyle \frac{3(\alpha ^{}+\varphi ^{})}{2}}\right]{\displaystyle \frac{S}{\alpha ^{}}},`$ (45)
and Figure 2 shows the field plot of this transformed planar system, using $`\sigma =d=1`$. There are periodic solutions without singularities which happen when the bohmian trajectories cross the lines $`(\alpha ^{}=3.73`$, $`\varphi ^{}<3.73)`$, or $`(\varphi ^{}=3.73`$, $`\alpha ^{}<3.73)`$, or equivalently $`(\alpha =\varphi 3.73`$, $`\alpha <3.73)`$. These oscillating trajectories can reach very negative values of $`\alpha `$ but their maximum size cannot exceed $`\alpha =0`$, or $`al_{pl}`$. Another behaviour is related to the trajectories shown in Figure 2 which are exclusively in the light gray region. They begin classically from a singularity, expand to a maximum value of $`\alpha `$, and then return classically to a singularity. Concluding, we have two types of trajectories in this case: one which is periodic due to quantum effects, and the other which exhibit the pattern of classical behaviour: expansion from a singularity until a maximum size followed by a contraction to a big crunch. The periodic solutions have maximum size around $`\alpha =0`$, or $`al_{pl}`$ and they cannot represent the Universe we live in.
### 3.3 Hypersurfaces with $`ϵ=1`$
In this case we will choose as $`A_k(\alpha )`$ the combination
$$A_k(\alpha )=[\mathrm{\Gamma }(1+\frac{ik}{2})J_{ik/2}(e^{2\alpha }/2)+\mathrm{\Gamma }(1\frac{ik}{2})J_{ik/2}(e^{2\alpha }/2)]$$
(46)
in order to get rid of the Gamma functions and obtain the preceding results for $`ϵ=0`$ when $`\alpha `$ is very negative because the Bessel function $`J`$ reduces in this limit to
$$J_{ik/2}(e^{2\alpha }/2)\frac{\mathrm{exp}[ik(\alpha \mathrm{ln}(2)]}{\mathrm{\Gamma }(1+\frac{ik}{2})}.$$
(47)
Taking this choice of $`A_k(\alpha )`$, Eq. (46), into the gaussian superposition
$$\mathrm{\Psi }_4(\alpha ,\varphi )=\mathrm{exp}[\frac{(kd)^2}{\sigma ^2}]A_k(\alpha )B_k(\varphi )dk,$$
(48)
the numerical calculations with respect to the $`ϵ=1`$ case show that the behaviour for very negative values of $`\alpha `$ is similar to the $`ϵ=0`$ case, as one can see by comparing Figure 3 with Figure 1. As $`\alpha `$ increases the regions with oscillating universes are squeezed and their separation decrease monotonically. Like the $`ϵ=0`$ case, there are periodic solutions without singularities and with amplitude of oscillation of order $`1`$. The other behaviour is described by trajectories that arise classically from a singularity, experiment a quantum halt at some maximum value of the scale factor, and then classically contracts to a big-crunch, contrary to the classical solutions of Eq. (17) which contract forever to or expand forever from a singularity.
## 4 Conclusion
The quantization of a scalar-tensor model in the minisuperspace leads to a separable partial differential equation, admitting analytical solutions, with positive and negative frequencies. In this work, we have studied gaussian superpositions of these different modes and the corresponding bohmian trajectories. Such analysis was performed for zero, positive and negative curvature spatial sections, which are considered to be compact. The bohmian trajectories in configuration space were calculated numerically, excepted for the flat case, where it is possible to reduce the equations for the bohmian trajectories to a two dimensional dynamical system.
The comparison of the trajectories in the configuration space of the variables $`a`$ and $`\varphi `$, which are the dynamical degrees of freedom of the minisuperspace, with the classical ones, allows one to identify the classical and quantum phases for the scalar-tensor cosmological models. For all three different values of the curvature of the spatial sections, the configuration space of the quantum solutions displays oscillating universes. However, these oscillating universes remain at the Planck scale and they cannot be considered as candidates for the description of the early Universe (they are more like baby universes), except for the unnatural choice $`|d|<<1`$ in the $`ϵ=0`$ case. There are also trajectories which correspond to universes which begin and end in singular states. Only for the flat case it is possible to have bouncing models.
In the bouncing models of the flat spatial section case, the scale factor has an infinite initial and final values, near which it behaves classically. As it approaches the singularity, the repulsive quantum effects lead to the bounce, avoiding the singularity. Such a scenario can be a candidate for the description of our early Universe, since it is free from the initial singularity and behaves classically in the asymptotic limit of large values for the scale factor. It is worth to note that this classical asymptotic limit corresponds to the stiff matter Friedmann universe which, according to Zel’dovich , is the most promising one to describe the very early Universe.
The free scalar field model considered in this paper, on the other hand, can be connected to a non-minimal coupled scalar field, with a coupling parameter $`\omega `$, like in the Brans-Dicke theory, by a conformal transformation. In a quantum analysis of these non-minimal models was performed, and it was shown that non-singular scenarios can be obtained when the parameter $`\omega `$ is negative. In fact, all quantum analysis performed here can be connected with a similar analysis in the non-minimal case through a conformal transformation, namely $`\alpha _{NMC}=\alpha +\varphi `$ and $`\varphi _{NMC}=\varphi `$. One can verify that when the minimal model displays singularities, it is possible to have non-singular solutions in the corresponding non-minimal case; but the non-singular solutions in the minimal case must also be non-singular in the corresponding non-minimal models.
An important generalization of the model studied here would be to consider self-interacting scalar fields. It was shown in that to each perfect fluid barotropic equation of state, it is possible to construct a self-interacting scalar field model leading to the same classical description in minisuperspace. We may argue if this correspondence remains at the quantum level. Note that in Ref. we have obtained bouncing universes for radiation fields with $`ϵ=0`$ and $`ϵ=1`$ with the same qualitative behaviour as the bouncing universes found here. They are also viable models for the early Universe. One should investigate if a scalar field model with a potential corresponding to the radiation fluid would give similar results.
## ACKNOWLEDGMENTS
We would like to thank Andrew Sornborger and the Cosmology Group of CBPF for many useful discussions, and CNPq and CAPES of Brazil for financial support. One of us (NPN) would like to thank the Laboratoire de Gravitation et Cosmologie Relativistes of Université Pierre et Marie Curie, and Fermilab, where part of this work has been done, for financial aid and hospitality.
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# 1 Introduction
## 1 Introduction
Weak decays of hyperons have been examined using effective field theory methods for more than three decades , but still a number of mysteries exist. The matrix elements of nonleptonic hyperon decays, e.g., can be described in terms of two amplitudes — the parity-violating s-wave and the parity-conserving p-wave. Chiral perturbation theory provides a framework whereby these amplitudes can be expanded in terms of small four-momenta and the current masses $`m_q`$ of the light quarks, $`q=u,d,s`$. At lowest order in this expansion the amplitudes are expressed in terms of two unknown coupling constants, so-called low-energy constants (LECs). However, there is no consensus for the determination of these two weak parameters. If one employs values which provide a good fit for the s-waves, one obtains a poor description of the p-waves. On the other hand, a good p-wave representation yields a poor s-wave fit. In order to overcome this problem, one must go beyond leading order. In Refs. , a first attempt was made in calculating the leading chiral corrections to such decays. But the inability to fit s- and p-waves simultaneously remains even after including the lowest nonanalytic contributions. In a recent paper a calculation was performed which included all terms at one-loop order. An exact fit to the data was possible but not unique, and other model-dependent assumptions had to be made in order to estimate the LECs. Another intriguing possibility was examined by Le Yaouanc et al., who asserted that a reasonable fit for both s- and p-waves can be provided by appending pole contributions from $`SU(6)(70,1^{})`$ states to the s-waves . Their calculations were performed in a simple quark model and appear to be able to provide a resolution of the s- and p-wave dilemma. The validity of this approach has been confirmed within the framework of chiral perturbation theory, but only after contributions from the lowest lying 1/2<sup>+</sup> baryon octet resonant states were also taken into account .
Another topic of interest are the radiative hyperon decays: $`\mathrm{\Sigma }^+p\gamma ,\mathrm{\Lambda }n\gamma ,`$ etc. Here, the primary problem has been and remains to understand the size of the asymmetry parameter in polarized $`\mathrm{\Sigma }^+p\gamma `$ decay: $`\alpha _\gamma =0.76\pm 0.08`$ . The difficulty here is associated with the restrictions posed by Hara’s theorem, which requires the vanishing of this asymmetry in the U-spin symmetric limit . Of course, in the real world, U-spin is broken and one should not be surprised to find a nonzero value for the asymmetry — what is difficult to understand is its size. Recent work involving the calculation of chiral loops has also not lead to a resolution, although slightly larger asymmetries can be accomodated . Within that work it was claimed, that in order to obtain a better understanding for the decays one should account for all terms at one-loop order, i.e., including all counterterms, not just the leading log corrections. For a different approach including explicitly baryon resonances which leads to improved agreement with experimental data from radiative hyperon decays see Refs. . This might indicate that weak hyperon decays cannot be described appropriately without the inclusion of baryon resonances. We will not elaborate on this possiblity in the present investigation. Rather, one of the main purposes of our paper is to estimate the size of higher order counterterms which might shed some light in the understanding of the origin of the discrepancy between theory and experiment when baryon resonances are not taken into account.
At one-loop order, divergences appear and are absorbed by infinite LECs from counterterms of the same chiral order. The renormalized low-energy constants are scale-dependent and measurable, i.e., they can in principle be determined from a fit to some observables. They satisfy renormalization group equations under scale changes and, therefore, the choice of another scale leads to modified values of the renormalized LECs. The divergences of the generating functional determine the renormalization group equations and the behavior of the renormalized LECs under scale changes. The sum of the irreducible one-loop functional and the counterterm functional is, of course, finite and scale-independent. Some of these divergences were treated in Ref. , which, to our knowledge, is the only work in the weak baryonic sector that performed renormalization explicitly. Obviously, only a subset of the leading divergences were treated in that work. It is our aim to work out all leading divergences in the generating functional of the weak baryon-meson Lagrangian, thus extending the work of Müller and Meißner in the strong sector, using two different techniques and to discuss a few applications. The complete divergence structure might be used as a check in future calculations.
In the next section we present the weak Lagrangian at lowest order. The generating functional to one-loop order is worked out in Sec. 3. The divergent parts of the irreducible tadpole, self-energy and the so-called eye-graph are isolated by using heat kernel techniques. Sec. 4 contains a sample calculation of the divergent pieces of the diagrams for the nonleptonic hyperon decay $`\mathrm{\Lambda }p\pi ^{}`$. The divergent pieces of these diagrams are compared with the results from the heat kernel technique. In Sec. 5 we give an estimate for some counterterm contributions to the s-wave amplitudes in nonleptonic hyperon decays and we summarize our findings in Sec.6. In the Appendix we extend the recently proposed super heat kernel technique to the weak effective Lagrangian.
## 2 Effective Lagrangian
We perform our calculations using the lowest order effective Lagrangian within the heavy baryon formalism. To this end, one writes down the most general relativistic Lagrangian which is invariant under chiral and $`CPS`$ transformations. Imposing invariance of the Lagrangian under the transformation $`S`$, which interchanges down and strange quarks in the Lagrangian, one can further reduce the number of counterterms. We will work in the $`CP`$ limit so that all LECs are real. This Lagrangian is then reduced to the heavy fermion limit by the use of path integral methods, which deliver the relativistic corrections as $`1/\stackrel{}{m}`$ terms in higher orders. The baryons are described by a four-velocity $`v_\mu `$ and a consistent chiral counting scheme emerges, i.e., a one-to-one correspondence between the Goldstone boson loops and the expansion in small momenta and quark masses. However, we will not present the relativistic Lagrangian explicitly here, but rather quote only the form of the heavy baryon limit.
The pseudoscalar Goldstone fields ($`\varphi =\pi ,K,\eta `$) are collected in the $`3\times 3`$ unimodular, unitary matrix $`U(x)`$,
$$U(\varphi )=u^2(\varphi )=\mathrm{exp}\{2i\varphi /F_0\}$$
(1)
with $`F_0`$ being the pseudoscalar decay constant (in the chiral limit), and
$`\varphi ={\displaystyle \frac{1}{\sqrt{2}}}\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\pi ^0+\frac{1}{\sqrt{6}}\eta & \pi ^+& K^+\\ \pi ^{}& \frac{1}{\sqrt{2}}\pi ^0+\frac{1}{\sqrt{6}}\eta & K^0\\ K^{}& \overline{K^0}& \frac{2}{\sqrt{6}}\eta \end{array}\right).`$
Under SU(3)$`{}_{L}{}^{}\times `$SU(3)<sub>R</sub>, $`U(x)`$ transforms as $`UU^{}=LUR^{}`$, with $`L,R`$ SU(3)<sub>L,R</sub>. The matrix $`B`$ denotes the baryon octet,
$`B=\left(\begin{array}{ccc}\frac{1}{\sqrt{2}}\mathrm{\Sigma }^0+\frac{1}{\sqrt{6}}\mathrm{\Lambda }& \mathrm{\Sigma }^+& p\\ \mathrm{\Sigma }^{}& \frac{1}{\sqrt{2}}\mathrm{\Sigma }^0+\frac{1}{\sqrt{6}}\mathrm{\Lambda }& n\\ \mathrm{\Xi }^{}& \mathrm{\Xi }^0& \frac{2}{\sqrt{6}}\mathrm{\Lambda }\end{array}\right),`$
which under $`SU(3)_L\times SU(3)_R`$ transforms as any matter field,
$$BB^{}=KBK^{},$$
(2)
with $`K(U,L,R)`$ the compensator field representing an element of the conserved subgroup SU(3)<sub>V</sub>. At the order we are working, the effective Lagrangian has the form
$$_{\text{eff}}=_{\varphi B}+_{\varphi B}^W+_\varphi ,$$
(3)
where $`_\varphi `$ is the usual (strong and electromagnetic) mesonic Lagrangian, see, e.g., Ref. .
For the strong meson-baryon Lagrangian $`_{\varphi B}`$ one writes
$`_{\varphi B}=_{\varphi B}^{(1)}`$ $`=`$ $`\text{i}<\overline{B}[v,B]>+D<\overline{B}S_\mu \{u^\mu ,B\}>+F<\overline{B}S_\mu [u^\mu ,B]>`$ (4)
with $`2S_\mu =\text{i}\gamma _5\sigma _{\mu \nu }v^\nu `$ denoting the Pauli-Lubanski spin vector, $`<\mathrm{}>`$ is defined as the trace in flavor space and the superscript denotes the chiral order. At lowest order the meson-baryon Lagrangian contains two axial-vector couplings, denoted by $`D`$ and $`F`$. The covariant derivative $`_\mu `$ of the baryons is decomposed as
$$[_\mu ,B]=_\mu B+[\mathrm{\Gamma }_\mu ,B]$$
(5)
with $`\mathrm{\Gamma }_\mu `$ being the so-called chiral connection
$$\mathrm{\Gamma }_\mu =\frac{1}{2}\left[u^{}(_\mu ir_\mu )u+u(_\mu il_\mu )u^{}\right].$$
(6)
The external fields $`v_\mu ,a_\mu `$ appear in the combinations $`r_\mu =v_\mu +a_\mu `$ and $`l_\mu =v_\mu a_\mu `$. The meson fields are summarized in the quantity
$$u_\mu =\frac{1}{2}\left[u^{}(_\mu ir_\mu )uu(_\mu il_\mu )u^{}\right].$$
(7)
We will also need the expression $`\chi _+=4B_0+𝒪(\varphi ^2)`$ with $`=\text{diag}(m_u,m_d,m_s)`$ being the quark mass matrix and $`B_0=0|\overline{q}q|0/F_0^2`$ the order parameter of the spontaneous symmetry violation.
Having dealt with its strong counterpart, the weak meson-baryon Lagrangian $`_{\varphi B}^W`$ at lowest order is
$$_{\varphi B}^{W(0)}=d<\overline{B}\{h_+,B\}>+f<\overline{B}[h_+,B]>.$$
(8)
Here, we have defined
$$h_+=u^{}hu+u^{}h^{}u,$$
(9)
with $`h_b^a=\delta _2^a\delta _b^3`$ being the weak transition matrix. Note that $`h_+`$ transforms as a matter field.
## 3 Renormalization of the weak one-loop generating functional using standard heat kernel techniques
In this section, we turn to the calculation of the complete one-loop generating functional in SU(3) heavy baryon chiral perturbation theory, i.e., to order $`q^3`$ for the strong interaction and to order $`q^2`$ for the weak sector in the small momentum expansion. This extends the renormalization of the strong interacting functional to the weak interaction. The method used in this section was first proposed by Ecker in the two flavor case. We will focus our calculation on applications that are of physical interest, so we neglect weak hyperon decays into two or more pions. Radiative and nonleptonic hyperon decays are of specific interest.
For the calculation it is useful to write the fields in terms of the physical basis, e.g., $`B=B^a\lambda ^a`$ with $`<\lambda _{}^{a}{}_{}{}^{}\lambda ^b>=\delta ^{ab}`$. We can then write the meson-baryon interaction in the form
$$S_{\varphi B}=d^4x\overline{B}^a(A_{\mathrm{str}}^{ab}+A_W^{ab})B^b$$
(10)
with the strong and weak interaction pieces $`A_{\mathrm{str}}^{ab}`$ and $`A_W^{ab}`$, respectively. Following Refs. and one has to expand
$$_\varphi ^{(2)}+_\varphi ^{(4)}\overline{R}^a[A_{(1),\mathrm{str}}^{ab}+A_{(0),W}^{ab}]^1R^b$$
(11)
in the functional integral around the classical solution, $`u_{\mathrm{cl}}=u_{\mathrm{cl}}[j]`$, which is obtained by the variation $`\delta d^4x_\varphi ^{(2)}/\delta U`$ at lowest order. Here, $`j`$ collectively denotes the external fields $`v_\mu ,a_\mu `$ and the quark mass matrix $``$. The baryon source fields have also been decomposed into a light component, denoted by $`R`$, and a heavy component where the heavy components are not needed for a consistent renormalization and will therefore be omitted. One thus arrives at a set of irreducible and reducible diagrams. From Eq. (11) we immediately derive
$$_\varphi ^{(2)}+_\varphi ^{(4)}\overline{R}^a[A_{(1),\mathrm{str}}^{ab}]^1R^b+\overline{R}^a[A_{(1),\mathrm{str}}^{ab}]^1[A_{(0),W}^{ab}][A_{(1),\mathrm{str}}^{ab}]^1]R^b,$$
(12)
where the first baryon term leads to the well known tadpole, self-energy contributions and the last baryon term is the new eye-graph part with a weak insertion on the intermediate baryon line. This problem was solved in Refs. and for the strong interaction with one insertion from the second order Lagrangian. As in the mesonic sector, we choose the fluctuation variables $`\xi `$ in a symmetric form ,
$$\xi _R=u_{\mathrm{cl}}\mathrm{exp}\{i\xi /2\},\xi _L=u_{\mathrm{cl}}^{}\mathrm{exp}\{i\xi /2\},$$
(13)
with $`\xi ^{}=\xi `$ traceless 3$`\times `$3 matrices. Consequently, we also have
$$U=u_{\mathrm{cl}}\mathrm{exp}\{i\xi \}u_{\mathrm{cl}}.$$
(14)
At second order in $`\xi `$, the covariant derivative $`_\mu `$, the chiral connection $`\mathrm{\Gamma }_\mu `$, the axial-vector $`u_\mu `$ and $`h_+`$ take the form
$`\mathrm{\Gamma }_\mu `$ $`=`$ $`\mathrm{\Gamma }_\mu ^{\mathrm{cl}}+{\displaystyle \frac{1}{4}}[u_\mu ^{\mathrm{cl}},\xi ]+{\displaystyle \frac{1}{8}}\xi \underset{\mu }{\overset{\mathrm{cl}}{\stackrel{}{}}}\xi +𝒪(\xi ^3),`$
$`[_\mu ^{\mathrm{cl}},\xi ]`$ $`=`$ $`_\mu \xi +[\mathrm{\Gamma }_\mu ^{\mathrm{cl}},\xi ],\xi \underset{\mu }{\overset{\mathrm{cl}}{\stackrel{}{}}}\xi =\xi [_\mu ^{\mathrm{cl}},\xi ][_\mu ^{\mathrm{cl}},\xi ]\xi ,`$
$`u_\mu `$ $`=`$ $`u_\mu ^{\mathrm{cl}}[_\mu ^{\mathrm{cl}},\xi ]+{\displaystyle \frac{1}{8}}[\xi ,[u_\mu ^{\mathrm{cl}},\xi ]]+𝒪(\xi ^3)`$
$`h_+`$ $`=`$ $`h_+^{\mathrm{cl}}+{\displaystyle \frac{i}{2}}[h_+^{\mathrm{cl}},\xi ]+{\displaystyle \frac{1}{8}}[\xi ,[h_+^{\mathrm{cl}},\xi ]]++𝒪(\xi ^3).`$ (15)
Inserting this into the expression for $`A^{ab}`$ and retaining only the terms up to and including order $`\xi ^2`$ gives
$`A_{(1),\mathrm{str}}^{ab}`$ $`=`$ $`A_{(1),\mathrm{str}}^{ab,\mathrm{cl}}+{\displaystyle \frac{i}{4}}<\lambda _{}^{a}{}_{}{}^{}[[vu_{\mathrm{cl}},\xi ],\lambda ^b]>D/F<\lambda _{}^{a}{}_{}{}^{}([S_{\mathrm{cl}},\xi ],\lambda ^b)_\pm >`$
$`+`$ $`{\displaystyle \frac{i}{8}}<\lambda _{}^{a}{}_{}{}^{}[\xi v\underset{\mathrm{cl}}{\overset{}{}}\xi ,\lambda ^b]>+{\displaystyle \frac{1}{8}}D/F<\lambda _{}^{a}{}_{}{}^{}([\xi ,[Su_{\mathrm{cl}},\xi ]],\lambda ^b)_\pm >+𝒪(\xi ^3),`$
$`A_{(0),W}^{ab}`$ $`=`$ $`A_{(1),W}^{ab,\mathrm{cl}}+{\displaystyle \frac{i}{2}}<\lambda _{}^{a}{}_{}{}^{}[h_+^{\mathrm{cl}},\xi ],\lambda ^b]>+{\displaystyle \frac{1}{8}}d/f<\lambda ^a^{}([\xi ,[h_+^{\mathrm{cl}},\xi ]],\lambda ^b)_\pm >+𝒪(\xi ^3),`$ (16)
where we have introduced the compact notation
$$D/F(\lambda ^a,\lambda ^b)_\pm =D\{\lambda ^a,\lambda ^b\}+F[\lambda ^a,\lambda ^b].$$
(17)
The corresponding generating functional reads
$`Z_{\mathrm{irr}}[j,\overline{R}^a,R^e]`$ $`=`$ $`{\displaystyle d^4xd^4x^{}d^4yd^4y^{}\overline{R}^a(x)S_{(1),\mathrm{str}}^{ac,\mathrm{cl}}(x,y)}`$
$`\times \left[\mathrm{\Sigma }_{\mathrm{tad}}^{cd}(y,y^{})\delta (yy^{})+\mathrm{\Sigma }_{\mathrm{self}}^{cd}(y,y^{})+\mathrm{\Sigma }_{\mathrm{eye}}^{cd}(y,y^{})\right]S_{(1),\mathrm{str}}^{de,\mathrm{cl}}(y^{},x^{})R^e(x^{})`$
in terms of the tadpole, self-energy and eye-graph functionals. Here, $`S_{(1),\mathrm{str}}^{\mathrm{cl}}`$ is the full classical fermion propagator with the weak interactions turned off. The functionals are given by
$`\mathrm{\Sigma }_{\mathrm{self}}^{ab}`$ $`=`$ $`{\displaystyle \frac{2}{F_0^2}}V_i^{ac}G_{ij}[A_{(1),\mathrm{str}}^{cd,cl}]^1V_j^{db}={\displaystyle \frac{2}{F_0^2}}V_i^{ac}G_{ij}S_{(1),\mathrm{str}}^{cd,cl}V_j^{db}`$
$`\mathrm{\Sigma }_{\mathrm{eye}}^{ab}`$ $`=`$ $`{\displaystyle \frac{2}{F_0^2}}V_i^{ac}G_{ij}[A_{(1),\mathrm{str}}^{ce,cl}][A_{(0),W}^{ef,cl}][A_{(1),\mathrm{str}}^{fd,cl}]^1V_j^{db}`$
$`\mathrm{\Sigma }_{\mathrm{tad}}^{ab}`$ $`=`$ $`{\displaystyle \frac{1}{8F_0^2}}\{D/F<\lambda _{}^{a}{}_{}{}^{}([\lambda _G^i,[Su^{cl},\lambda _G^j]],\lambda ^b)_\pm >G_{ij}`$ (19)
$`+d/f<\lambda _{}^{a}{}_{}{}^{}([\lambda _G^i,[h_+,\lambda _G^j]],\lambda ^b)_\pm >G_{ij}`$
$`+i<\lambda _{}^{a}{}_{}{}^{}[\lambda _G^i(G_{ij}v\underset{jk}{\overset{}{d}}vd_{ij}G_{jk})\lambda _G^k,\lambda ^b]>\}`$
with the following definitions of the vertices
$`V_i^{ab}`$ $`=`$ $`V_{i,\mathrm{str}}^{ab}+V_{i,W}^{ab},V_{i,\mathrm{str}}^{ab}=V_{i,\mathrm{str}}^{ab(1)}+V_{i,\mathrm{str}}^{ab(2)}`$
$`V_{i,\mathrm{str}}^{ab(1)}`$ $`=`$ $`{\displaystyle \frac{i}{4\sqrt{2}}}<\lambda _{}^{a}{}_{}{}^{}[[vu^{cl},\lambda _G^i],\lambda ^b]>,V_{i,\mathrm{str}}^{ab(2)}={\displaystyle \frac{D/F}{\sqrt{2}}}<\lambda _{}^{a}{}_{}{}^{}(\lambda _G^jSd_{ji},\lambda ^b)_\pm >`$
$`V_{i,W}^{ab}`$ $`=`$ $`{\displaystyle \frac{d/f}{2\sqrt{2}}}<\lambda _{}^{a}{}_{}{}^{}([h_+^{\mathrm{cl}},\lambda _G^j],\lambda ^b)_\pm >`$ (20)
where $`i,j,k,a,b,c=1,\mathrm{},8`$ and $`\lambda _G^i`$ denote Gell-Mann’s SU(3) matrices, which are related to the ones in the physical basis by $`\lambda _p=(\lambda _G^4+i\lambda _G^5)/2`$, $`\lambda _n=(\lambda _G^6+i\lambda _G^7)/2`$, etc. The quantity $`G_{ij}`$ is the full meson propagator
$$G_{ij}=(d_\mu d^\mu \delta ^{ij}+\sigma ^{ij})^1$$
(21)
with
$`[_{\mathrm{cl}}^\mu ,\xi ]`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}}}\lambda _G^jd_{jk}^\mu \xi _k,\xi ={\displaystyle \frac{1}{\sqrt{2}}}\lambda _G^i\xi _i,`$
$`d_{ij}^\mu `$ $`=`$ $`\delta _{ij}^\mu +\gamma _{ij}^\mu ,\underset{ij}{\overset{\mu }{\stackrel{}{d}}}=\delta _{ij}\stackrel{\mu }{\stackrel{}{}}\gamma _{ij}^\mu ,`$
$`\gamma _{ij}^\mu `$ $`=`$ $`{\displaystyle \frac{1}{2}}<\mathrm{\Gamma }_{\mathrm{cl}}^\mu [\lambda _G^i,\lambda _G^j]>,`$
$`\sigma ^{ij}`$ $`=`$ $`{\displaystyle \frac{1}{8}}<[u_\mu ^{\mathrm{cl}},\lambda _G^i][\lambda _G^j,u_{\mathrm{cl}}^\mu ]+\chi _+\{\lambda _G^i,\lambda _G^j\}>.`$ (22)
Note that the differential operator $`d_{ij}`$ is related to the covariant derivative $`_{\mathrm{cl}}^\mu `$ and it acts on the meson propagator $`G_{ij}`$. The connection $`\gamma _\mu `$ defines a field strength tensor,
$`\gamma _{\mu \nu }`$ $`=`$ $`_\nu \gamma _\mu _\mu \gamma _\nu +[\gamma _\mu ,\gamma _\nu ],`$ (23)
where we have omitted the flavor indices. We are now in a position to derive the divergences of the one-loop generating functional for weak and strong interaction after contracting trace indices in flavor space. In the heat kernel representation, the divergences appear as simple poles in $`ϵ=4d`$. The beta functions of the strong interacting functional are given in Ref. and the new divergent contribution of the weak interaction can be cast in the form (rotated back to Minkowski space)
$$\mathrm{\Sigma }_{\mathrm{weak}}^{ab,\mathrm{div}}(y,y)=\frac{1}{(4\pi F_0)^2}\frac{1}{ϵ}[\widehat{\mathrm{\Sigma }}_{\mathrm{tad}}^{ab}(y,y)+\widehat{\mathrm{\Sigma }}_{\mathrm{self}}^{ab}(y,y)+\widehat{\mathrm{\Sigma }}_{\mathrm{eye}}^{ab}(y,y)],$$
(24)
where $`\widehat{\mathrm{\Sigma }}^{ab}(y,y)`$ are finite monomials in the fields of chiral dimension two. Let us start with the tadpole contribution which is given by
$`\widehat{\mathrm{\Sigma }}_{\mathrm{tad}}^{ab}(y,y)`$ $`=`$ $`{\displaystyle \frac{1}{4}}d/f\{{\displaystyle \frac{3}{2}}<\lambda _{}^{a}{}_{}{}^{}(\{h_+,\chi _+\},\lambda ^b)_\pm >`$ (25)
$``$ $`<\lambda _{}^{a}{}_{}{}^{}(h_+,\lambda ^b)_\pm ><\chi _+>\}{\displaystyle \frac{1}{2}}d<\lambda _{}^{a}{}_{}{}^{}\lambda ^b><h_+\chi _+>.`$
Notice that we have neglected terms of order $`u^2`$ which contain at least two external pions, since they do not contribute to nonleptonic and radiative hyperon decays. The self-energy contribution is given by
$`\widehat{\mathrm{\Sigma }}_{\mathrm{self}}^{ab}(y,y)`$ $`=`$ $`6(Df+Fd)<\lambda _{}^{a}{}_{}{}^{}\{[h_+,S^\mu \mathrm{\Gamma }_{\mu \nu }v^\nu ],\lambda ^b\}>`$ (26)
$`+`$ $`({\displaystyle \frac{10}{3}}Dd+6Ff)<\lambda _{}^{a}{}_{}{}^{}[[h_+,S^\mu \mathrm{\Gamma }_{\mu \nu }v^\nu ],\lambda ^b]>`$
$`+`$ $`{\displaystyle \frac{3}{2}}f<\lambda _{}^{a}{}_{}{}^{}\{vu,\{h_+,[iv,\lambda ^b]\}\}>+\text{ h.c.}`$
$`+`$ $`{\displaystyle \frac{3}{2}}f<\lambda _{}^{a}{}_{}{}^{}[vu,[h_+,[iv,\lambda ^b]]]>+\text{ h.c.}`$
$`+`$ $`{\displaystyle \frac{3}{2}}d<\lambda _{}^{a}{}_{}{}^{}\{vu,[h_+,[iv,\lambda ^b]]\}>+\text{ h.c.}`$
$`+`$ $`{\displaystyle \frac{3}{2}}d<\lambda _{}^{a}{}_{}{}^{}[vu,\{h_+,[iv,\lambda ^b]\}]>+\text{ h.c.}`$
$`+`$ $`3d<\lambda _{}^{a}{}_{}{}^{}\{[h_+,[iv,vu]],\lambda ^b\}>+3d<\lambda _{}^{a}{}_{}{}^{}\{[[iv,h_+],vu],\lambda ^b\}>`$
$`+`$ $`3f<\lambda _{}^{a}{}_{}{}^{}[[h_+,[iv,vu]],\lambda ^b]>+3f<\lambda _{}^{a}{}_{}{}^{}[[[iv,h_+],vu],\lambda ^b]>`$
$`+`$ $`4f<\lambda _{}^{a}{}_{}{}^{}[iv,\lambda ^b]><vuh_+>+\text{ h.c.}`$
$`+`$ $`4f<\lambda _{}^{a}{}_{}{}^{}h_+><vu[iv,\lambda ^b]>+\text{ h.c. }.`$
Here, $`\mathrm{\Gamma }_{\mu \nu }`$ denotes the field strength tensor of the fields $`\mathrm{\Gamma }_\mu `$. The last contribution stems from the eye-graph and takes the form
$`\widehat{\mathrm{\Sigma }}_{\mathrm{eye}}^{ab}(y,y)`$ $`=`$ $`[S^\mu ,S^\nu ]\{(4(D^2+F^2)f8DFd)<\lambda _{}^{a}{}_{}{}^{}\lambda ^b><h_+\mathrm{\Gamma }_{\mu \nu }>`$ (27)
$`+`$ $`4f(D^2F^2)<\lambda _{}^{a}{}_{}{}^{}h_+><\mathrm{\Gamma }_{\mu \nu }\lambda ^b>+\text{ h.c.}`$
$`+`$ $`(3DFd{\displaystyle \frac{3}{2}}f(D^2+F^2))<\lambda _{}^{a}{}_{}{}^{}\{\mathrm{\Gamma }_{\mu \nu }\{h_+,\lambda ^b\}\}>+\text{ h.c.}`$
$`+`$ $`({\displaystyle \frac{1}{3}}DFd+{\displaystyle \frac{3}{2}}f(D^2+F^2))<\lambda _{}^{a}{}_{}{}^{}[\mathrm{\Gamma }_{\mu \nu }[h_+,\lambda ^b]]>+\text{ h.c.}`$
$`+`$ $`(3DFf{\displaystyle \frac{3}{2}}d(D^2+F^2))<\lambda _{}^{a}{}_{}{}^{}\{\mathrm{\Gamma }_{\mu \nu }[h_+,\lambda ^b]\}>+\text{ h.c.}`$
$`+`$ $`({\displaystyle \frac{17}{3}}DFf+{\displaystyle \frac{3}{2}}d(D^2+F^2))<\lambda _{}^{a}{}_{}{}^{}[\mathrm{\Gamma }_{\mu \nu }\{h_+,\lambda ^b\}]>+\text{ h.c. }\}`$
$`+`$ $`(6D^2d18F^2d+36DFf)<[\lambda _{}^{a}{}_{}{}^{},v\stackrel{}{}]\{h_+,[v,\lambda ^b]\}>`$
$`+`$ $`(10D^2f18F^2f+20DFd)<[\lambda _{}^{a}{}_{}{}^{},v\stackrel{}{}][h_+,[v,\lambda ^b]]>`$
$`+`$ $`(2D^2d+6F^2d12DFf)<\lambda _{}^{a}{}_{}{}^{}\{[v,[v,h_+]],\lambda ^b\}>`$
$`+`$ $`({\displaystyle \frac{10}{3}}D^2f+6F^2f{\displaystyle \frac{20}{3}}DFd)<\lambda _{}^{a}{}_{}{}^{}[[v,[v,h_+]],\lambda ^b]>`$
$``$ $`{\displaystyle \frac{3}{16}}\{({\displaystyle \frac{68}{9}}D^2d+4F^2d8DFf)<\lambda _{}^{a}{}_{}{}^{}\{h_+,\lambda ^b\}><\chi _+>`$
$`+`$ $`({\displaystyle \frac{68}{9}}D^2f+4F^2f8DFd)<\lambda _{}^{a}{}_{}{}^{}[h_+,\lambda ^b]><\chi _+>`$
$`+`$ $`(8(D^2+F^2)d+8DFf)<\lambda _{}^{a}{}_{}{}^{}\lambda ^b><h_+\chi _+>`$
$`+`$ $`({\displaystyle \frac{136}{9}}D^2d8F^2d)<\lambda _{}^{a}{}_{}{}^{}h_+><\chi _+\lambda ^b>+\text{ h.c.}`$
$`+`$ $`({\displaystyle \frac{23}{3}}D^2d6DFf+3F^2d)<\lambda _{}^{a}{}_{}{}^{}\{\chi _+\{h_+,\lambda ^b\}\}>+\text{ h.c.}`$
$`+`$ $`(3D^2d+{\displaystyle \frac{2}{3}}DFf3F^2d)<\lambda _{}^{a}{}_{}{}^{}[\chi _+[h_+,\lambda ^b]]>+\text{ h.c.}`$
$`+`$ $`({\displaystyle \frac{7}{3}}D^2f{\displaystyle \frac{2}{3}}DFd+3F^2f)<\lambda _{}^{a}{}_{}{}^{}\{\chi _+[h_+,\lambda ^b]\}>+\text{ h.c.}`$
$`+`$ $`(3D^2f+{\displaystyle \frac{2}{3}}DFd3F^2f)<\lambda _{}^{a}{}_{}{}^{}[\chi _+\{h_+,\lambda ^b\}]>+\text{ h.c. }\}`$
$`+`$ $`({\displaystyle \frac{68}{9}}D^3d12D^2Ff12DF^2d4F^3f)<\lambda _{}^{a}{}_{}{}^{}[iv,\lambda ^b]><Suh_+>+\text{ h.c.}`$
$`+`$ $`(4D^3d+4D^2Ff+4DF^2d4F^3f)<\lambda _{}^{a}{}_{}{}^{}h_+><Su[iv,\lambda ^b]>+\text{ h.c.}`$
$`+`$ $`(3D^3d+3D^2Ff+3DF^2d3F^3f)<\lambda _{}^{a}{}_{}{}^{}\{h_+,\{Su[iv,\lambda ^b]\}\}>+\text{ h.c.}`$
$`+`$ $`({\displaystyle \frac{1}{3}}D^3d+5D^2Ff+{\displaystyle \frac{13}{3}}DF^2d3F^3f)<\lambda _{}^{a}{}_{}{}^{}[h_+,[Su[iv,\lambda ^b]]]>+\text{ h.c.}`$
$`+`$ $`({\displaystyle \frac{7}{3}}D^3f+D^2Fd+3DF^2f3F^3d)<\lambda _{}^{a}{}_{}{}^{}\{h_+,[Su[iv,\lambda ^b]]\}>+\text{ h.c.}`$
$`+`$ $`(D^3f{\displaystyle \frac{1}{3}}D^2Fd+3DF^2f3F^3d)<\lambda _{}^{a}{}_{}{}^{}[h_+,\{Su[iv,\lambda ^b]\}]>+\text{ h.c.}`$
$`+`$ $`({\displaystyle \frac{2}{9}}D^3f{\displaystyle \frac{26}{9}}D^2Fd2DF^2f+2F^3d)<\lambda _{}^{a}{}_{}{}^{}\{[[iv,h_+],Su],\lambda ^b\}>`$
$`+`$ $`({\displaystyle \frac{2}{3}}D^3d2D^2Ff+{\displaystyle \frac{14}{9}}DF^2d+2F^3f)<\lambda _{}^{a}{}_{}{}^{}[[[iv,h_+],Su],\lambda ^b]>`$
$`+`$ $`({\displaystyle \frac{2}{9}}D^3f+{\displaystyle \frac{26}{9}}D^2Fd+2DF^2f2F^3d)<\lambda _{}^{a}{}_{}{}^{}\{[h_+,[iv,Su]],\lambda ^b\}>`$
$`+`$ $`({\displaystyle \frac{2}{3}}D^3d+2D^2Ff+{\displaystyle \frac{22}{9}}DF^2d2F^3f)<\lambda _{}^{a}{}_{}{}^{}[[h_+,[iv,Su]],\lambda ^b]>.`$
It is now straightforward to pin down the full weak counterterm Lagrangian at order $`q^2`$. We also use the curvature relation for $`\mathrm{\Gamma }_{\mu \nu }`$
$$\mathrm{\Gamma }_{\mu \nu }=\frac{1}{4}[u_\mu ,u_\nu ]\frac{i}{2}F_{\mu \nu }^+$$
(28)
where $`F_{\mu \nu }^+=uF_{\mu \nu }^Lu^{}+u^{}F_{\mu \nu }^Ru`$ with $`F_{\mu \nu }^{L/R}`$ being the field strength tensors related to the external fields $`l_\mu `$ and $`r_\mu `$, respectively. The generating functional can be renormalized by introducing the counterterm Lagrangian for the strong interaction
$$_{\varphi B}^{(3)}(x)=\frac{1}{(4\pi F_0)^2}\underset{i}{}d_i\overline{B}^a(x)\stackrel{~}{O}_{i,\mathrm{str}}^{ab}(x)B^b(x)$$
(29)
where the $`d_i`$ are dimensionless coupling constants and the field monomials $`\stackrel{~}{O}_{i,\mathrm{str}}^{ab}(x)`$ are of order $`q^3`$ and by introducing for the weak interaction the counterterm Lagrangian
$$_{\varphi B}^{(2)W}(x)=\frac{1}{(4\pi F_0)^2}\underset{i}{}h_i\overline{B}^a(x)\stackrel{~}{O}_{i,W}^{ab}(x)B^b(x)$$
(30)
where the $`h_i`$ are dimensionless coupling constants and the field monomials $`\stackrel{~}{O}_{i,W}^{ab}(x)`$ are of order $`q^2`$. The low-energy constants are decomposed into
$`d_i`$ $`=`$ $`d_i^r(\mu )+(4\pi )^2\beta _iL(\mu )`$
$`h_i`$ $`=`$ $`h_i^r(\mu )+(4\pi )^2\beta _iL(\mu ),`$ (31)
with $`\mu `$ being the scale introduced in dimensional regularization and
$$L(\mu )=\frac{\mu ^{d4}}{(4\pi )^2}\left(\frac{1}{d4}\frac{1}{2}[\mathrm{log}(4\pi )+1\gamma ]\right),$$
(32)
where $`\gamma =0.5772215\mathrm{}`$ is the Euler-Mascheroni constant. The $`\beta _i`$ are dimensionless functions of $`F`$, $`D`$ and $`f`$, $`d`$ constructed such that they cancel the divergences of the one-loop functional. The renormalized LECs $`d_i^r(\mu )(h_i^r(\mu ))`$ are measurable quantities. They satisfy the renormalization group equations and therefore, the choice of another scale leads to modified values of the renormalized LECs. We remark that the scale-dependence in the counterterm Lagrangian is, of course, balanced by the scale-dependence of the renormalized finite one-loop functional for observable quantities.
The next two sections are devoted to the application of the formulae derived in the present investigation. First, we will show how the divergences can be used to check results obtained from the ordinary computation of Feynman diagrams. Second, the scale-dependence of the renormalized LECs as given in Eq. (3) can be employed to make an estimate of the size of such higher order counterterms.
## 4 A sample calculation
In this section we will discuss the renormalization of the nonleptonic hyperon decay $`\mathrm{\Lambda }p\pi ^{}`$ by calculating explicitly the Feynman graphs and comparing them with the results obtained in the previous section. In the rest frame of the heavy baryon, $`v_\mu =(1,0,0,0)`$, the decay amplitude reduces to the non-relativistic form
$`𝒜(B_iB_j\pi )`$ $`=`$ $`\overline{u}_{B_j}\left\{𝒜_{ij}^{(s)}+Sk𝒜_{ij}^{(p)}\right\}u_{B_i},`$ (33)
where $`k`$ is the outgoing momentum of the pion. Here, $`𝒜_{ij}^{(s)}`$ is the parity-violating s-wave amplitude and $`𝒜_{ij}^{(p)}`$ is the corresponding parity-conserving p-wave term. In this frame, the energy of the outgoing pion is
$$vk=\frac{1}{2m_i}\left(m_i^2m_j^2+M_\pi ^2\right)$$
(34)
and the energy of the decaying hyperon in the heavy baryon formulation can be written as
$$vp=m_i\stackrel{}{m}.$$
(35)
Here, $`m_{i,j}`$ are the physical masses of the baryons and $`\stackrel{}{m}`$ is the mass of the baryon octet in the chiral limit. Since the baryon masses are analytic to linear order in the quark masses we see that $`vp`$ and $`vk`$ count effectively as order $`𝒪(q^2)`$. We will therefore restrict ourselves to the computation of the one-loop divergences to this decay which are proportional to the quark mass matrix and do not contain the momenta of the external particles. In our example we present only the renormalization of the weak vertices. The divergent structure of the purely strong baryonic Lagrangian which also contributes to this decay is already given in Ref. .
We start by renormalizing the p-wave amplitude. There are four Feynman graphs contributing to the p-wave and containing mass-dependent divergences (see Fig. 1). For the renormalization procedure it is sufficient to consider only the irreducible tadpole from Figures 1a and 1b and the irreducible eye-graph from Figs. 1c and 1d by neglecting the parts from the internal baryon propagator and the strong vertex. The contributions $`P^a`$ and $`P^b`$ of the irreducible tadpole in Figures 1a and 1b to the decay then read
$`P^a`$ $`=`$ $`{\displaystyle \frac{i}{4F_0^2}}(df)L(\mu )\left(3M_\pi ^2+6M_K^2+3M_\eta ^2\right)`$
$`P^b`$ $`=`$ $`{\displaystyle \frac{i}{4F_0^2}}{\displaystyle \frac{1}{\sqrt{6}}}(d+3f)L(\mu )\left(3M_\pi ^2+6M_K^2+3M_\eta ^2\right).`$ (36)
The divergent pieces can be recovered by using the results from the previous section. Employing the counterterms from Eq. (25) and using the Gell-Mann-Okubo relation for the pseudoscalar mesons, $`4M_K^2=3M_\eta ^2+M_\pi ^2`$, which is consistent to the order we are working, cancels the divergences from the calculation of the one-loop graphs. Note that the last term in Eq. (25) does not contribute here, since $`h_+\chi _+=0`$.
The mass-dependent divergences from diagrams 1c and 1d read after neglecting the internal baryon propagator and the strong vertex
$`P^c`$ $`=`$ $`{\displaystyle \frac{i}{F_0^2}}L(\mu )(M_\pi ^2[{\displaystyle \frac{2}{3}}D^2d8DFf+6F^2d{\displaystyle \frac{10}{3}}D^2f+4DFd6F^2f]`$
$`+M_K^2[{\displaystyle \frac{7}{3}}D^2d10DFf+3F^2d{\displaystyle \frac{5}{3}}D^2f+6DFd3F^2f])`$
$`P^d`$ $`=`$ $`{\displaystyle \frac{i\sqrt{6}}{F_0^2}}L(\mu )(M_\pi ^2[{\displaystyle \frac{14}{9}}D^2d2DFf{\displaystyle \frac{4}{3}}D^2f+{\displaystyle \frac{4}{3}}DFd]`$ (37)
$`+M_K^2[{\displaystyle \frac{19}{18}}D^2d+5DFf{\displaystyle \frac{3}{2}}F^2d{\displaystyle \frac{7}{6}}D^2f+{\displaystyle \frac{11}{3}}DFd{\displaystyle \frac{9}{2}}F^2f]).`$
The divergent pieces are in agreement with the mass-dependent divergences in Eq. (27).
The renormalization of the s-wave is somewhat more subtle. In this case it turns out that in addition to the simple irreducible tadpole diagram in Fig. 2a contributing to the s-wave there are three more diagrams (2b,c,d) which lead to the same divergent structure. The pertinent divergent contributions which are proportional to the quark mass matrix for Fig. 2a read
$$S^a=\frac{1}{24\sqrt{3}F_0^3}(d+3f)L(\mu )\left(5M_\pi ^2+16M_K^2+9M_\eta ^2\right)$$
(38)
and for Figures 2c and 2d
$$S^c+S^d=\frac{1}{4\sqrt{3}F_0^3}(d+3f)L(\mu )\left(2M_\pi ^2+M_K^2\right).$$
(39)
The sum of graphs 2a, 2c and 2d leads to a divergence structure different than that obtained from the tadpole occuring in the p-wave. Agreement between s- and p-waves is only achieved if one includes the $`Z`$-factor of the pion which is divergent and, therefore, contributes to the divergences for nonleptonic hyperon decay. This was not treated correctly in Ref. . From an explicit calculation of the pion $`Z`$-factor including the mesonic Lagrangian of fourth chiral order $`_\varphi ^{(4)}`$ one obtains
$$\sqrt{Z_\pi }=\frac{1}{F_0^2}L(\mu )\frac{2}{3}\left(2M_\pi ^2+M_K^2\right)+\text{finite term}$$
(40)
which is multiplied by the pertinent tree level contribution to the decays
$$\frac{1}{2\sqrt{3}F_0}(d+3f).$$
(41)
Expanding $`h_+`$ in Eq. (25) in terms of the meson fields
$$h_+=h\frac{i}{F_0}[\varphi ,h]+𝒪(\varphi ^2,h^{})$$
(42)
the term which is linear in the meson fields reproduces the divergences of the one-loop contributions to the s-wave discussed above. Note that the pseudoscalar decay constant in the chiral limit $`F_0`$ is finite and does not lead to additional divergences.
Finally, Fig. 2e leads to the divergences
$`S^e`$ $`=`$ $`{\displaystyle \frac{\sqrt{3}}{F_0^3}}L(\mu )(M_\pi ^2[{\displaystyle \frac{14}{9}}D^2d2DFf{\displaystyle \frac{4}{3}}D^2f+{\displaystyle \frac{4}{3}}DFd]`$ (43)
$`+M_K^2[{\displaystyle \frac{19}{18}}D^2d+5DFf{\displaystyle \frac{3}{2}}F^2d{\displaystyle \frac{7}{6}}D^2f+{\displaystyle \frac{11}{3}}DFd{\displaystyle \frac{9}{2}}F^2f]).`$
This result agrees with the mass dependent divergences in Eq. (27).
## 5 Importance of counterterms
Finally, we would like to address the issue of the importance of higher order counterterms for nonleptonic hyperon decays. There exist seven such transitions: $`\mathrm{\Sigma }^+n\pi ^+,\mathrm{\Sigma }^+p\pi ^0,\mathrm{\Sigma }^{}n\pi ^{},\mathrm{\Lambda }p\pi ^{},\mathrm{\Lambda }n\pi ^0,\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{},\text{and}\mathrm{\Xi }^0\mathrm{\Lambda }\pi ^0`$. Isospin symmetry of the strong interactions implies the relations
$`𝒜(\mathrm{\Lambda }p\pi ^{})+\sqrt{2}𝒜(\mathrm{\Lambda }n\pi ^0)=0`$
$`𝒜(\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{})+\sqrt{2}𝒜(\mathrm{\Xi }^0\mathrm{\Lambda }\pi ^0)=0`$
$`\sqrt{2}𝒜(\mathrm{\Sigma }^+p\pi ^0)+𝒜(\mathrm{\Sigma }^{}n\pi ^{})𝒜(\mathrm{\Sigma }^+n\pi ^+)=0`$ (44)
which hold both for s- and p-waves. We choose $`\mathrm{\Sigma }^+n\pi ^+,\mathrm{\Sigma }^{}n\pi ^{},\mathrm{\Lambda }p\pi ^{},\text{and}\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{}`$ to be the four independent decay amplitudes which are not related by isospin. As mentioned in the introduction, both a good s- and p-wave fit cannot be achieved just by taking the lowest order couplings $`d,f`$ and chiral corrections into account. This suggests that one has to consider higher order counterterms, but there are so many of them that they cannot be determined uniquely . However, as we will show in this section the divergent structure of the one-loop functional can be used to give an estimate of the size of these counterterm contributions. In order to keep the presentation more compact, we restrict ourselves to the case of the s-waves. Neglecting divergences which are proportional to $`vk`$ and $`vq`$ (see last section) the possible counterterms are linear in $`\chi _+`$ and read
$`_{\varphi B}^{W(2,br)}`$ $`=`$ $`h_1\left\{\text{tr}\left(\overline{B}\{h_+,\{\chi _+,B\}\}\right)+\text{tr}\left(\overline{B}\{\chi _+,\{h_+,B\}\}\right)\right\}`$ (45)
$`+`$ $`h_2\left\{\text{tr}\left(\overline{B}[h_+,[\chi _+,B]]\right)+\text{tr}\left(\overline{B}[\chi _+,[h_+,B]]\right)\right\}`$
$`+`$ $`h_3\left\{\text{tr}\left(\overline{B}[h_+,\{\chi _+,B\}]\right)+\text{tr}\left(\overline{B}\{\chi _+,[h_+,B]\}\right)\right\}`$
$`+`$ $`h_4\left\{\text{tr}\left(\overline{B}\{h_+,[\chi _+,B]\}\right)+\text{tr}\left(\overline{B}[\chi _+,\{h_+,B\}]\right)\right\}`$
$`+`$ $`h_5\left\{\text{tr}\left(\overline{B}h_+\right)\text{tr}\left(\chi _+B\right)+\text{tr}\left(\overline{B}\chi _+\right)\text{tr}\left(h_+B\right)\right\}`$
$`+`$ $`h_6\text{tr}\left(\overline{B}[h_+,B]\right)\text{tr}\left(\chi _+\right)+h_7\text{tr}(\overline{B}\{h_+,B]\}\left)\text{tr}\right(\chi _+).`$
Note that we did not make any use of Cayley-Hamilton identities in order to have the same set of counterterms as in Eqs. (25) and (27). After renormalization of the mass-dependent divergences the contributions of these counterterms read
$`𝒜_{\mathrm{\Sigma }^+n}^{(s)}`$ $`=`$ $`0`$
$`𝒜_{\mathrm{\Sigma }^{}n}^{(s)}`$ $`=`$ $`{\displaystyle \frac{2\sqrt{2}}{F_\pi }}(M_\pi ^2[h_1^rh_2^rh_3^r+h_4^r{\displaystyle \frac{1}{2}}h_6^r+{\displaystyle \frac{1}{2}}h_7^r]`$
$`+M_K^2[h_1^r+h_2^rh_3^rh_4^rh_6^r+h_7^r])`$
$`𝒜_{\mathrm{\Lambda }p}^{(s)}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{3}F_\pi }}(M_\pi ^2[3h_1^r3h_2^r+h_3^rh_4^r+2h_5^r{\displaystyle \frac{3}{2}}h_6^r{\displaystyle \frac{1}{2}}h_7^r]`$
$`+M_K^2[5h_1^r+3h_2^r7h_3^r+h_4^r2h_5^r3h_6^rh_7^r])`$
$`𝒜_{\mathrm{\Xi }^{}\mathrm{\Lambda }}^{(s)}`$ $`=`$ $`{\displaystyle \frac{2}{\sqrt{3}F_\pi }}(M_\pi ^2[3h_1^r3h_2^rh_3^r+h_4^r+2h_5^r+{\displaystyle \frac{3}{2}}h_6^r{\displaystyle \frac{1}{2}}h_7^r]`$ (46)
$`+M_K^2[5h_1^r+3h_2^r+7h_3^rh_4^r2h_5^r+3h_6^rh_7^r])`$
where we have not shown explicitly the scale-dependence of the $`h_i^r`$. At the order we are working, we replace $`F_0`$ by the pion decay constant $`F_\pi =92.4`$ MeV since the differences between the two show up at higher orders only. The dependence on a chosen scale $`\mu _1`$ of these counterterm contributions is, of course, compensated by the scale-dependence of the renormalized finite one-loop functional for observable quantities. The choice of another scale $`\mu _2`$ leads to modified values of the renormalized LECs according to the equation
$$h_i^r(\mu _2)=h_i^r(\mu _1)+\beta _i\mathrm{log}\frac{\mu _1}{\mu _2}.$$
(47)
Varying the scale from the $`\rho `$-meson mass $`M_\rho =770`$ MeV to the $`\mathrm{\Delta }`$ mass $`M_\mathrm{\Delta }=1232`$ MeV one obtains the following differences in the counterterm contributions for the s-waves, in units of 10<sup>-7</sup>,
$`𝒜_{\mathrm{\Sigma }^{}n}^{(s)}(\mu _\rho )𝒜_{\mathrm{\Sigma }^{}n}^{(s)}(\mu _\mathrm{\Delta })`$ $`=`$ $`0.77(4.27)`$
$`𝒜_{\mathrm{\Lambda }p}^{(s)}(\mu _\rho )𝒜_{\mathrm{\Lambda }p}^{(s)}(\mu _\mathrm{\Delta })`$ $`=`$ $`0.60(3.25)`$
$`𝒜_{\mathrm{\Xi }^{}\mathrm{\Lambda }}^{(s)}(\mu _\rho )𝒜_{\mathrm{\Xi }^{}\mathrm{\Lambda }}^{(s)}(\mu _\mathrm{\Delta })`$ $`=`$ $`1.94(4.51)`$ (48)
which is about 19% for $`\mathrm{\Sigma }^{}n\pi ^{},\mathrm{\Lambda }p\pi ^{}`$ and even 43% for $`\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{}`$; the experimental values for the decays being given in brackets. There are no contributions at one-loop order for the decay $`\mathrm{\Sigma }^+n\pi ^+`$ , neither from the counterterms nor from the chiral loops; we have therefore omitted its presentation here. Assuming that the differences of these counterterm contributions for different realistic choices of the scale give an estimate of their absolute value $`\delta 𝒜^{(s)}`$, we obtain, in units of 10<sup>-7</sup>,
$`|\delta 𝒜_{\mathrm{\Sigma }^{}n}^{(s)}|`$ $`=`$ $`0.77`$
$`|\delta 𝒜_{\mathrm{\Lambda }p}^{(s)}|`$ $`=`$ $`0.60`$
$`|\delta 𝒜_{\mathrm{\Xi }^{}\mathrm{\Lambda }}^{(s)}|`$ $`=`$ $`1.94.`$ (49)
While the counterterm contributions to the s-waves seem to be well behaved for $`\mathrm{\Sigma }^{}n\pi ^{}`$ and $`\mathrm{\Lambda }p\pi ^{}`$, our calculation indicates that they might be significant for $`\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{}`$.
## 6 Summary
In this paper, we have performed the chiral invariant renormalization of the weak effective baryon-meson Lagrangian up to one-loop order within heavy baryon chiral perturbation theory. The complete set of counterterms at leading one-loop order $`q^2`$ with $`q`$ being an external momentum or meson mass has been constructed. This extends work by Müller and Meißner , who considered the strong $`SU(3)`$ baryon-meson Lagrangian. The present calculation has been performed both using the standard heat kernel formalism and the recently proposed super heat kernel method which has the advantage of simplifying the calculation in intermediate steps. We also compared our results with a direct calculation of the Feynman graphs for the nonleptonic hyperon decay $`\mathrm{\Lambda }p\pi ^{}`$. It turns out that the divergences contained in the $`Z`$-factor of the pion are essential for achieving agreement between the renormalization of s- and p-wave amplitudes. Since, to our knowledge, there exists only one calculation in the weak baryon-meson sector where some of the divergences at order $`𝒪(q^2)`$ have been evaluated , our work might serve as a reference to check further calculations in the future.
The low-energy constants of higher order counterterms contain, in general, divergent pieces which cancel the divergences from one-loop graphs. The finite remainder of the counterterms is scale-dependent and compensates the scale-dependence of the one-loop functional to give scale independent expressions for physical quantities. For radiative and nonleptonic hyperon decays there exist more counterterms than there are experimental data , so that the counterterms cannot be determined uniquely from experiment. Nevertheless, the scale-dependence of the finite remainder of the LECs after renormalization can be used to give an estimate of the size of the counterterm contributions. Since the divergent structure of the one-loop functional determines the scale-dependence of the renormalized LECs, we are able to give such estimates and as an example we have chosen the contributions of the counterterms linear in the quark mass matrix $``$ to the s-waves of the independent nonleptonic hyperon decays that are not related by isospin. We find that for the s-waves the contributions from these counterterms relatively to the experimental values range from 19% for $`\mathrm{\Sigma }^{}n\pi ^{}`$ and $`\mathrm{\Lambda }p\pi ^{}`$ up to 43% for the decay $`\mathrm{\Xi }^{}\mathrm{\Lambda }\pi ^{}`$.
## Acknowledgments
Useful discussions with H. Neufeld and S. Steininger are gratefully acknowledged. This work was supported in part by the Deutsche Forschungsgemeinschaft and the BMBF.
## Appendix: Renormalization using the super heat kernel formalism
We provide an alternative renormalization prescription by using the super heat kernel formalism as proposed by Ecker, Gasser, and Neufeld . This formalism has already been applied to the effective two- and three-flavor Lagrangian in Refs. and , respectively, where it served as a check for previous work . We use this approach to calculate the divergences of the weak effective Lagrangian. This method allows us to verify our previous calculation presented in Sec.3.
The fluctuation action generated by the lowest order meson-baryon Lagrangian has the general form
$$S^{(2)}=\frac{1}{2}\xi ^T(D_\mu D^\mu +Y)\xi +\overline{\eta }(\alpha +\beta _\mu D^\mu )\xi +\xi ^T(\overline{\delta }\overline{\beta }_\mu 𝒟^\mu )\eta +\overline{\eta }iv^\mu 𝒟_\mu \eta ,$$
(50)
where $`\psi (\eta )`$ are the bosonic (fermionic) fluctuations and
$$\begin{array}{cc}D_\mu =_\mu +X_\mu ,\hfill & 𝒟_\mu =_\mu +f_\mu ,\hfill \\ \overline{\delta }=\overline{\alpha }\overline{𝒟_\mu \beta ^\mu },\hfill & \\ v^2=1,\hfill & v\beta =0.\hfill \end{array}$$
(51)
$`X_\mu `$, $`Y`$, and $`f_\mu `$ are bosonic (matrix) fields, whereas $`\alpha `$ and $`\beta _\mu `$ are fermionic objects. The form of $`\overline{\delta }`$ in Eq. (51) is required by the reality of Eq. (50). Apart from the condition $`v\beta =0`$ no further assumption about the terms entering in Eq. (50) is made. To be specific, in HBCHPT we find:
$`X_\mu `$ $`=`$ $`\gamma _\mu +g_\mu ,(g_\mu )^{ij}=i{\displaystyle \frac{v_\mu }{8F_0^2}}<\overline{B}[[\lambda _G^i,\lambda _G^j],B]>,i,j=1,\mathrm{},8,`$
$`Y`$ $`=`$ $`\sigma +s,`$
$`s^{ij}`$ $`=`$ $`{\displaystyle \frac{D/F}{4F_0^2}}<\overline{B}([\lambda _G^i,[Su,\lambda _G^j]],B)_\pm >{\displaystyle \frac{d/f}{4F_0^2}}<\overline{B}([\lambda _G^i,[Su,\lambda _G^j]],B)_\pm >,`$
$`f_\mu `$ $`=`$ $`f_\mu ^{\mathrm{str}}+f_\mu ^W,`$
$`f_\mu ^{ab,\mathrm{str}}`$ $`=`$ $`<\lambda _{}^{a}{}_{}{}^{}[\mathrm{\Gamma }_\mu ,\lambda ^b]>iv_\mu D/F<\lambda _{}^{a}{}_{}{}^{}(Su,\lambda ^b)_\pm >,`$
$`f_\mu ^{ab,\mathrm{W}}`$ $`=`$ $`iv_\mu d/f<\lambda _{}^{a}{}_{}{}^{}(h_+,\lambda ^b)_\pm >,a=1,\mathrm{},8,`$
$`\alpha ^{ai}`$ $`=`$ $`{\displaystyle \frac{i}{4F_0}}<\lambda _{}^{a}{}_{}{}^{}[[vu,\lambda _G^i],B]>+{\displaystyle \frac{i}{2F_0}}<\lambda _{}^{a}{}_{}{}^{}[[h_+,\lambda _G^i],B]>,`$
$`(\beta _\mu )^{ai}`$ $`=`$ $`{\displaystyle \frac{D/F}{F_0}}S_\mu <\lambda _{}^{a}{}_{}{}^{}(\lambda _G^i,B)_\pm >`$ (52)
with the definitions of Section 3.
Details of super heat kernel calculation in SU(3) can be found in Ref. . We only use the final result to calculate the structure of the divergences. Using these definitions we easily derive from the following formula the pertinent beta functions and the counterterms
$`W_{L=1}^{\mathrm{div}}|_{\overline{\mathrm{\Gamma }}\mathrm{}\mathrm{\Gamma }}`$ $`=`$ $`{\displaystyle \frac{i}{48\pi ^2(d4)}}{\displaystyle }d^4x\text{ tr }\{12\overline{\alpha }v\widehat{}\alpha +6[\overline{\alpha }\beta _\mu X^{\mu \nu }v_\nu +\overline{\beta }_\mu \alpha X^{\mu \nu }v_\nu ]`$ (53)
$`3\left[\overline{\beta }\beta v\widehat{}Y+2\overline{\beta }_\mu (v\widehat{}\beta ^\mu )Y\right]4\overline{\beta }_\mu (v\widehat{})^3\beta ^\mu +\overline{\beta }\beta \widehat{}_\mu X^{\mu \nu }v_\nu `$
$`+6\overline{\beta }_\mu (v\widehat{}\beta _\nu )X^{\mu \nu }+4\overline{\beta }_\mu \beta _\nu v\widehat{}X^{\mu \nu }+2\overline{\beta }_\mu \beta _\nu \widehat{}^\mu X^{\nu \rho }v_\rho \}.`$
where $`X^{\mu \nu }`$ is the field strength tensor of the fields $`X^\mu `$. This result can be used since the weak interaction has a relatively simple structure and does not contain any divergences which act on the bosonic or fermionic fluctuation variables. The new eye-graph contribution can be evaluated by redefining the covariant derivative in the most economic way, i.e., $`f_\mu =f_\mu ^{\mathrm{str}}+f_\mu ^W`$. This shows that the eye-graph without any derivative is immediately obtained in the super heat kernel formalism as it is the case in the standard heat kernel method. But the situation changes dramatically in the presence of derivatives. So far the result is given in a very compact form, and both the beta functions and the corresponding monomials must be evaluated by contracting the various trace indices. In the SU(3) case this is the most tedious part of the calculation and therefore, we skip the presentation of these technicalities. After a lot of cumbersome algebra one ends up with the result presented in the last section.
## Figure captions
1. Given are the diagrams which contribute to the mass-dependent divergences of the p-wave in the nonleptonic hyperon decay $`\mathrm{\Lambda }p\pi ^{}`$. Solid and dashed lines denote baryons and pseudoscalar mesons, respectively. The solid square represents a weak vertex and the solid circle denotes a strong vertex.
2. Given are the diagrams which contribute to the mass-dependent divergences of the s-wave in the nonleptonic hyperon decay $`\mathrm{\Lambda }p\pi ^{}`$. Solid and dashed lines denote baryons and pseudoscalar mesons, respectively. The solid square represents a weak vertex and the solid circle denotes a strong vertex.
Figure 1
Figure 2
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# Attractor states and infrared scaling in de Sitter space
## I Introduction
Quantum field theory in curved spacetimes does not contain in itself a unique specification of the quantum state of the system . Even in Minkowski spacetime, where the existence of the Poincaré group singles out a special state, the Minkowski vacuum, it is certainly of interest to consider states that are non-invariant under Poincaré transformations, since they contain all the information about the physical excitations and dynamics of the theory. Such non-vacuum states are also necessary in a general initial value formulation of the back-reaction problem in both curved and flat spacetimes. In flat space the initial value problem for arbitrary physically allowable states has been formulated and studied for both QED and scalar $`\mathrm{\Phi }^4`$ theory in the large $`N`$ limit, principally for time varying but spatially homogeneous mean fields .
The simplest situation in which the back-reaction problem can be studied in curved spacetime is that of a free scalar field in a spatially homogeneous and isotropic Robertson-Walker (RW) cosmology, where the geometry is characterized by just one non-trivial function of time. The wave equation for a free scalar field in such a geometry can be separated and expressed in terms of a complete set of time dependent mode functions. The general initial value problem is specified by giving initial data for this complete set at a given initial time. The back-reaction of the quantum scalar field(s) on the RW geometry can be studied then by constructing the renormalized expectation value of the energy-momentum tensor $`T_{ab}`$ of the field(s) and solving (numerically) the semi-classical Einstein equations, augmented by higher derivative terms required by renormalization. As in the flat space examples, this semi-classical back-reaction problem becomes exact in the large $`N`$ limit, with $`N`$ the number of identical scalar fields .
As a prelude to the dynamical back-reaction problem in cosmological spacetimes it is necessary to study non-vacuum states first in fixed RW backgrounds. The maximally symmetric de Sitter spacetime is of particular interest. Most previous work has focused on maximally symmetric $`O(4,1)`$ de Sitter invariant states or the special $`O(4)`$ invariant state found by Allen . Since the universe is not globally $`O(4,1)`$ invariant, a more generic set of initial conditions, consistent only with RW symmetry and general principles of renormalization of $`T_{ab}`$ is required for cosmology. The investigation of these much weaker requirements and specification of the general initial value problem for back-reaction calculations was initiated in Ref. .
In this paper we study the behavior of the renormalized $`T_{ab}`$ for arbitrary physically admissable spatially homogeneous and isotropic states in a fixed de Sitter background. We argue in Section III that such states must be fourth order adiabatic states that also possess an infrared finite two-point function. In de Sitter space the wave equation for free scalar fields can be solved exactly for arbitrary values of the mass and the curvature coupling. Its solutions depend only on the wave number $`k`$ of the mode and the parameter $`\nu ^2=\frac{9}{4}m^2\alpha ^212\xi `$, with $`R=12\alpha ^2`$ the constant scalar curvature of de Sitter spacetime. For $`\mathrm{}(\nu )<\frac{3}{2}`$, corresponding to $`m^2+\xi R>0`$, we prove that for all UV and IR physically allowed initial states the renormalized value of $`T_{ab}`$ at late times asymptotically approaches that of the Euclidean or Bunch-Davies de Sitter invariant state . The conformally invariant scalar field ($`m=0`$$`\xi =\frac{1}{6}`$) falls into this class.
The case $`\nu =\frac{3}{2}`$ corresponding to $`m^2\alpha ^2+12\xi =0`$ is more delicate. If $`m`$ and $`\xi `$ are separately zero (the massless, minimally coupled case), then we prove that the renormalized $`T_{ab}`$ for all physically admissable states approaches the Allen-Folacci de Sitter invariant value . Numerical evidence for this result was found previously in Ref. . In this paper we provide an analytic proof that late time attractor behavior occurs for all physically admissable RW states, when $`m^20`$ and $`\xi 0`$. If $`m^2\alpha ^2+12\xi =0`$ but $`m^2`$ and $`\xi `$ are not separately zero (so that one of them is negative) we prove that $`T_{ab}`$ grows linearly in RW comoving (cosmic) time without bound, and this asymptotic behavior is independent of the state of the field. Finally, and in contrast, in the case $`\nu >\frac{3}{2}`$, corresponding to $`m^2+\xi R<0`$, $`T_{ab}`$ depends sensitively on the state and, for most values of $`m`$ and $`\xi `$, grows exponentially at late times for all states. This case is of considerably less physical relevance, since it corresponds to a tachyonic field theory with no stable vacuum state.
The asymptotic approach of $`T_{ab}`$ to a de Sitter invariant form, independently of the lower symmetry of the initial data when $`\mathrm{}(\nu )\frac{3}{2}`$ is a striking result. Certainly no such attractor behavior of $`T_{ab}`$, independent of initial conditions occurs in Minkowski space for any mass. One may regard this result as a kind of cosmic “no hair” theorem for scalar quantum fields in de Sitter space. For $`\mathrm{}(\nu )<\frac{3}{2}`$ it is in accord with one’s classical intuition that any initial energy density satisfying the weak energy condition ($`\epsilon +p>0`$) is redshifted away by the exponential de Sitter expansion, although as we will see, the redshifting of the quantum $`T_{ab}`$ is not that of classical matter or radiation. At asymptotically late times what is left behind is a kind of frozen quantum vacuum energy “condensate,” satisfying the de Sitter invariant equation of state $`p=\epsilon `$. This result justifies the choice of the Bunch-Davies vacuum in calculations of quantum fluctuations of free fields, i.e. without back-reaction, in a long-lived de Sitter expansion phase of inflationary cosmological models. For $`\nu =\frac{3}{2}`$, $`m=\xi =0`$, the approach of $`T_{ab}`$ to the de Sitter invariant Allen-Folacci value is perhaps more surprising. As shown in Section IV one expects the leading order contribution of the modes to $`T_{ab}`$ in this case to be constant in comoving time at late times. In fact this occurs if $`m`$ and $`\xi `$ are not separately zero for all the modes. However when $`m`$ and $`\xi `$ are both zero, the leading order contributions to $`T_{ab}`$ of all the modes except the spatially homogeneous one, for an arbitrary physically admissable state have exactly zero coefficient, the subleading contributions redshift away, and we are left only with the with the de Sitter invariant Allen-Folacci constant value at late times. The finite difference from the Bunch-Davies value may be attributed entirely to the constant behavior of the spatially homogeneous mode contributing to the vacuum energy condensate in de Sitter space.
In all those cases for which $`\mathrm{}(\nu )\frac{3}{2}`$ when $`T_{ab}`$ approaches a de Sitter invariant value at late times, the quantum expectation value loses all its initial state dependence and hence its asymptotic value must be determined purely by the background geometry. When the mass of the field vanishes, the existence of only one covariantly conserved local geometrical tensor of adiabatic order four in de Sitter space, namely $`{}_{}{}^{(3)}H_{ab}^{}`$ given by Eq. (146) below, permits us to identify the asymptotic value of $`T_{ab}`$ in the vacuum energy condensate with this tensor. Since $`{}_{}{}^{(3)}H_{ab}^{}`$ cannot be derived by variation of a covariant local action, but corresponds instead to a certain well defined non-local term in the quantum effective action , the asymptotic vacuum energy condensate of the quantum field is determined by the global or extreme infrared properties of de Sitter space. The form of the non-local effective action is determined by the trace anomaly in conformally flat spacetimes. Since the approach of $`T_{ab}`$ to a de Sitter invariant value occurs for all $`\xi 0`$, the existence of this term in the effective action for massless fields is much more general than the strict definition of the trace anomaly in the conformally invariant case. Hence the asymptotic late time behavior of $`T_{ab}`$ in de Sitter space can be used to define a generalized trace “anomaly” coefficient in the massless but non-conformally invariant cases, $`\xi \frac{1}{6}`$. As we show by consideration of the covariant $`\zeta `$ function method , this coefficient is exactly the same as that which determines the infrared response of the vacuum condensate to global Weyl rescalings. Hence the significance of the state-independence of the vacuum energy condensate in de Sitter space is that it determines certain conformal properties of non-conformal field theories in the extreme infrared.
The paper is organized as follows. In Section II we give the expectation value of the energy-momentum tensor as a mode sum for an arbitrary homogeneous and isotropic, physically admissable state with a non-zero initial particle number of the scalar field. In Section III we analyze the late time behavior of expectation value of the energy-momentum tensor in de Sitter space in flat spatial sections and show that it approaches the de Sitter invariant Bunch-Davies value for all $`\mathrm{}(\nu )<\frac{3}{2}`$, independently of the initial state. In Section IV we analyze the limit $`\mathrm{}(\nu )\frac{3}{2}`$ in closed spatial sections in order to keep careful track of the spatially homogeneous mode in a discrete basis, and show how the Allen-Folacci fixed point at late times is obtained for the massless minimally coupled field. We also investigate the asymptotic behavior of the energy-momentum tensor for arbitrary mass and curvature coupling when $`\nu \frac{3}{2}`$. In Section V we illustrate the analytic results with numerical studies, investigating in particular the interesting case when $`\nu `$ is slightly smaller than $`\frac{3}{2}`$. We find that for many states when $`\nu `$ is only slightly smaller than $`\frac{3}{2}`$ the energy-momentum tensor first approaches the Allen-Folacci value and only much later approaches the Bunch-Davies value. In Section VI we consider the geometric significance of the state independent asymptotic behavior of $`T_{ab}`$, relating it to the quantum effective action which determines the behavior of $`S_{\mathrm{e}ff}`$ under global Weyl scaling and providing the generalization of the trace “anomaly” in the infrared, for $`\xi \frac{1}{6}`$. Section VII contains some discussion and final conclusions. There are two Appendices. Appendix A completes the proof of the Bunch-Davies attractor behavior in the cases of integer and pure imaginary $`\nu `$, while Appendix B contains a discussion of the simple harmonic oscillator in the limit of vanishing frequency, which shares many features with the spatially constant mode in de Sitter spacetime.
## II Scalar field in a RW background
The metric for a general RW spacetime can be written in conformal time $`\eta `$ in the form
$$\mathrm{d}s^2=a^2(\eta )\left(\mathrm{d}\eta ^2+\frac{\mathrm{d}r^2}{1\kappa r^2}+r^2\mathrm{d}\mathrm{\Omega }^2\right).$$
(1)
Here $`a(\eta )`$ is the scale factor and $`\kappa =0,+1,1`$ corresponds to the cases of flat, spherical, and hyperbolic spatial sections, respectively. Throughout we use units such that $`\mathrm{}=c=1`$ and the Misner, Thorne, and Wheeler conventions for the curvature tensors, $`R_{bcd}^a=\mathrm{\Gamma }_{bd,c}^a\mathrm{}`$ and $`R_{ab}=R_{acb}^c`$.
We consider in this paper a free quantum scalar field $`\mathrm{\Phi }`$ with the quadratic action
$`S={\displaystyle \frac{1}{2}}{\displaystyle \mathrm{d}^4x\sqrt{g}\left[(_a\mathrm{\Phi })g^{ab}(_b\mathrm{\Phi })+m^2\mathrm{\Phi }^2+\xi R\mathrm{\Phi }^2\right]},`$ (2)
where $`_a`$ denotes the covariant derivative, $`R`$ is the scalar curvature, and $`gdet(g_{ab})`$. The mass $`m`$ and curvature coupling $`\xi `$ are allowed to have any real value. The wave equation for $`\mathrm{\Phi }`$ obtained by varying this action is
$`\left[\text{.09}.09+m^2+\xi R\right]\mathrm{\Phi }(\eta ,𝐱)=\left[{\displaystyle \frac{1}{a^4}}{\displaystyle \frac{}{\eta }}\left(a^2{\displaystyle \frac{}{\eta }}\right){\displaystyle \frac{1}{a^2}}\mathrm{\Delta }^{(3)}+m^2+\xi R\right]\mathrm{\Phi }=0,`$ (3)
with $`\mathrm{\Delta }^{(3)}`$ the covariant spatial Laplacian. For spacetimes with the metric (1) the field $`\mathrm{\Phi }`$ can be expanded as a mode sum in the form,
$$\mathrm{\Phi }(\eta ,𝐱)=\frac{1}{a(\eta )}d\stackrel{~}{\mu }(𝐤)\left[a_𝐤Y_𝐤(𝐱)\psi _k(\eta )+a_𝐤^{}Y_𝐤^{}(𝐱)\psi _k^{}(\eta )\right],$$
(4)
where the integration measure is given by
$`{\displaystyle d\stackrel{~}{\mu }(𝐤)}\{\begin{array}{ccc}\mathrm{d}^3𝐤\hfill & \mathrm{if}\hfill & \kappa =0,\hfill \\ _0^{\mathrm{}}dk_{l,m}\hfill & \mathrm{if}\hfill & \kappa =1,\hfill \\ _{k,l,m}\hfill & \mathrm{if}\hfill & \kappa =+1,\hfill \end{array}`$ (8)
and the spatial part of the mode functions $`Y_𝐤(𝐱)`$ obeys the equation
$$\mathrm{\Delta }^{(3)}Y_𝐤(𝐱)=(k^2\kappa )Y_𝐤(𝐱),$$
(9)
with $`k=1,2,\mathrm{}`$ in the case of closed spatial sections, $`\kappa =+1`$. The time dependent part of the mode functions $`\psi _k`$ obeys the equation
$$\psi _{k}^{}{}_{}{}^{\prime \prime }+\left[k^2+m^2a^2+\left(\xi \frac{1}{6}\right)a^2R\right]\psi _k=0,$$
(10)
where primes denote derivatives with respect to the conformal time variable $`\eta `$, and the scalar curvature in a general RW spacetime is given by
$$R=6\left(\frac{a^{\prime \prime }}{a^3}+\frac{\kappa }{a^2}\right).$$
(11)
For the quantum field to satisfy the canonical commutation relations, the creation and annihilation operators are required to obey the commutation relations $`[a_𝐤,a_𝐤^{}^{}]=\delta _{\mathrm{𝐤𝐤}^{}}`$, whereupon the $`\psi _k`$ must obey the Wronskian condition
$$\psi _k\psi _{k}^{}{}_{}{}^{}\psi _k^{}\psi _k^{}=i.$$
(12)
The components of the unrenormalized energy-momentum tensor (energy density and trace) are given by
$`\epsilon _u`$ $`=`$ $`T_{}^{0}{}_{0}{}^{}_u={\displaystyle \frac{1}{4\pi ^2a^4}}{\displaystyle }\mathrm{d}\mu (k)(2n_k+1)\{|\psi _k^{}|^2+(k^2+m^2a^2)|\psi _k|^2`$ (15)
$`+(6\xi 1)[{\displaystyle \frac{a^{}}{a}}(\psi _k\psi _{k}^{}{}_{}{}^{}+\psi _k^{}\psi _k^{})({\displaystyle \frac{a^2}{a^2}}\kappa )|\psi _k|^2]\},`$
and
$`\epsilon _u+3p_u`$ $`=`$ $`T_u={\displaystyle \frac{1}{2\pi ^2a^4}}{\displaystyle }\mathrm{d}\mu (k)(2n_k+1)\{m^2a^2|\psi _k|^2+(6\xi 1)[|\psi _k^{}|^2+{\displaystyle \frac{a^{}}{a}}(\psi _k\psi _{k}^{}{}_{}{}^{}+\psi _k^{}\psi _k^{})]`$ (17)
$`+(6\xi 1)[k^2+m^2a^2+({\displaystyle \frac{a^{\prime \prime }}{a}}{\displaystyle \frac{a^2}{a^2}})+(\xi {\displaystyle \frac{1}{6}})a^2R]|\psi _k|^2]\}.`$
where we have allowed an arbitrary number of particles in the initial state, $`n_k=a_𝐤^{}a_𝐤`$, and the scalar measure $`\mathrm{d}\mu (k)`$ is given by
$`{\displaystyle d\mu (k)}\{\begin{array}{ccc}_0^{\mathrm{}}dkk^2\hfill & \mathrm{if}\hfill & \kappa =0,1,\hfill \\ _1^{\mathrm{}}k^2\hfill & \mathrm{if}\hfill & \kappa =+1.\hfill \end{array}`$
As we are considering spatially homogeneous and isotropic initial states (consistent with the RW symmetry), $`n_k`$ depends only on the magnitude $`k`$ of the spatial wave vector $`𝐤`$. Expectation values of the bilinears $`a_𝐤a_𝐤`$ and $`a_𝐤^{}a_𝐤^{}`$ in a general state need not be considered since they can be removed by a time-independent Bogoliubov transformation at the initial time . Hence these initial state correlations may be parameterized instead by the initial data on the mode functions $`\psi _k`$, together with the non-negative set of $`n_k`$, with no loss of generality.
Since the expectation value of the unrenormalized energy-momentum tensor $`T_{ab}_u`$ is quartically divergent, a procedure for defining a finite, renormalized expectation value must be given. We will follow the adiabatic regularization method . In this method the renormalization counterterms are constructed using a fourth order WKB expansion for the mode functions. We denote these counterterms by $`T_{ab}_{ad}`$. They are given in Refs. and . The renormalized energy-momentum tensor is then
$$T_{ab}_{ren}=T_{ab}_uT_{ab}_{ad}.$$
(19)
This subtraction scheme is not manifestly covariant in form, since space and time are treated quite differently. However, adiabatic regularization is equivalent to a covariant point splitting procedure in which the points are split only in the spacelike hypersurface of constant $`\eta `$ , and the value of the renormalized $`T_{ab}`$ obtained by this procedure is the same as in a strictly covariant one. Hence this subtraction procedure does correspond to adjustment of counterterms to the quantum effective action, and $`T_{ab}_{ren}`$ is covariantly conserved. As discussed in detail in Ref. the adiabatic terms in all cases consist of an integral rather than a sum over $`k`$. The reason is that subtraction corresponds to purely local counterterms in the effective action, and thus must be independent of the global compactness or non-compactness of the spatial sections.
For an arbitrary homogeneous and isotropic state to be physically admissible the renormalized energy-momentum tensor defined by the adiabatic order four subtractions in (19) must be both ultraviolet and infrared finite. In a general RW spacetime, ultraviolet finiteness for a field with non-conformal coupling to the scalar curvature requires that the particular solution of the mode equation (10) in a general physical state must match the fourth order adiabatic form at large $`k`$, with the deviations from the fourth order WKB form falling faster than $`k^4`$. Likewise the initial number $`n_k`$ must fall faster than $`k^4`$ at large $`k`$, for the mode sums (integrals) to be ultraviolet convergent. This is equivalent to the requirement that the two-point function of the scalar field have the vacuum Hadamard form to sufficiently high order in the short distance expansion as the points approach one another. As long as there are two linearly independent complex oscillatory solutions to the equation (10), the Wronskian normalization condition (12) can be imposed and the state will be free of any infrared divergences. However, when $`m^2+\xi R0`$ for some low $`k`$ in de Sitter space no such complex oscillatory solutions to (10) exist and the infrared finiteness requirement on the both the energy-momentum tensor and the two-point function of the physical state becomes non-trivial, as we discuss in detail in Section IV.
A useful variation of the method of adiabatic regularization has been developed by two of us . In this method one first computes a quantity $`T_{ab}_d`$, obtained by expanding the adiabatic counterterms $`T_{ab}_{ad}`$ in inverse powers of $`k`$ and truncating at order $`k^3`$. The same renormalized energy-momentum tensor defined in Eq. (19) is separated into the sum of two finite terms by adding and subtracting the simplified form of the divergent counterterms $`T_{ab}_d`$
$`T_{ab}_{ren}`$ $`=`$ $`T_{ab}_n+T_{ab}_{an},`$ (20)
$`T_{ab}_n`$ $`=`$ $`T_{ab}_uT_{ab}_d,`$ (21)
$`T_{ab}_{an}`$ $`=`$ $`T_{ab}_dT_{ab}_{ad}.`$ (22)
The full expressions for $`T_{ab}_d`$ and $`T_{ab}_{an}`$ are given in Ref. for a general RW spacetime. The advantage of this splitting is that $`T_{ab}_n`$ and $`T_{ab}_{an}`$ are separately conserved, and moreover, $`T_{ab}_{an}`$ may be computed analytically in terms of the scale factor $`a(\eta )`$ and its derivatives . Thus the state dependence of the renormalized $`T_{ab}_{ren}`$ resides completely in $`T_{ab}_n`$, which can be computed numerically.
## III The Bunch-Davies Attractor for $`\mathrm{}(\nu )<\frac{3}{2}`$
We restrict our consideration henceforth to the particular maximally symmetric RW background of de Sitter space. The geometry of de Sitter spacetime can be described in a number of different coordinate systems. If $`\kappa =0`$ the spatial sections are flat and the scale factor is
$$a(\eta )=\frac{\alpha }{\eta },\mathrm{}<\eta <0,\kappa =0,$$
(23)
with $`\alpha `$ a real, positive constant, and $`R=12\alpha ^2`$. If $`\kappa =+1`$ then the scale factor is
$$a(\eta )=\alpha \mathrm{sec}\eta ,\frac{\pi }{2}<\eta <\frac{\pi }{2},\kappa =+1,$$
(24)
which is equivalent to $`a(\eta )=\alpha \mathrm{csc}\eta `$ with $`0<\eta <\pi `$. To simplify the notation we will generally use dimensionless units where $`\alpha =1`$ and $`R=12`$, restoring the dimensionful quantities when it is instructive to do so.
The asymptotic behavior of the energy-momentum tensor does not depend on $`\kappa `$ so the bulk of the analysis will be carried out in the flat ($`\kappa =0`$) coordinates. However, when we turn to the massless, minimally coupled limit in the next section, it will become useful to have a discrete $`k`$ basis in order to separate out the $`k=1`$ spatially homogeneous mode explicitly, since it is the most infrared sensitive. No confusion should be caused by our use of the same symbol $`\eta `$ for conformal time in both cases of flat and closed spatial sections, since we make use of only the $`\kappa =0`$ coordinates in this Section and only the $`\kappa =+1`$ coordinates in the next Section. We will not make use of the spatially open ($`\kappa =1`$) coordinates in this paper.
For the case of Eq. (23) the general solution to the mode equation can be written as <sup>*</sup><sup>*</sup>*In Ref. the arguments of the Hankel functions are given as $`k\eta `$ rather than $`k\eta `$. We have chosen to use non-negative arguments to avoid complications that result from the fact that these functions have branch cuts along the negative real axis.
$$\psi _k(\eta )=\frac{1}{2}(\pi \eta )^{\frac{1}{2}}e^{\frac{i\nu \pi }{2}}\left[c_1(k)H_\nu ^{(1)}(k\eta )+c_2(k)H_\nu ^{(2)}(k\eta )\right],$$
(25)
where the $`H_\nu ^{(1),(2)}`$ are Hankel functions and
$$\nu ^2\frac{9}{4}m^2\alpha ^212\xi \gamma ^2.$$
(26)
The latter notation is useful in the case $`\nu ^2<0`$ so that $`\nu =i\gamma `$ is purely imaginary. When $`\nu ^2>0`$ we will choose $`\nu `$ to be the positive root of (26). From Eq. (25) we see that solutions to the mode equation in de Sitter space depend on $`m`$ and $`\xi `$ only through their dependence on the parameter $`\nu `$. Note that because of the minus sign in the arguments of the Hankel functions, it is the function $`H_\nu ^{(1)}`$ that corresponds to a positive frequency mode in the large $`k`$ limit. The normalization of the mode function in (25) has been chosen so that the Wronskian condition (12) becomes simply
$$|c_1(k)|^2|c_2(k)|^2=1.$$
(27)
The Bunch-Davies (BD) state is defined by the choice, $`c_1=1`$ and $`c_2=0`$ (with $`n_k=0`$) for all $`k`$. The renormalized energy-momentum tensor in the BD state is given by
$`T_{ab}_{BD}`$ $`=`$ $`{\displaystyle \frac{g_{ab}}{64\pi ^2}}\{m^2[m^2+(\xi {\displaystyle \frac{1}{6}})R][\psi ({\displaystyle \frac{3}{2}}+\nu )+\psi ({\displaystyle \frac{3}{2}}\nu )\mathrm{log}\left({\displaystyle \frac{12m^2}{R}}\right)]`$ (28)
$``$ $`m^2(\xi {\displaystyle \frac{1}{6}})R{\displaystyle \frac{1}{18}}m^2R{\displaystyle \frac{1}{2}}(\xi {\displaystyle \frac{1}{6}})^2R^2+{\displaystyle \frac{R^2}{2160}}\},`$ (29)
where $`\psi (z)=\frac{d\mathrm{log}\mathrm{\Gamma }(z)}{dz}`$ is the digamma function.
That this finite value of $`T_{ab}_{BD}`$ coincides with the renormalized $`T_{ab}_{ren}`$ defined by the adiabatic subtraction in (19) follows from the fact that the BD state is an allowed fourth order adiabatic state. This may be checked by comparing the asymptotic expansion of the exact BD mode function, $`\frac{(\pi \eta )^{\frac{1}{2}}}{2}H_\nu ^{(1)}(k\eta )`$, for large $`k`$ with the fourth order adiabatic mode function
$$\psi _k^{(4)}(\eta )=\frac{1}{\sqrt{2W_k^{(4)}}}\mathrm{exp}\left(i^\eta W_k^{(4)}(\eta ^{})d\eta ^{}\right),$$
(30)
with $`W_k^{(4)}(\eta )`$ the fourth order adiabatic frequency. It is given explicitly in de Sitter space in flat conformal time coordinates by
$$W_k^{(4)}(\eta )=k+\frac{1}{2k\eta ^2}\left(\frac{1}{4}\nu ^2\right)\frac{1}{8k^3\eta ^4}\left(\frac{1}{4}\nu ^2\right)\left(\frac{25}{4}\nu ^2\right)+𝒪\left(\frac{1}{k^5}\right),$$
(31)
up to the required order at large $`k`$.
For the general state with $`c_20`$ to remain fourth order adiabatic, we must have for large values of $`k`$
$$c_2(k)=\frac{C(k)}{k^4},$$
(32)
for some complex function $`C(k)`$ which vanishes in the limit $`k\mathrm{}`$. This condition is necessary for an arbitrary (spatially homogeneous) state to posses a finite energy-momentum tensor after the fourth order adiabatic subtraction defined by (19). Likewise the same condition of finite $`T_{ab}`$ requires us to restrict the average number of particles $`a_𝐤^{}a_𝐤=n_k`$ by
$$n_k=\frac{N(k)}{k^4},$$
(33)
for some real function $`N(k)`$ which vanishes in the limit $`k\mathrm{}`$. The two ultraviolet conditions
$$\underset{k\mathrm{}}{lim}|C(k)|=\underset{k\mathrm{}}{lim}N(k)=0,$$
(34)
on the physically allowed states guarantee that the Green’s function for the scalar field is locally of the Hadamard form , and that the divergences of $`T_{ab}`$ match those of the fourth order adiabatic vacuum, and are removed by the adiabatic subtraction procedure. We will call any state which satisfies the conditions (32), (33), and (34), together with the Wronskian condition (27), a UV admissible physical state In fact, the requirement of fourth order adiabatic states is slightly more restrictive than UV finiteness of the energy-momentum tensor in de Sitter space, since $`C(k)`$ can go to a non-zero constant at large $`k`$ and the mode sums still converge. This is associated with the vanishing of the tensor, $`{}_{}{}^{(1)}H_{ab}^{}`$, defined in Eq. (145) in de Sitter space. If the field is non-conformally coupled, a state with $`C(k)`$ constant would lead to a divergent energy-momentum tensor if the spacetime is not exactly de Sitter. If the field is conformally coupled and the spacetime is not exactly homogeneous and isotropic, then the energy-momentum tensor would again be divergent. Thus the most general physically acceptable UV states are fourth order adiabatic states..
We shall require also that the arbitrary physical state possess a two-point function and energy-momentum tensor which are free of any infrared divergences. Because the canonical dimension of $`T_{ab}`$ is four whereas that of $`\mathrm{\Phi }(x)\mathrm{\Phi }(x^{})`$ is two, the conditions (32) and (33) which require finiteness of $`T_{ab}_{ren}`$ are more restrictive in the UV, whereas the condition of finiteness of $`\mathrm{\Phi }(x)\mathrm{\Phi }(x^{})`$ is more restrictive in the IR. These two sets of conditions will be sufficient to demonstrate that the energy-momentum tensor for any UV and IR admissible physical state approaches the BD value at late times for $`\mathrm{}(\nu )<\frac{3}{2}`$.
To understand why such a result is to be expected and outline the more detailed proof which we give below, let us observe that at late times $`\eta 0^{}`$, the general state mode function (25) behaves like
$$\psi _k(\eta )^{\frac{1}{2}\nu }a^{\nu \frac{1}{2}}.$$
(35)
Substituting this into (15) and (17) shows that to leading order at late times the contributions to the mode sums of the unrenormalized energy-momentum tensor behave like $`(\eta )^{32\nu }a^{2\nu 3}`$ for $`\nu `$ real. Since the renormalization counterterms are state independent , the state dependent terms are the same in the unrenormalized and renormalized energy-momentum tensor. One can perform all the UV renormalization in the BD state at a fixed time and collect the remaining finite state dependent terms which are unaffected by the subtraction procedure, and they all fall off at least as fast as $`(\eta )^{32\nu }`$ as $`\eta 0^{}`$ for $`\mathrm{}(\nu )<\frac{3}{2}`$.
These remaining finite state dependent terms in the energy density and trace are expressible as integrals over the wave number $`k`$ with the general form
$$I(\eta )=_0^{\mathrm{}}\frac{\mathrm{d}k}{k}R(k)S(k\eta ),$$
(36)
where $`R(k)`$ is one of the four state dependent, but time independent functions
$`|c_2(k)|^2`$ $`(1+2n_k),`$ (37)
$`\mathrm{}[c_1(k)c_2(k)]`$ $`(1+2n_k),`$ (38)
$`\mathrm{}[c_1(k)c_2(k)]`$ $`(1+2n_k),`$ (40)
$`n_k,`$
and $`S(z=k\eta )`$ is a product of the state independent Bessel functions and their derivatives. An explicit basis for the twelve products of Bessel functions $`S_i(z)`$ for $`i=1,\mathrm{},12`$ which appear in the integrals is given in Table 1. The essential point is that all the state dependent mode integrals of the form (36) are uniformly convergent for all $`\eta `$ (including $`\eta =0`$) at both their lower limit, $`k=0`$, and their upper limit $`k=\mathrm{}`$, due to the IR and UV finiteness of the state. Hence the limit of $`\eta 0^{}`$ can be taken inside the integral over $`k`$. Since, as Table 1 shows, all the $`S_i(z)`$ behave like
$$S_i(z)s_{i,0}z^{\beta _i},\mathrm{with}\beta _i>0,$$
(41)
as $`z0`$, for $`\mathrm{}(\nu )<\frac{3}{2}`$, it follows that
$$\underset{\eta 0^{}}{lim}I(\eta )=_0^{\mathrm{}}\frac{\mathrm{d}k}{k}R(k)\underset{\eta 0^{}}{lim}S(k\eta )=0,$$
(42)
and all the state dependent contributions to $`T_{ab}_{ren}`$ vanish at late times.
The validity of bringing the limit inside the integral depends on the uniform convergence of the integral at both its upper and lower limits. In the form (36) the behavior of the $`S(k\eta )`$ factor at small arguments (and the absence of any IR divergence from the $`R(k)`$ factor) clearly guarantees the uniform convergence at the lower limit. However, the change of variables $`z=k\eta `$ brings the integral (36) into the form
$$I(\eta )=_0^{\mathrm{}}\frac{\mathrm{d}z}{z}R\left(\frac{z}{\eta }\right)S(z).$$
(43)
In this form it is clear that the uniform convergence of the integral at its upper limit is guaranteed by the falloff of the state dependent mode functions at large $`k`$, namely
$$\underset{k\mathrm{}}{lim}R(k)=\underset{\eta 0^{}}{lim}R\left(\frac{z}{\eta }\right)=0.$$
(44)
Equations (41) and (44) guarantee that both the IR and UV contributions go to zero as $`\eta 0^{}`$ for $`\mathrm{}(\nu )<\frac{3}{2}`$.
In fact, we can go one step further by using the information from the fourth order adiabatic nature of the state
$$\underset{k\mathrm{}}{lim}k^4R(k)=0,$$
(45)
to conclude that the UV contribution to the integral from $`z1`$, and very large $`k(\eta )^1`$ in (43) falls faster than $`(\eta )^4`$ as $`\eta 0^{}`$. This is faster than the IR contribution which falls only as $`(\eta )^{\beta _i}`$ if $`\beta _i4`$. If $`\beta _i>4`$ then both the IR and UV contributions fall faster than $`(\eta )^4`$. Hence, the $`S_i(z)`$ for which $`\beta _i>4`$ are subdominant at late times. Thus we conclude that the leading order state dependent terms at late times are those with the smallest $`\beta _i`$. These give an IR dominant, (i.e. finite $`k`$, $`z1`$) contribution of the form,
$$I(\eta )(\eta )^{\beta _i}_0^{\mathrm{}}dkk^{\beta _i1}R(k)(\eta )^{32\nu }_0^{\mathrm{}}dkk^{22\nu }R(k),$$
(46)
as $`\eta 0^{}`$. This last integral is guaranteed to converge at its upper limit for all $`\mathrm{}(\nu )\frac{1}{2}`$, and in particular for $`0\mathrm{}(\nu )\frac{3}{2}`$ by (45). Hence for all $`\mathrm{}(\nu )<\frac{3}{2}`$ the state dependent contributions go to zero as $`a^{2\nu 3}`$ at late times. The limiting case when $`\nu \frac{3}{2}`$ is special because then the state dependent IR contributions apparently does not go to zero at late times. This case will be considered separately in the next Section.
We will now make the proof more explicit by giving the form of all of the state dependent terms of the scalar field in de Sitter space, and analyzing the IR and UV contributions in detail. If we make use of Eq. (27), we find that in an arbitrary state
$`T_{ab}_{ren}`$ $`=`$ $`T_{ab}_{BD}+T_{ab}_{SD},`$ (47)
where $`T_{ab}_{SD}`$ is the finite state dependent term, depending on the coefficients $`c_1(k),c_2(k)`$, and $`n_k`$, which may be expressed as an integral over the wave number $`k`$ in the form
$`T_{ab}_{SD}`$ $`=`$ $`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle _0^{\mathrm{}}}dkI_{ab}(k,\eta ).`$ (48)
The explicit expressions for the integrand $`I_{ab}`$ depend on whether $`\nu `$ is real and not an integer, real and an integer, or pure imaginary, although the result will be the same for all $`\mathrm{}(\nu )<\frac{3}{2}`$. In the rest of this Section we restrict our discussion to the case where $`\nu `$ is real and not equal to an integer. The cases of integer and imaginary values of $`\nu `$ are covered in Appendix A. For real $`\nu `$ after some regrouping of terms we have
$`I_{}^{0}{}_{0}{}^{}`$ $`=`$ $`A_1(k)\left[S_1+{\displaystyle \frac{1}{4}}(9+4m^248\xi )S_4{\displaystyle \frac{3}{2}}(4\xi 1)S_7+S_{10}\right]`$ (52)
$`+A_2(k)\left[S_2+{\displaystyle \frac{1}{4}}(9+4m^248\xi )S_5{\displaystyle \frac{3}{2}}(4\xi 1)S_8+S_{11}\right]`$
$`+A_3(k)\left[S_3+{\displaystyle \frac{1}{4}}(9+4m^248\xi )S_6{\displaystyle \frac{3}{2}}(4\xi 1)S_9+S_{12}\right],\mathrm{and}`$
$`I`$ $`=`$ $`A_1(k)\left[2(6\xi 1)S_1{\displaystyle \frac{1}{2}}\left(9+8(3\xi 1)m^2102\xi +288\xi ^2\right)S_4+3(6\xi 1)S_7+2(6\xi 1)S_{10}\right]`$ (55)
$`+A_2(k)\left[2(6\xi 1)S_2{\displaystyle \frac{1}{2}}\left(9+8(3\xi 1)m^2102\xi +288\xi ^2\right)S_5+3(6\xi 1)S_8+2(6\xi 1)S_{11}\right]`$
$`+A_3(k)\left[2(6\xi 1)S_3{\displaystyle \frac{1}{2}}\left(9+8(3\xi 1)m^2102\xi +288\xi ^2\right)S_6+3(6\xi 1)S_9+2(6\xi 1)S_{12}\right],`$
with
$`A_1(k)`$ $`=`$ $`{\displaystyle \frac{\pi }{2k}}[\mathrm{csc}^2(\nu \pi )((1+2n_k)|c_2|^2+n_k)+{\displaystyle \frac{1}{2}}(1\mathrm{cot}^2(\nu \pi ))(1+2n_k)(c_1c_2^{}+c_1^{}c_2)`$ (58)
$`+{\displaystyle \frac{i}{2}}\mathrm{cot}(\nu \pi )(1+2n_k)(c_1c_2^{}c_1^{}c_2)],`$
$`A_2(k)`$ $`=`$ $`{\displaystyle \frac{\pi }{k}}[\mathrm{cot}(\nu \pi )\mathrm{csc}(\nu \pi )((1+2n_k)|c_2|^2+n_k){\displaystyle \frac{1}{2}}\mathrm{cot}(\nu \pi )\mathrm{csc}(\nu \pi )(1+2n_k)(c_1c_2^{}+c_1^{}c_2)`$ (60)
$`+{\displaystyle \frac{i}{2}}\mathrm{csc}(\nu \pi )(1+2n_k)(c_1c_2^{}c_1^{}c_2)],`$
$`A_3(k)`$ $`=`$ $`{\displaystyle \frac{\pi }{2k}}\left[\mathrm{csc}^2(\nu \pi )\left((1+2n_k)|c_2|^2+n_k\right){\displaystyle \frac{1}{2}}\mathrm{csc}^2(\nu \pi )(1+2n_k)(c_1c_2^{}+c_1^{}c_2)\right],`$ (61)
and the $`S_i(k\eta )`$ composed of various products of $`J_\nu (k\eta )`$, $`J_\nu (k\eta )`$ and their derivatives. Making use of the general formula for the product of two Bessel functions ,
$$J_\mu (z)J_\nu (z)=\underset{p=0}{\overset{\mathrm{}}{}}\frac{(1)^p\left(\frac{z}{2}\right)^{\nu +\mu +2p}\mathrm{\Gamma }(\nu +\mu +2p+1)}{p!\mathrm{\Gamma }(\nu +\mu +p+1)\mathrm{\Gamma }(\nu +p+1)\mathrm{\Gamma }(\mu +p+1)},$$
(62)
we can expand the $`S_i`$ in power series of the form,
$$S_i(z)=\underset{p=0}{\overset{\mathrm{}}{}}s_{i,p}z^{2p+\beta _i},$$
(63)
with $`z=k\eta `$. The explicit expressions for $`\beta _i`$ and $`S_i`$ for the case of real, non-integer $`\nu `$ are given in Table 1.
| $`i`$ | $`\beta _i`$ | $`S_i(z)`$ |
| --- | --- | --- |
| 1 | $`5+2\nu `$ | $`z^5J_\nu ^2(z)`$ |
| 2 | $`5`$ | $`z^5J_\nu (z)J_\nu (z)`$ |
| 3 | $`52\nu `$ | $`z^5J_\nu ^2(z)`$ |
| 4 | $`3+2\nu `$ | $`z^3J_\nu ^2(z)`$ |
| 5 | $`3`$ | $`z^3J_\nu (z)J_\nu (z)`$ |
| 6 | $`32\nu `$ | $`z^3J_\nu ^2(z)`$ |
| 7 | $`3+2\nu `$ | $`z^4\frac{\mathrm{d}}{\mathrm{d}z}J_\nu ^2(z)`$ |
| 8 | $`3`$ | $`z^4\frac{\mathrm{d}}{\mathrm{d}z}\left(J_\nu (z)J_\nu (z)\right)`$ |
| 9 | $`32\nu `$ | $`z^4\frac{\mathrm{d}}{\mathrm{d}z}J_\nu ^2(z)`$ |
| 10 | $`3+2\nu `$ | $`z^5\left(\frac{\mathrm{d}}{\mathrm{d}z}J_\nu (z)\right)^2`$ |
| 11 | $`3`$ | $`z^5\left(\frac{\mathrm{d}}{\mathrm{d}z}J_\nu (z)\right)\left(\frac{\mathrm{d}}{\mathrm{d}z}J_\nu (z)\right)`$ |
| 12 | $`32\nu `$ | $`z^5\left(\frac{\mathrm{d}}{\mathrm{d}z}J_\nu (z)\right)^2`$ |
Table 1
We are interested in the behavior of the finite state dependent terms, $`T_{ab}_{SD}`$, in the limit $`\eta 0^{}`$. To investigate this limit in detail it is useful to break up the integral in Eq. (48) into the three parts
$$T_{ab}_{SD}=\frac{1}{4\pi ^2}_0^\lambda dkI_{ab}(k,\eta )+\frac{1}{4\pi ^2}_\lambda ^{Z/\eta }dkI_{ab}(k,\eta )+\frac{1}{4\pi ^2}_{Z/\eta }^{\mathrm{}}dkI_{ab}(k,\eta ).$$
(64)
Here $`\lambda `$ is a finite positive constant. For $`k>\lambda `$ we can make use of the UV conditions (32), (33), and (34) in the second and third integrals. Thus, the most infrared sensitive integral is the first one. The positive constant $`Z`$ is arbitrary, provided only that $`Z>\lambda \eta `$, which is always satisfied for fixed $`\lambda `$ and small enough $`\eta `$. Hence the second integral provides the bulk of the contribution of the state dependent wave numbers that have redshifted outside the de Sitter horizon at late times. The third integral is the contribution of the state dependent terms still within the horizon which go to zero very rapidly at late times due to the UV conditions. If the $`\eta 0^{}`$ limit is uniform then all three integrals should vanish unconditionally in this limit, i.e. without any restrictions on the arbitrary parameters $`\lambda `$ and $`Z`$.
We begin by analyzing the first integral in (64). We are considering only fourth order adiabatic states that are IR admissible. Hence the $`k`$ integration converges at its lower limit and we can expand the integrand for this first integral in a series of the form (63) and interchange the order of summation and integration. Each term in the resulting sums contains an integral over $`k`$ that is finite by assumption and a factor of $`(\eta )^{2p+\beta _i}`$ where $`p=0,1,2\mathrm{}`$ is a non-negative integer. In fact, since at late times, $`\lambda \eta 1`$ for any finite $`\lambda `$, it is sufficient to consider only the leading $`p=0`$ term. From Table 1 it is clear that for $`\nu <\frac{3}{2}`$, $`\beta _i>0`$ for all $`i`$. Hence in the limit $`\eta 0^{}`$, the integrand vanishes and therefore, the first integral on the right hand side of Eq. (64) vanishes in this late time limit.
For the second integral in (64) we may utilize the expansion (62) again. Since the integration is between finite limits for finite $`\eta `$ we can exchange the order of summation over $`p`$ and integration over $`k`$. Having done so we can then bound each term in the sums by taking the absolute value of its factors. Inspection of the expressions (55) and (61) indicates that the result is a linear combination of terms involving integrals of the three possible forms
$`_1`$ $`=`$ $`{\displaystyle _\lambda ^{Z/\eta }}{\displaystyle \frac{\mathrm{d}k}{k}}|c_2(k)|^2(1+2n_k)(k\eta )^{2p+\beta _i},`$ (65)
$`_2`$ $`=`$ $`{\displaystyle _\lambda ^{Z/\eta }}{\displaystyle \frac{\mathrm{d}k}{k}}|c_1(k)c_2(k)|(1+2n_k)(k\eta )^{2p+\beta _i},`$ (66)
$`_3`$ $`=`$ $`{\displaystyle _\lambda ^{Z/\eta }}{\displaystyle \frac{\mathrm{d}k}{k}}n_k(k\eta )^{2p+\beta _i},`$ (67)
multiplied by constant coefficients. Because of (34) for any $`\lambda `$
$`|C(k)|`$ $`<`$ $`C,`$ (68)
$`N(k)`$ $`<`$ $`N,`$ (69)
for some real positive numbers $`C`$ and $`N`$, depending on $`\lambda `$. With these bounds we can bound the values of $`|c_2(k)|`$ and $`n_k`$, and use Eq. (27) to bound $`|c_1(k)|`$ for all $`k\lambda `$ as follows
$`|c_2(k)|`$ $`<`$ $`{\displaystyle \frac{C}{k^4}},`$ (70)
$`1|c_1(k)|`$ $`=`$ $`\left[1+|c_2(k)|^2\right]^{\frac{1}{2}}1+|c_2(k)|^2<1+{\displaystyle \frac{C^2}{k^8}}1+{\displaystyle \frac{C^2}{\lambda ^8}},`$ (71)
$`n_k`$ $`<`$ $`{\displaystyle \frac{N}{k^4}}`$ (72)
$`1+2n_k`$ $`<`$ $`1+{\displaystyle \frac{2N}{k^4}}1+{\displaystyle \frac{2N}{\lambda ^4}}.`$ (73)
Using these bounds the integrals in Eq. (67) can be bounded as follows
$`_1`$ $`<`$ $`C^2(1+2N/\lambda ^4){\displaystyle _\lambda ^{Z/\eta }}dkk^{2p+\beta _i9}(\eta )^{2p+\beta _i}`$ (74)
$`=`$ $`{\displaystyle \frac{C^2(1+2N/\lambda ^4)}{\beta _i+2p8}}\left[(\eta )^8Z^{2p+\beta _i8}(\eta )^{2p+\beta _i}\lambda ^{2p+\beta _i8}\right]`$ (75)
$`_2`$ $`<`$ $`C(1+C^2/\lambda ^8)(1+2N/\lambda ^4){\displaystyle _\lambda ^{Z/\eta }}dkk^{2p+\beta _i5}(\eta )^{2p+\beta _i}`$ (76)
$`=`$ $`{\displaystyle \frac{C(1+C^2/\lambda ^8)(1+2N/\lambda ^4)}{\beta _i+2p4}}\left[(\eta )^4Z^{2p+\beta _i4}(\eta )^{2p+\beta _i}\lambda ^{2p+\beta _i4}\right]`$ (77)
$`_3`$ $`<`$ $`N{\displaystyle _\lambda ^{Z/\eta }}dkk^{2p+\beta _i5}(\eta )^{2p+\beta _i}`$ (78)
$`=`$ $`{\displaystyle \frac{N}{\beta _i+2p4}}\left[(\eta )^4Z^{2p+\beta _i4}\lambda ^{2p+(\eta )^{2p+\beta _i}\beta _i4}\right].`$ (79)
Each of these bounds vanishes in the limit $`\eta 0^{}`$. Therefore each of the terms appearing in the second integral in (64) vanishes in the late time limit for $`\mathrm{}(\nu )<\frac{3}{2}`$.
If $`\nu =\frac{1}{2}`$ some of the terms will have vanishing denominators and should be interpreted according to the limiting relation
$$\underset{q0}{lim}\frac{\left(\frac{Z}{\eta }\right)^q\lambda ^q}{q}=\mathrm{log}\left(\frac{Z}{\eta \lambda }\right),$$
(80)
but the appearance of these logarithms does not change the result since they are always multiplied by at least $`(\eta )^{2p+\beta _i}`$ which vanishes for $`\mathrm{}(\nu )<\frac{3}{2}`$ for $`p=0,1,2\mathrm{}`$. The cases $`\nu =0,1`$ also involve logarithms in the Bessel function expansions but the result that the second integral in (64) vanishes in the late time limit is unchanged. One may consider also the case when $`\nu =i\gamma `$ is pure imaginary, where the forms of the coefficients (61) and bilinears $`S_i(z)`$ in the table change somewhat, with again the same result. For completeness these cases are treated in detail in Appendix A.
Finally, for the third integral on the right hand side of Eq. (64) we do not expand the Bessel functions in powers of $`(k\eta )`$. Instead we change the integration variable to $`z=k\eta `$ and find integrals of the forms
$`𝒥_1`$ $`=`$ $`(\eta )^8{\displaystyle _Z^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}z}{z^9}}\left|C\left({\displaystyle \frac{z}{\eta }}\right)\right|^2\left[1+2\left({\displaystyle \frac{\eta }{z}}\right)^4N\left({\displaystyle \frac{z}{\eta }}\right)\right]S_i(z),`$ (81)
$`𝒥_2`$ $`=`$ $`(\eta )^4{\displaystyle _Z^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}z}{z^5}}(\mathrm{}\mathrm{or}\mathrm{})\left\{c_1\left({\displaystyle \frac{z}{\eta }}\right)\left[C^{}\left({\displaystyle \frac{z}{\eta }}\right)\right]\right\}\left[1+2\left({\displaystyle \frac{\eta }{z}}\right)^4N\left({\displaystyle \frac{z}{\eta }}\right)\right]S_i(z),`$ (82)
$`𝒥_3`$ $`=`$ $`(\eta )^4{\displaystyle _Z^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}z}{z^5}}N\left({\displaystyle \frac{z}{\eta }}\right)S_i(z).`$ (83)
Since we are considering fourth order adiabatic states these integrals are finite for all $`\mathrm{}<\eta 0`$. Furthermore, the integrands of these integrals are finite throughout the entire integration range for $`\mathrm{}<\eta 0`$. With the lower limit $`Z>0`$ fixed, it is clear that we can evaluate the integrals at $`\eta =0`$ by first taking the limit $`\eta 0^{}`$ of the integrands and then computing the integrals. Since the integrands all vanish in the limit $`\eta 0^{}`$, the integrals do as well. Then the third integral in (64) vanishes in the late time limit for $`\mathrm{}(\nu )<\frac{3}{2}`$. From the adiabatic four conditions (34) these UV contributions go to zero faster than $`(\eta )^4`$ for any fourth order adiabatic state.
Therefore we have proven that for $`\mathrm{}(\nu )<\frac{3}{2}`$ all the integrals in Eq. (64) vanish in the limit $`\eta 0^{}`$ for an arbitrary state that is both infrared and ultraviolet finite. The fact that all three integrals vanish at late times, independently of the parameters $`\lambda `$ and $`Z`$ which we introduced to control the IR and UV contributions of the mode integral verifies that the total mode integral does converge uniformly in $`\eta `$, as expected. Hence all the state dependent contributions vanish at late times, $`\eta 0^{}`$, and we have proven that in this limit $`T_{ab}_{ren}T_{ab}_{BD}`$. Thus, we conclude that for $`\mathrm{}(\nu )<\frac{3}{2}`$ the value of the energy-momentum tensor in any physically admissable homogeneous and isotropic RW quantum state asymptotically approaches the Bunch-Davies value in de Sitter space at late times.
Moreover, from inspection of (79) we observe that the contribution from the upper limit at $`Z`$ in the second integral always falls faster than $`(\eta )^4`$, while the contribution from the lower limit at $`\lambda `$ falls off only as $`(\eta )^{\beta _i}`$ as $`\eta 0^{}`$. This must be the same order as the first IR integral so that the arbitrary parameter $`\lambda `$ drops out of the final result. Hence our detailed evaluation has verified that the leading order state dependent corrections to the BD expectation value at late times come from the terms in $`T_{ab}_{ren}`$ with the smallest $`\beta _i`$. From Table 1 these are the $`i=6,9`$, and $`12`$ terms. Collecting these terms from (55) and (61), and the numerical coefficients from the Bessel function product formula (62), we conclude that the leading order behavior of the energy density and trace at late times are given by
$`\epsilon `$ $``$ $`a^{2\nu 3}\left[\left(\nu {\displaystyle \frac{1}{2}}\right)^2+m^2+2(6\xi 1)(\nu 1)\right]{\displaystyle _0^{\mathrm{}}}dkk^{22\nu }R(k),\mathrm{and}`$ (85)
$`T`$ $``$ $`2a^{2\nu 3}\left\{m^2+(6\xi 1)\left[\left(\nu {\displaystyle \frac{1}{2}}\right)^2+2\nu +m^2+2(6\xi 1)\right]\right\}{\displaystyle _0^{\mathrm{}}}dkk^{22\nu }R(k),`$ (86)
respectively, with
$$R(k)=\frac{2^{2\nu 3}\mathrm{csc}^2(\pi \nu )}{\pi \left[\mathrm{\Gamma }(\nu +1)\right]^2}\left\{(1+2n_k)\left[|c_2(k)|^2\mathrm{}(c_1c_2^{})\right]+n_k\right\}.$$
(87)
The leading order state dependent contribution at late times is $`a^{2\nu 3}`$ from the IR part of the mode integral, which falls off very slowly if $`\mathrm{}(\nu )\frac{3}{2}`$. We examine this latter limit in detail in the next Section.
## IV The Allen-Folacci Attractor for $`\mathrm{}(\nu )\frac{3}{2}`$
In the analysis of the previous Section we saw that the terms with the slowest falloff at late times were those with the smallest $`\beta _i=32\nu `$ which behave like $`a^{2\nu 3}`$ for $`\nu `$ real and positive, and the coefficient of this falloff is controlled by the finite $`k`$, IR part of the mode integral. To examine the limit $`\nu \frac{3}{2}`$ carefully, it is easiest to work with closed spatial sections and a discrete set of mode functions in order to treat the most infrared sensitive, spatially homogeneous $`k=1`$ mode separately from the rest, instead of dealing with an infrared sensitive continuous mode integral.
The scale factor for $`\kappa =+1`$ is given by Eq. (24) and under the variable substitution $`\zeta =i\mathrm{tan}\eta `$ the mode equation (10) becomes Legendre’s differential equation. Hence the general solutions may be expressed in terms of associated Legendre functions $`P_{\frac{1}{2}+\nu }^{\pm k}(\zeta )`$. Since as conventionally defined these functions have a cut discontinuity on the real axis from $`1`$ to $`1`$ if $`k`$ is an odd integer, we write the fundamental complex valued solution for real $`\nu `$ in the form
$`f_k(\eta )`$ $`=`$ $`\left[{\displaystyle \frac{\mathrm{\Gamma }\left(k+\frac{1}{2}+\nu \right)\mathrm{\Gamma }\left(k+\frac{1}{2}\nu \right)}{2}}\right]^{\frac{1}{2}}\mathrm{exp}\left({\displaystyle \frac{ik\pi }{2}}ϵ(\eta )\right)P_{\frac{1}{2}+\nu }^k(i\mathrm{tan}\eta )`$ (88)
$`=`$ $`\left[{\displaystyle \frac{\mathrm{\Gamma }\left(k+\frac{1}{2}+\nu \right)\mathrm{\Gamma }\left(k+\frac{1}{2}\nu \right)}{2}}\right]^{\frac{1}{2}}{\displaystyle \frac{e^{ik\eta }}{k!}}F({\displaystyle \frac{1}{2}}+\nu ,{\displaystyle \frac{1}{2}}\nu ;k+1;{\displaystyle \frac{1i\mathrm{tan}\eta }{2}}),`$ (89)
where $`ϵ(\eta )=\theta (\eta )\theta (\eta )=\pm 1`$ is the sign function and $`F`$ is the hypergeometric function. The phase factor depending on the sign of $`\eta `$ for odd $`k`$ removes the discontinuity in the $`P_{\frac{1}{2}+\nu }^k`$ function as $`\eta `$ approaches zero from positive or negative values, as the second form of (89) makes clear, since the $`F`$ function is an analytic function of $`\frac{1\zeta }{2}`$, with no branch cuts for $`\zeta `$ on the imaginary axis. With the normalization factors chosen as in (89) the Wronskian condition
$$f_kf_{k}^{}{}_{}{}^{}f_k^{}f_k^{}=i,$$
(90)
is satisfied. Hence the general solution of the mode equation (10) is
$$\psi _k(\eta )=\alpha _kf_k(\eta )+\beta _kf_k^{}(\eta ),$$
(91)
with
$$|\alpha _k|^2|\beta _k|^2=1.$$
(92)
The Bunch-Davies state is given by $`\alpha _k=1`$ and $`\beta _k=0`$.
Now as $`\nu \frac{3}{2}`$ inspection of (89) shows that all the $`f_k`$ for $`k>1`$ are regular. In fact, in that limit the hypergeometric series for $`F`$ terminates and the $`f_k(\eta )`$ become the elementary functions
$$\underset{\nu \frac{3}{2}}{lim}f_k(\eta )=\frac{e^{ik\eta }}{[2k(k^21)]^{\frac{1}{2}}}(k+i\mathrm{tan}\eta ),k=2,3,\mathrm{},$$
(93)
so they can be treated in the Bunch-Davies state $`\alpha _k=1,\beta _k=0`$ with no difficulty. However in this limit the $`k=1`$ mode function is singular and must be treated separately. The behavior of the $`k=1`$ mode as $`\nu \frac{3}{2}`$ is similar to that of a simple harmonic oscillator mode as its frequency goes to zero, i.e. in the limit where the harmonic oscillator becomes a free particle. The zero frequency limit in this simple flat space analogy is reviewed in Appendix B. Just as in that case, one can construct regular solutions to the mode equation in the limit $`\nu \frac{3}{2}`$ by taking suitable linear combinations of $`f_1`$ and $`f_1^{}`$. In fact, the limiting form of $`f_1(\eta )`$ can be found from Sec. 2.3.1 of Ref. , which gives
$$f_1(\eta )\frac{\mathrm{sec}\eta }{2\sqrt{\frac{3}{2}\nu }}\frac{i}{2}\sqrt{\frac{3}{2}\nu }\mathrm{sec}\eta (\eta +\mathrm{sin}\eta \mathrm{cos}\eta )+\mathrm{}.$$
(94)
We have neglected terms in the real part of $`f_1`$ that are of order $`\sqrt{\frac{3}{2}\nu }`$. We have also neglected all terms that go to zero faster than this as $`\nu \frac{3}{2}`$. Extracting the scale factor $`a(\eta )=\mathrm{sec}\eta `$ we now define the real functions $`u`$ and $`v`$ by
$`u(\eta )`$ $``$ $`{\displaystyle \frac{1}{a(\eta )}}\underset{\nu \frac{3}{2}}{lim}\left\{\sqrt{{\displaystyle \frac{3}{2}}\nu }(f+f^{})\right\}=1,`$ (95)
$`v(\eta )`$ $``$ $`{\displaystyle \frac{i}{a(\eta )}}\underset{\nu \frac{3}{2}}{lim}\left\{{\displaystyle \frac{1}{2\sqrt{\frac{3}{2}\nu }}}(ff^{})\right\}={\displaystyle \frac{\eta +\mathrm{sin}\eta \mathrm{cos}\eta }{2}}.`$ (96)
We next define new coefficients
$`A`$ $``$ $`i\underset{\nu \frac{3}{2}}{lim}\left\{\sqrt{{\displaystyle \frac{3}{2}}\nu }(\alpha _1\beta _1)\right\},`$ (97)
$`B`$ $``$ $`\underset{\nu \frac{3}{2}}{lim}\left\{{\displaystyle \frac{1}{2\sqrt{\frac{3}{2}\nu }}}(\alpha _1+\beta _1)\right\},`$ (98)
which are finite in this limit. With these definitions the normalization is such that
$$A^{}BB^{}A=i,$$
(99)
and the limit of the general $`k=1`$ mode function may be written
$$\underset{\nu \frac{3}{2}}{lim}\psi _1(\eta )=\mathrm{sec}\eta (Av+Bu)=a(\eta )\left[\frac{A}{2}(\eta +\mathrm{sin}\eta \mathrm{cos}\eta )+B\right].$$
(100)
One can also define the time-independent Hermitian operators,
$`Q`$ $``$ $`{\displaystyle \frac{1}{2\sqrt{\frac{3}{2}\nu }}}\left[(\alpha _1+\beta _1)a_1+(\alpha _1+\beta _1)^{}a_1^{}\right]Ba_1+B^{}a_1^{},`$ (101)
$`P`$ $``$ $`i\sqrt{{\displaystyle \frac{3}{2}}\nu }\left[(\alpha _1\beta _1)a_1(\alpha _1\beta _1)^{}a_1^{}\right]Aa_1+A^{}a_1^{},`$ (102)
obeying the canonical commutation relations, $`[Q,P]=i`$. They also remain finite in the limit $`\nu \frac{3}{2}`$.
We will consider the limit $`\nu \frac{3}{2}`$ of the energy density in the particular order of fixing the mass of the scalar field to be $`m=0`$ and letting the dimensionless parameter $`\xi `$ approach zero, since this case is relevant for the infrared scaling analysis of Section VI. From (26) with $`m=0`$ and $`\xi `$ small we have
$$\xi \frac{(32\nu )}{8}0.$$
(103)
The complementary case, $`\xi =0`$ and $`m^20`$ is similar and has been discussed in Refs. and . From (15) we read the energy density in the $`k=1`$ mode for $`m=0`$
$$\epsilon _1=T_{}^{0}{}_{0}{}^{}_{k=1}=\frac{(1+2n_1)}{4\pi ^2a^4}\left\{|\psi _1^{}|^2+|\psi _1|^2+(6\xi 1)[\mathrm{tan}\eta (\psi _1\psi _1^{}+\psi _1^{}\psi _1^{})(\mathrm{tan}^2\eta 1)|\psi _1|^2]\right\}.$$
(104)
We are interested first in the asymptotic form of this energy density at late times, $`\eta \frac{\pi }{2}`$, and then in the limiting form of the resulting expression as $`\xi 0`$ according to (103). The asymptotic late time limit of the mode function $`f_1(\eta )`$ for any $`\nu >0`$ can be found from the inversion transformation of the hypergeometric function, given by formula 2.1.4 (17) of . We find that
$`f_1\left(\eta {\displaystyle \frac{\pi }{2}}\right)`$ $``$ $`\left[{\displaystyle \frac{\mathrm{\Gamma }\left(\frac{3}{2}\nu \right)}{2\mathrm{\Gamma }\left(\frac{3}{2}+\nu \right)}}\right]^{\frac{1}{2}}{\displaystyle \frac{\mathrm{\Gamma }(2\nu )}{\mathrm{\Gamma }(\frac{1}{2}+\nu )}}(i)\left({\displaystyle \frac{i\mathrm{sec}\eta }{2}}\right)^{\nu \frac{1}{2}}a^{\nu \frac{1}{2}},`$ (105)
$`f_1^{}`$ $``$ $`\left(\nu {\displaystyle \frac{1}{2}}\right)\mathrm{sec}\eta f_1a^{\nu +\frac{1}{2}},`$ (106)
as $`\eta \frac{\pi }{2}`$, which is the same late time behavior in terms of the scale factor that we found in the spatially flat coordinates. Notice that at $`\nu =\frac{3}{2}`$ the phase factor cancels and these limiting forms of the oscillatory mode functions become real, while the $`\mathrm{\Gamma }`$ function has a pole singularity there. Thus, keeping only the leading behavior as $`\nu `$ approaches $`\frac{3}{2}`$ we find that the late time limits of the mode function and its derivative are
$`\psi _1\left(\eta {\displaystyle \frac{\pi }{2}}\right)`$ $``$ $`{\displaystyle \frac{\alpha _1+\beta _1}{2\sqrt{\frac{3}{2}\nu }}}a^{\nu \frac{1}{2}},`$ (107)
$`\psi _1^{}\left(\eta {\displaystyle \frac{\pi }{2}}\right)`$ $``$ $`{\displaystyle \frac{\alpha _1+\beta _1}{2\sqrt{\frac{3}{2}\nu }}}a^{\nu +\frac{1}{2}}.`$ (108)
Thus, the dominant terms in the $`k=1`$ energy density $`\epsilon _1`$ in this limit are $`|\psi _1^{}|^2`$ and the terms involving $`\frac{a^{}}{a}=\mathrm{tan}\eta `$, and we obtain
$`\underset{\eta \frac{\pi }{2}}{lim}\epsilon _1=\left[{\displaystyle \frac{3(1+2n_1)}{32\pi ^2}}|\alpha _1+\beta _1|^2+𝒪(32\nu )\right]a^{2\nu 3}.`$ (109)
The singularity in the $`\mathrm{\Gamma }`$ function has canceled against the $`\xi `$ in the numerator. The remaining coefficient is clearly state dependent. From this asymptotic form of the energy density in the $`k=1`$ mode at late times it is clear that for $`\nu `$ close to but still less than $`\frac{3}{2}`$, the state-dependent energy density $`\epsilon _1`$ goes to zero at late times, albeit very slowly, which is consistent with our previous flat section analysis.
The corresponding expression for the trace in the $`k=1`$ mode is
$$\epsilon _1+3p_1=\frac{(1+2n_1)}{2\pi ^2a^4}(6\xi 1)\left\{|\psi _1^{}|^2+[\mathrm{tan}\eta (\psi _1\psi _1^{}+\psi _1^{}\psi _1^{})+[\mathrm{sec}^2\eta +2(6\xi 1)\mathrm{sec}^2\eta ]|\psi _1|^2]\right\}.$$
(110)
Substituting the late time asymptotic forms (108) to leading order in $`\xi `$ as before yields
$$\epsilon _1+3p_1\left[\frac{3(1+2n_1)}{8\pi ^2}|\alpha _1+\beta _1|^2+𝒪(32\nu )\right]a^{2\nu 3}4\epsilon _1.$$
(111)
Hence, $`p_1\epsilon _1`$ and the contribution from this mode is de Sitter invariant at late times for any initial physical state.
Since the Bunch-Davies state is $`\alpha _k=1`$ and $`\beta _k=n_k=0`$ for all $`k`$ and this state has a renormalized $`T_{ab}`$ which is strictly time independent for all $`\nu `$, the time dependent contribution in (109) and (111) for $`\nu <\frac{3}{2}`$ in the $`k=1`$ mode must be canceled by a time dependent contribution from all the other modes in the renormalized energy-momentum tensor. In other words, the late time behavior of all the $`k>1`$ modes in the BD state must be
$$\epsilon _{BD}|_{k>1}\epsilon _{BD}\left[\frac{3}{32\pi ^2}+𝒪(32\nu )\right]a^{2\nu 3},$$
(112)
for $`\nu `$ close to $`\frac{3}{2}`$. The pressure for the $`k>1`$ modes is obtained from this by the de Sitter invariant relation $`p=\epsilon `$.
The subtraction of the second term in (112) can be understood in a different way. Consider the short distance expansion of the BD two-point function for $`\nu <\frac{3}{2}`$, namely
$$G_{BD}(x,x^{})\frac{1}{8\pi ^2}\left[\frac{1}{1Z}\mathrm{log}(1Z)+\frac{1}{\frac{3}{2}\nu }\right],$$
(113)
for the de Sitter invariant bi-scalar, $`Z(x,x^{})1`$ and $`\nu `$ close to $`\frac{3}{2}`$. The constant term is singular at $`\nu =\frac{3}{2}`$, but it gives a finite contribution to the energy density from the $`\xi G_{ab}\varphi ^2`$ term in the energy-momentum tensor. In the dimensionless units we are using this is equal to
$$\frac{3}{8\pi ^2}\xi \left(\frac{R}{12}\right)^2\frac{1}{\frac{3}{2}\nu }=\frac{3}{32\pi ^2},$$
(114)
for $`\xi 0`$ according to (103). The last constant term in (113) is absent in the short distance expansion of the Allen-Folacci two-point function , and in the Bunch-Davies two-point function it comes entirely from the $`k=1`$ mode in the BD state. Hence the contribution to the energy density of the $`k>1`$ modes in the BD state does not contain (114), and must equal the full $`\epsilon _{BD}`$ minus (114) from the $`k=1`$ mode, which is equivalent to (112) at $`\nu =\frac{3}{2}`$.
With the $`k=1`$ and $`k>1`$ mode contributions separated we are now in a position to take the limit of $`\nu \frac{3}{2}`$. Apparently we would obtain a state dependent contribution from (109). However, inspection of (100) shows that requiring the $`k=1`$ mode function to be regular as $`\nu \frac{3}{2}`$ is equivalent to requiring that $`A`$ and $`B`$ remain finite in this limit. From (98) we then have
$$\alpha _1+\beta _12B\sqrt{\frac{3}{2}\nu }0.$$
(115)
Therefore the apparently state dependent contribution (109) from the $`k=1`$ mode vanishes in any infrared finite state parameterized by finite $`A`$ and $`B`$, while the $`k>1`$ contribution is still given by (112), and we have the result
$$\epsilon _R=\underset{\xi 0}{lim}\epsilon _{BD}(\xi )|_{m=0}\frac{3}{32\pi ^2}\left(\frac{R}{12}\right)^2=\epsilon _{AF},$$
(116)
in any $`(A,B)`$ “vacuum” state with $`\alpha _k=1,\beta _k=n_k=0`$ for $`k>1`$, upon restoring physical units. The difference here is just that required to give the Allen-Folacci (AF) renormalized expectation value, provided the limit of the BD value is taken in the same order of $`m=0`$, $`\xi 0`$ that we have evaluated (109). Since from (29)
$$\epsilon _{BD}(\xi )|_{m=0}=\frac{3}{16\pi ^2}\left(\frac{R}{12}\right)^2\left[\frac{1}{180}\frac{1}{6}(6\xi 1)^2\right],$$
(117)
and both the BD and AF values are de Sitter invariant,
$`\underset{\xi 0}{lim}T_{ab}_{ren}(\xi )|_{m=0}`$ $`=`$ $`{\displaystyle \frac{3g_{ab}}{16\pi ^2}}\left({\displaystyle \frac{R}{12}}\right)^2\left({\displaystyle \frac{1}{180}}{\displaystyle \frac{1}{6}}{\displaystyle \frac{1}{2}}\right)`$ (118)
$`=`$ $`{\displaystyle \frac{119}{960\pi ^2}}\left({\displaystyle \frac{R}{12}}\right)^2g_{ab}={\displaystyle \frac{119R^2}{\mathrm{138\hspace{0.17em}240}\pi ^2}}g_{ab}=T_{ab}_{AF},`$ (119)
and we obtain the de Sitter invariant AF result for all $`(A,B)`$ “vacuum” states in the massless, minimally coupled case.
With the above careful analysis of the spatially homogeneous mode and the independence of the asymptotic value of its energy-momentum tensor in an arbitrary $`k=1`$ infrared finite state characterized by $`A`$ and $`B`$, it is now straightforward to carry out the proof of the attractor behavior of the AF state for an arbitrary UV finite physical state with $`m=\xi =0`$, by allowing the $`k>1`$ modes to have $`\beta _k`$ and $`n_k`$ different from zero. Substituting (91), (93), (100), and (108) into (15), (17) with $`m=\xi =0`$, and using (92) gives
$`\epsilon =T_{}^{0}{}_{0}{}^{}_{ren}`$ $`=`$ $`T_{}^{0}{}_{0}{}^{}_{AF}+(1+2n_1){\displaystyle \frac{|A|^2\mathrm{cos}^6\eta }{\pi ^2}}`$ (124)
$`+{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \underset{k=2}{\overset{\mathrm{}}{}}}\{[2n_k+2(1+2n_k)|\beta _k|^2][k^3\mathrm{cos}^4\eta +k(\mathrm{cos}^4\eta +{\displaystyle \frac{1}{2}}\mathrm{cos}^2\eta )]`$
$`+(1+2n_k)[(\beta _k\alpha _k^{}e^{2ik\eta }+\beta _k^{}\alpha _ke^{2ik\eta })k(\mathrm{cos}^4\eta +{\displaystyle \frac{1}{2}}\mathrm{cos}^2\eta )`$
$`+i(\beta _k\alpha _k^{}e^{2ik\eta }\beta _k^{}\alpha _ke^{2ik\eta })k^2\mathrm{cos}^3\eta \mathrm{sin}\eta ]\},`$
$`T_{ren}`$ $`=`$ $`T_{AF}+(1+2n_1){\displaystyle \frac{2|A|^2\mathrm{cos}^6\eta }{\pi ^2}}`$ (128)
$`{\displaystyle \frac{1}{4\pi ^2}}{\displaystyle \underset{k=2}{\overset{\mathrm{}}{}}}\{[2n_k+2(1+2n_k)|\beta _k|^2]k\mathrm{cos}^2\eta `$
$`+(1+2n_k)[(\beta _k\alpha _k^{}e^{2ik\eta }+\beta _k^{}\alpha _ke^{2ik\eta })(2k^3\mathrm{cos}^4\eta +k\mathrm{cos}^2\eta )`$
$`+\mathrm{\hspace{0.33em}2}i(\beta _k\alpha _k^{}e^{2ik\eta }\beta _k^{}\alpha _ke^{2ik\eta })k^2\mathrm{cos}^3\eta \mathrm{sin}\eta ]\}.`$
Provided the $`k`$ sums converge, it is clear that all the state dependent terms contain at least one factor of $`a^2=\mathrm{cos}^2\eta `$, and so vanish in the limit of $`\eta \frac{\pi }{2}`$. However, the requirement that the state be fourth order adiabatic just guarantees this convergence, for the same reason as in the previous analysis in spatially flat coordinates. Indeed we have
$`|\beta _k|`$ $`=`$ $`{\displaystyle \frac{C(k)}{k^4}},`$ (129)
$`n_k`$ $`=`$ $`{\displaystyle \frac{N(k)}{k^4}},`$ (130)
for some $`C(k)`$ and $`N(k)`$ that vanish as $`k\mathrm{}`$. This is sufficient to guarantee the absolute convergence of all terms in the sums. Since all state dependent terms are multiplied by at least two powers of $`\mathrm{cos}\eta =a^1`$, which vanishes in the late time limit $`\eta \frac{\pi }{2}`$, we conclude that any UV and IR admissible state of the massless, minimally coupled scalar field has an energy-momentum tensor which approaches the AF value, $`T_{ab}_{AF}`$ in the late time limit $`\eta \frac{\pi }{2}`$.
If one considers the contributions of the state dependent terms to the energy density and trace from the $`k>1`$ modes, for $`\nu `$ not exactly $`\frac{3}{2}`$, the kinematics is essentially the same as our previous analysis of the $`k=1`$ mode. Their contribution also falls off proportional to $`\xi |\alpha _k+\beta _k|^2a^{2\nu 3}`$ at late times for $`\nu `$ close to $`\frac{3}{2}`$. However, there is no compensating large factor coming from the pole in the $`\mathrm{\Gamma }`$ function normalization constant as there is for the $`k=1`$ mode. Hence, the coefficient of this slow fall off goes to zero as the massless, minimally coupled limit is approached. That is, exactly at $`m=\xi =0`$ when the contribution of the $`k>1`$ modes no longer falls off and they could in principle contribute to the asymptotic value of $`T_{ab}_{ren}`$ at late times, at that very point their coefficient vanishes identically and they make no contribution at all. Thus, at precisely $`m=0`$ and $`\xi =0`$ the difference between the BD and AF values can be attributed entirely to the additional condensate in the spatially homogeneous $`k=1`$ mode alone, and there are no slow transient modes in our explicit analysis of the massless, minimally coupled energy and trace.
These considerations are relevant to the case when $`\nu =\frac{3}{2}`$ but $`m`$ and $`\xi `$ are separately non-zero. The calculation is almost identical but the conclusion is different, since now the finite $`k>1`$ mode sum is multiplied by a coefficient $`\xi `$ which does not vanish. The entire mode sum of state independent terms from $`k=2`$ up to $`a(\eta )`$ do not fall off at late times and in fact all add up, to give a contribution proportional to $`\mathrm{\Phi }^2_{k=2}^ak^1\mathrm{log}a`$ in $`T_{ab}_{ren}`$, which grows linearly in comoving time. The explicit expression is most conveniently calculated by using (22) to divide the energy-momentum tensor into a state dependent numerical part, $`T_{ab}_n`$ and a state independent analytic part, $`T_{ab}_{an}`$. They are separately conserved . The quantities $`T_{ab}_d`$ and $`T_{ab}_{an}`$ are given in Ref. . By substituting Eqs. (108) and (93) into Eqs. (15) and (17), and subtracting the expression for $`T_{ab}_d`$ given in Ref. , we find that $`T_{ab}_n`$ approaches a state dependent, finite constant in the limit $`\eta \frac{\pi }{2}`$. However, the quantity $`T_{ab}_{an}`$ has a term that is proportional to $`\mathrm{\Phi }^2`$. At late times this term behaves as $`\mathrm{log}a`$ and dominates. In fact
$`\epsilon _{ren}`$ $``$ $`{\displaystyle \frac{3\xi }{4\pi ^2}}\mathrm{log}(\mathrm{cos}\eta )+q_1,`$ (132)
$`T_{ren}`$ $``$ $`{\displaystyle \frac{3\xi }{\pi ^2}}\mathrm{log}(\mathrm{cos}\eta )+q_2,`$ (133)
with the constants $`q_1`$ and $`q_2`$ dependent on the state of the field and constrained by the conservation equation. Since $`\mathrm{log}(\mathrm{cos}\eta )t`$ in comoving time, the energy-momentum tensor grows linearly with $`t`$ at late times. For $`\xi >0`$ (and $`m^2<0`$) the linear growth in comoving time decreases the effective cosmological “constant”, whereas for $`\xi <0`$ (and $`m^2>0`$) it increases it. In either case the back-reaction of the energy-momentum of the quantum scalar field for $`m^2+\xi R=0`$ (but $`m^2=\xi R0`$) on the geometry certainly cannot be neglected at late times since $`T_{ab}`$ grows without bound for any physical state. The fact that a massive non-minimal field with $`\xi <0`$ can induce an effective cosmological “constant” due to inflationary particle production was noted in Ref. using a different approach. It is an interesting open question whether this linearly growing behavior (in proper time) carries over to the physically more relevant case of one-loop gravitons, since the mode functions for gravitons in a particular gauge obey the same equation in a RW spacetime as do the mode functions for a massless minimally coupled scalar field .
Finally, we note that for all real values of $`\nu `$, the analysis of the $`k=1`$ mode in this Section may be extended to all of the higher $`k`$ modes, since the late time behavior of the higher $`k`$ modes is determined by that of the hypergeometric function in (89), which gives
$$f_k\left(\eta \frac{\pi }{2}\right)\left[\frac{\mathrm{\Gamma }\left(k+\frac{1}{2}\nu \right)}{2\mathrm{\Gamma }\left(k+\frac{1}{2}+\nu \right)}\right]^{\frac{1}{2}}\frac{\mathrm{\Gamma }(2\nu )}{\mathrm{\Gamma }(\frac{1}{2}+\nu )}\frac{(i)^k}{k!}\left(\frac{i\mathrm{sec}\eta }{2}\right)^{\nu \frac{1}{2}}.$$
(134)
Thus all of the higher $`k`$ modes behave as $`a^{\nu \frac{1}{2}}`$ and give the (unrenormalized) energy density the leading order late time behavior
$$a^{2\nu 3}\left[\left(\nu \frac{3}{2}\right)^2+12\xi (\nu 1)+m^2\right]\left[\frac{\mathrm{\Gamma }(2\nu )}{2^\nu \mathrm{\Gamma }\left(\nu +\frac{1}{2}\right)}\right]^2\underset{k=1}{\overset{\mathrm{}}{}}(1+2n_k)\frac{\mathrm{\Gamma }\left(k+\frac{1}{2}\nu \right)}{\mathrm{\Gamma }\left(k+\frac{1}{2}+\nu \right)(k!)^2}|\alpha _k+e^{i\pi (k+\nu \frac{1}{2})}\beta _k|^2,$$
(135)
which grows (or shrinks) like $`a^{2\nu 3}`$ for an arbitrary physical state unless either the coefficient vanishes or $`\alpha _k+e^{i\pi \left(k+\nu \frac{1}{2}\right)}\beta _k=0`$ for all $`k`$. However, the latter is impossible since it is inconsistent with the requirement that all physically allowable states be fourth order adiabatic states, which requires that $`\alpha _k1`$ and $`\beta _k0`$ sufficiently fast as $`k\mathrm{}`$, in order for the renormalized energy density to be UV finite. Since the fourth order adiabatic subtraction behaves at most like a constant at late times, the $`a^{2\nu 3}`$ behavior is not affected by the UV subtraction, and indeed the sum in (135) converges at large $`k`$ for any UV admissable state. Hence after the adiabatic four subtraction the leading order late time behavior indicated by (135) survives unless the first bracket in front of the entire expression vanishes. The quantity in this bracket is identical to the factor in (85) found in the flat spatial section analysis of the last Section. The corresponding quantity for the trace is given by the factor in curly brackets in (86), which is a similar combination of $`\nu ,\xi `$, and $`m^2`$. If we require that both of these factors vanish identically, to eliminate the leading order behavior in all components of $`T_{ab}_{ren}`$, then these two conditions plus the defining relation (26) give
$`m^2`$ $`=`$ $`{\displaystyle \frac{\nu (2\nu 3)(2\nu 1)}{4(\nu 2)}},`$ (136)
$`\xi `$ $`=`$ $`{\displaystyle \frac{(2\nu 3)}{8(\nu 2)}}.`$ (137)
Thus, except for $`\nu =2`$, for any given $`\nu `$ there is always one value of $`m^2`$ and one value of $`\xi `$ for which the coefficients of these leading order terms in $`T_{ab}_{ren}`$ vanish. The next to leading order terms go like $`a^{2\nu 5}`$, and grow without bound at late times in any case when $`\nu >\frac{5}{2}`$.
The above analysis implies that, for most values of $`m^2`$ and $`\xi `$, when $`\nu >\frac{3}{2}`$ the leading order terms in the components of the energy-momentum tensor grow without bound like $`a^{2\nu 3}`$ in de Sitter space for any physically admissable initial state of the scalar field. These values of $`\nu `$ correspond to the purely tachyonic cases $`m^2+\xi R<0`$.
## V Numerical Studies
In this Section we display numerical results for various values of $`\nu `$. All results are given in dimensionless units where $`\alpha =1`$ and $`R=12`$. The quantum state in each case is a fourth order adiabatic state matched to the vacuum at some initial time $`\eta _0`$. That is, we choose initial conditions for the mode function and its first derivative to be
$`\psi _k(\eta _0)`$ $`=`$ $`\psi _k^{(4)}(\eta _0),`$ (138)
$`{\displaystyle \frac{\mathrm{d}\psi _k}{\mathrm{d}\eta }}(\eta _0)`$ $`=`$ $`{\displaystyle \frac{\mathrm{d}\psi _k^{(4)}}{\mathrm{d}\eta }}(\eta _0),`$ (139)
with $`\psi _k^{(4)}`$ the fourth order adiabatic mode function defined by (30) with phase measured from the initial time $`\eta _0`$. The initial time was chosen to be $`t_0=1`$ in comoving coordinates, i.e. $`\mathrm{sec}\eta _0=\mathrm{cosh}(1)=1.54308\mathrm{}`$.
For $`\mathrm{}(\nu )<\frac{3}{2}`$ the proof in Section III states that the energy-momentum tensor approaches the Bunch-Davies value (29) at late times for an arbitrary physically admissable state. This occurs due to a redshifting of the state dependent part of the energy-momentum tensor. The initial state dependent transient contributions fall off like
$$a^{2\nu 3}=(\mathrm{cosh}t)^{2\nu 3}.$$
(140)
Thus the characteristic time to approach the BD value is
$$\tau =\frac{1}{32\nu }.$$
(141)
In Fig. 1 we plot the renormalized energy density for these adiabatic initial conditions with $`m=0`$ and $`\xi =\frac{1}{7}`$. For this relatively large value of $`\xi `$, the characteristic time $`\tau `$ is of order one and the energy density approaches the BD value within one expansion time.
For smaller values of $`\xi `$ the initial value transients persist for longer times. Our analysis in the previous Section shows that the $`k=1`$ contribution is essential for the shift from the BD to AF value as $`\xi 0`$. Since the $`k=1`$ mode in an arbitrary physical state contributes to the energy density the value $`\epsilon _1`$ given by (109), up to terms which fall off like $`a^2`$, for arbitrary $`\alpha _1,\beta _1`$, and $`n_1`$, while the BD state has $`\alpha _1=1,\beta _1=n_1=0`$, it is clear that the difference of the renormalized energy density from the BD value coming from this mode is
$$\epsilon _R\epsilon _{BD}=\epsilon _1\frac{3}{32\pi ^2}\left(\frac{R}{12}\right)^2a^{2\nu 3}.$$
(142)
For the adiabatic initial conditions here we find that
$$\psi _1(\eta _0)=\frac{1}{\sqrt{2W_1^{(4)}(\eta _0)}}\frac{(\alpha _1+\beta _1)\mathrm{sec}\eta _0}{2\sqrt{\frac{3}{2}\nu }}.$$
(143)
Since $`W_1^{(4)}(\eta _0)`$ remains finite as $`\xi 0`$, provided that $`\eta _00`$, it follows that in this limit $`|\alpha _1+\beta _1|^232\nu \xi `$ also goes to zero. Hence for small $`\xi `$, $`\epsilon _R`$ goes to a value close to ( but slightly larger than) the AF value after a time of order one. This is observed in both Figs. 2 and 3 for $`\xi =\frac{1}{100}`$ and $`\xi =\frac{1}{1000}`$, respectively.
Further, the energy-momentum tensor contains a term proportional to $`\xi \mathrm{\Phi }^2`$ which would grow linearly in comoving time for $`\nu `$ very close to $`\frac{3}{2}`$, except for the factor of $`a^{2\nu 3}`$ that damps it to zero at very late times $`t\tau `$. Hence on times $`1<t\tau `$, where the $`a^{2\nu 3}`$ factor is essentially constant, we should expect this linear growth of $`\epsilon `$ in comoving time with a slope proportional to $`\xi `$, given by (132). This behavior is demonstrated in Figs. 2 and 3. In the latter case $`\xi `$ is so small that $`\alpha _1+\beta _1`$ and the slope are nearly vanishing and the energy density stays close to the AF value until times of order $`\tau =125`$, which is much larger than the times shown in Fig. 3. When $`\xi =0`$ (still keeping $`m=0`$) both $`\alpha _1+\beta _1`$ and the slope vanishes identically, so the energy-momentum tensor goes to the AF value and remains there. This demonstrates that the limit $`\mathrm{}(\nu )\frac{3}{2}`$ is quite continuous when viewed at finite times with well-defined physical initial conditions, although the late time limit is discontinuous.
In the case $`\nu =\frac{3}{2}`$ there are two distinct behaviors depending on whether $`m`$ and $`\xi `$ are separately vanishing or not. In the massless, minimally coupled case we proved in Section IV that the Allen-Folacci value is a fixed point at late times. In Figure 4 we show the approach of the energy density, pressure, and trace to their Allen-Folacci values for the massless minimally coupled field. The field is in an “n-particle” state with $`n_1=n_2=2`$ and $`n_k=0`$ for all $`k>2`$. In Figure 5 we show the behavior of the energy density and trace for the case $`m=0.3`$ and $`\xi =0.0075`$ when the field is in a fourth order adiabatic vacuum state.
When $`\mathrm{}(\nu )>\frac{3}{2}`$ the analysis at the end of Section IV shows that the leading order state dependent terms in the energy-momentum tensor will generally grow exponentially with time. This is illustrated in Fig. 6. The exponential growth that occurs here is similar to the well known one for the classical scalar field when $`m`$ and $`\xi `$ have values such that $`\nu >\frac{3}{2}`$ . Such fields are tachyonic and presumbably of little physical interest unless interactions are added to stabilize them.
## VI Infrared Scaling and the Generalized Anomaly
We have found that for all $`\mathrm{}(\nu )<\frac{3}{2}`$ the renormalized expectation value of $`T_{ab}`$ approaches the de Sitter invariant Bunch-Davies value for any physically admissable initial state, whereas it approaches the de Sitter invariant Allen-Folacci value for any physically admissable initial state in the massless, minimally coupled case. Since all the initial state dependence vanishes asymptotically, these state independent de Sitter invariant fixed point values for $`T_{ab}_{ren}`$ must be purely geometrical in origin. Indeed, both the Bunch-Davies point-splitting calculation and the Hadamard calculation of Allen-Folacci rely only on the properties of the two-point function of the scalar field $`G(x,x^{})`$ for $`x^{}x`$. Hence, the BD and AF asymptotic values of $`T_{ab}_{ren}`$ are certainly “pseudo-local” in the terminology of Ref. , i.e. they are expressible in terms of purely local functions of the RW scale factor $`a(\eta )`$ and its derivatives.
If we specialize now to zero mass, $`m=0`$, then on simple dimensional grounds the asymptotic $`T_{ab}_{ren}`$ can be expressed purely in terms of local conserved tensors of fourth adiabatic order. Although we have used adiabatic subtraction methods to renormalize $`T_{ab}`$ it is known that the value of $`T_{ab}_{ren}`$ so obtained is equal to that in a fully covariant procedure such as dimensional regularization or covariant point-splitting . In a fully covariant procedure, which yields a local conserved tensor of fourth adiabatic order, only the Riemann tensor together with its covariant derivatives and contractions can appear. Hence $`T_{ab}_{ren}`$ for $`m=0`$ must be expressible entirely in terms of such local geometrical tensors.
In four dimensions the only such local tensors are linear combinations of $`{}_{}{}^{(1)}H_{ab}^{}`$, $`{}_{}{}^{(2)}H_{ab}^{}`$, and $`{}_{}{}^{(3)}H_{ab}^{}`$, where
$`{}_{}{}^{(1)}H_{ab}^{}`$ $``$ $`{\displaystyle \frac{1}{\sqrt{g}}}{\displaystyle \frac{\delta }{\delta g^{ab}}}{\displaystyle \sqrt{g}R^2\mathrm{d}^4x}`$ (144)
$`=`$ $`2g_{ab}\text{ }R2_a_bR+2RR_{ab}{\displaystyle \frac{1}{2}}g_{ab}R^2,`$ (145)
vanishes in de Sitter spacetime and
$${}_{}{}^{(3)}H_{ab}^{}=R_a^cR_{cb}\frac{2}{3}RR_{ab}\frac{1}{2}R_{cd}R^{cd}g_{ab}+\frac{1}{4}R^2g_{ab}.$$
(146)
In RW spacetimes, which are all conformally flat, the tensor $`{}_{}{}^{(2)}H_{ab}^{}`$ is proportional to $`{}_{}{}^{(1)}H_{ab}^{}`$ and hence vanishes as well. Therefore the only non-trivial fourth order conserved geometrical tensor in de Sitter spacetime is $`{}_{}{}^{(3)}H_{ab}^{}`$ and we conclude that the fixed point BD and AF values found in our previous analysis are proportional to
$${}_{}{}^{(3)}H_{ab}^{}=\frac{R^2}{48}g_{ab}=3\left(\frac{R}{12}\right)^2g_{ab}.$$
(147)
Furthermore, since
$${}_{}{}^{(3)}H_{ab}^{}g^{ab}=R_{ab}R^{ab}+\frac{R^2}{3}=\frac{1}{2}\left(R_{abcd}R^{abcd}4R_{ab}R^{ab}+R^2\right)\frac{G}{2},$$
(148)
in RW spacetimes (where the Weyl tensor vanishes), the coefficient of $`{}_{}{}^{(3)}H_{ab}^{}`$ is proportional to the coefficient of the Gauss-Bonnet integrand in the trace of $`T_{ab}`$. Such a term in the trace is known to correspond to a non-local but nevertheless fully covariant action and this action is precisely the same as that generated by the trace anomaly of free conformal fields . Since we have obtained fixed point results for the asymptotic values of $`T_{ab}_{ren}`$ for massless fields in de Sitter space which are purely geometrical and proportional to $`{}_{}{}^{(3)}H_{ab}^{}`$, even for non-conformal massless fields, we can give a definite meaning to the value of the proportionality coefficient and the non-local anomaly-like term in the effective action even when $`\xi \frac{1}{6}`$.
Let us define the generalized anomaly coefficient by fixing the normalization
$$\underset{t\mathrm{}}{lim}T_{ab}_{ren}=\frac{Q^2}{16\pi ^2}^{(3)}H_{ab}=\frac{3Q^2}{16\pi ^2}\left(\frac{R}{12}\right)^2g_{ab}.$$
(149)
With this normalization we find from the asymptotic value of $`T_{ab}_{ren}`$ for a scalar field in de Sitter space that
$$Q^2=\{\begin{array}{cc}Q_{BD}^2\hfill & =\frac{1}{180}\frac{1}{6}(6\xi 1)^2,m=0,\xi >0,\hfill \\ Q_{AF}^2\hfill & =\frac{1}{180}\frac{1}{6}\frac{1}{2}=\frac{119}{180},m=0,\xi =0.\hfill \end{array}$$
(150)
The value of $`Q^2`$ for a conformally invariant field ($`m=0,\xi =\frac{1}{6}`$) is $`\frac{1}{180}`$ and corresponds to the pure trace anomaly coefficient. The first member of (150) provides the generalization of this coefficient away from the conformal case. The discontinuous behavior at $`\xi =0`$ has been discussed in Ref. . As we have seen it arises from the singular behavior of the spatially constant zero mode of the massless, minimally coupled field, which is non-oscillatory and hence cannot be quantized as a Fock mode in the same fashion as the oscillatory modes. We discussed this discontinuity in detail in Section IV.
The connection of the tensor $`{}_{}{}^{(3)}H_{ab}^{}`$ with the trace anomaly may be seen from the general form of the effective action for the anomaly in a conformally flat space with metric
$$g_{ab}=e^{2\sigma }\overline{g}_{ab},$$
(151)
namely
$$S_{\mathrm{e}ff}=\frac{Q^2}{16\pi ^2}\mathrm{d}^4x\sqrt{\overline{g}}\left[\sigma \overline{\mathrm{\Delta }}_4\sigma +\frac{1}{2}\left(\overline{G}\frac{2}{3}\overline{\text{.09}.09}\overline{R}\right)\sigma \right],$$
(152)
where $`\mathrm{\Delta }_4=\text{ }\text{ }\text{ }\text{ }\text{ }^2+2R^{ab}_a_b\frac{2}{3}R\text{ }\text{ }\text{ }\text{ }\text{ }^2+\frac{1}{3}(^aR)_a`$ is the unique fourth order differential operator acting on scalars which is conformally covariant. A fully covariant but non-local form of the effective action (152) can be obtained by solving $`\sqrt{g}\left(G\frac{2}{3}\text{ }\text{ }\text{ }\text{ }\text{ }R\right)=\sqrt{\overline{g}}\left(\overline{G}\frac{2}{3}\overline{\text{ }\text{ }\text{ }\text{ }\text{ }}\overline{R}\right)+4\sqrt{\overline{g}}\overline{\mathrm{\Delta }}_4\sigma `$ for $`\sigma `$. In that non-local form all reference to the separation of the metric into background and conformal factor as in (151) disappears.
The energy-momentum tensor following from the variation of the local form of the effective action (152) with respect to the background metric $`\overline{g}_{ab}`$, is given by Eq. (2.9) of Ref. . The form of this energy-momentum tensor simplifies considerably on the Einstein space $`R\times S^3`$, where $`\overline{G}=\overline{\text{.09}.09}\overline{R}=0`$. If we set $`\sigma =\mathrm{log}a(\eta )`$ then this is equivalent to evaluating the components of this energy-momentum tensor in a general RW spacetime with closed spatial sections ($`\kappa =+1`$). Using the expressions for $`{}_{}{}^{(1)}H_{ab}^{}`$ and $`{}_{}{}^{(3)}H_{ab}^{}`$ in a general RW space in terms of $`a(\eta )`$ and its derivatives, we quickly find that
$$T_{ab}[\sigma ]=\frac{2}{\sqrt{g}}\frac{\delta S_{\mathrm{e}ff}}{\delta \overline{g}^{ab}}=\frac{Q^2}{16\pi ^2}\left[\frac{1}{18}^{(1)}H_{ab}^{(3)}H_{ab}+^{(3)}H_{ab}|_{R\times S^3}\right].$$
(153)
Hence the tensor $`{}_{}{}^{(3)}H_{ab}^{}`$ which is called “accidentally conserved” in Ref. is associated in fact with the existence of a non-local covariant effective action related to the trace anomaly, which has the local form (152) when the metric is conformally decomposed as in (151).
The last term in (153) would not have been present had we varied the fully covariant but non-local form of the anomalous effective action, in which all reference to the background $`\overline{g}_{ab}`$ drops out. Equivalently, it is just canceled if we add to (153) the Casimir energy on $`R\times S^3`$, which is determined by the anomaly by a further conformal transformation from flat space . In either case, we should drop this last term in (153) which depends on the arbitrary background $`\overline{g}_{ab}`$. Evaluating the first two terms on de Sitter space we find that it is exactly the value of the asymptotic form of the renormalized energy-momentum tensor $`T_{ab}`$ in de Sitter space (149), found previously, with the value of $`Q^2`$ for the massless scalar given by (150). Thus, the effective action (152), associated with the conformal trace anomaly, appears in the effective action for a massless scalar field, even for non-conformal couplings, $`\xi \frac{1}{6}`$.
Since it is associated with the anomaly, the physical significance of $`S_{\mathrm{e}ff}`$ in the quantum effective action for the scalar field is that it determines the scaling behavior of the field theory under global Weyl transformations of the background space. Using the fact that the Euclidean effective action is given by $`I_{\mathrm{e}ff}=S_{\mathrm{e}ff}`$ and that the Euclidean continuation of de Sitter space is $`S^4`$ with Euler number $`\chi =2`$, we can vary $`I_{\mathrm{e}ff}`$ with respect to a constant $`\sigma =\sigma _0`$ and obtain
$$\frac{I_{\mathrm{e}ff}}{\sigma _0}=\alpha \frac{\mathrm{d}I_{\mathrm{eff}}}{\mathrm{d}\alpha }=\frac{Q^2}{32\pi ^2}_{S^4}\mathrm{d}^4x\sqrt{g}G=Q^2\chi =2Q^2,$$
(154)
since $`\frac{\mathrm{d}}{\mathrm{d}\sigma _0}=\frac{\mathrm{d}}{\mathrm{d}\mathrm{log}\alpha }`$ is a global rescaling of the $`S^4`$. As a check, this relation can be verified explicitly by the $`\zeta `$ function evaluation of the Euclidean effective action,
$$I_{\mathrm{e}ff}=\frac{1}{2}\mathrm{Tr}\mathrm{log}\left(\text{.09}.09+m^2+\xi R\right)=\frac{1}{2}\frac{\mathrm{d}\zeta }{\mathrm{d}s}|_{_{s=0}},$$
(155)
where the generalized zeta function for the Euclidean continuation of the wave operator appearing in the Tr $`\mathrm{log}`$ is
$$\zeta (s)=\underset{n=0}{\overset{\mathrm{}}{}}d_n\lambda _n^s,$$
(156)
in terms of its eigenvalues $`\lambda _n`$ with degeneracy $`d_n`$ on $`S^4`$. This sum is convergent for $`\mathrm{}(s)>2`$ and defines a meromorphic function of $`s`$ which is analytic near $`s=0`$, where its derivative is required. Introducing a mass scale $`\mu `$ to keep $`\zeta (s)`$ dimensionless for all $`s`$ leads in a standard calculation to
$$\frac{1}{2}\frac{\mathrm{d}\zeta }{\mathrm{d}s}|_{_{s=0}}=\zeta (0)\mathrm{log}(\mu \alpha )+I_1(\nu ),$$
(157)
where $`I_1(\nu )`$ is a certain finite function of $`\nu `$, which from the definition of $`\nu `$ (26) becomes independent of $`\alpha `$ when $`m=0`$, and the value $`\zeta (0)`$ is given by
$$\zeta (0)=\frac{1}{12}\left(\nu ^4+\frac{1}{2}\nu ^2+\frac{17}{240}\right)=\frac{1}{90}\frac{1}{3}(6\xi 1)^2,$$
(158)
when $`m=0`$. By making use of (150), (155), and (157), we find that
$$\alpha \frac{\mathrm{d}}{\mathrm{d}\alpha }I_{\mathrm{e}ff}(m=0)=\zeta (0)=2Q^2,$$
(159)
is the behavior of the effective action for a massless scalar field with $`\xi >0`$ under global Weyl rescaling of the metric, exactly as predicted by(154) and the previous discussion based on the anomalous action (152).
When $`m=0`$ and $`\xi 0`$, the integral representation of the function $`I_1(\nu )`$ develops a logarithmic singularity, which can be traced to the vanishing of the $`n=0`$ eigenvalue in the expression (156) for $`\zeta (s)`$. In this case the $`n=0`$ mode must be excluded from the ultraviolet regulated sum over modes, which has the effect of adding one unit to $`\zeta (0)`$ in the infrared scaling behavior of the effective action , and accounts for the addition of $`\frac{1}{2}`$ in $`Q^2`$ in the minimally coupled case. Hence the discontinuous behavior of $`Q^2`$ found in (150) by our analysis of the asymptotic attractor behavior of the energy-momentum tensor in de Sitter space is precisely the same as that occuring in the effective action under global Weyl rescalings when $`m=0`$ and $`\xi 0`$.
Note also that the global Weyl variation is given in terms of the trace of the energy-momentum tensor by
$$\alpha \frac{\mathrm{d}}{\mathrm{d}\alpha }I_{\mathrm{e}ff}(m=0)=_{S^4}\mathrm{d}^4x\sqrt{g}T_{ren}=\frac{8\pi ^2}{3}\alpha ^4T_{ren}.$$
(160)
Since the $`\zeta `$ function method is fully covariant, the renormalized trace $`T_{ren}`$ computed in this way cannot contain non-covariant contributions, and must be expressible entirely in terms of local curvature invariants.
Finally we may consider relaxing the condition $`m=0`$. If $`m^2>0`$, the asymptotic form of the energy-momentum tensor for an arbitrary initial state is given by the BD value. However, if we expand the BD result (29) in powers of $`R/m^2`$, we find that it contains no adiabatic order four $`R^2`$ terms, beginning instead with $`R^3/m^2`$. Mathematically, this is because all terms up to fourth adiabatic order have been removed by the ultraviolet regulating procedure of point splitting or adiabatic subtraction. Hence the coefficient of $`{}_{}{}^{(3)}H_{ab}^{}`$ at asymptotically late times in an arbitrary physical state is given by
$$Q^2=0\mathrm{for}m^2>0,$$
(161)
and no anomalous $`S_{\mathrm{e}ff}`$ term appears in the quantum effective action for a massive field. This is consistent with the interpretation of the anomalous term (152) in the effective action of the scalar field as an infrared effect, since the fluctuations of a massive field decouple at large distances or late times, and should induce only strictly irrelevant operators in the effective action in the far infrared, which are suppressed by positive powers of $`R/m^2`$.
Only in the strictly conformal case, $`m=0`$ and $`\xi =\frac{1}{6}`$ is the infrared effective action equal to that obtained by ultraviolet methods, such as the $`a_2`$ coefficient in the Schwinger-DeWitt proper time expansion. However, our analysis of the fixed point behavior of $`T_{ab}_{ren}`$ in de Sitter space shows that its coefficient is connected with the global or extreme infrared scaling of the effective action, and this asymptotic behavior does not depend on the field being conformally invariant. The asymptotic attractor behavior of the energy-momentum tensor in de Sitter space defines an infrared scaling coefficient that reduces to the trace anomaly coefficient in the conformal case, but is a much more general concept than the trace anomaly coefficient, since it is well-defined for all massless fields, conformal or not. It is well-defined even for massive fields, although as (161) shows, it vanishes in this case.
## VII Discussion and Conclusions
In this paper we have considered a quantum scalar field in a fixed de Sitter background. We have studied the late time behavior of the renormalized energy-momentum tensor and have found two important cases in which, for arbitrary physically admissable states, the energy-momentum tensor approaches a particular value at late times. The values approached are those of the energy-momentum tensor in the Bunch-Davies and Allen-Folacci states. Thus, in this sense, these special quantum states behave as fixed point attractors.
In the case $`\mathrm{}(\nu )<\frac{3}{2}`$ we have shown that for all fourth order adiabatic states that are infrared finite the energy-momentum tensor approaches the BD value at late times. The longest time scale for the state dependent terms to redshift away is $`\tau =(32\nu )^1`$. This has been numerically verified for various values of $`m`$ and $`\xi `$ when the fields are in a fourth order adiabatic state. For the case in which $`0<\frac{3}{2}\nu <<1`$ and $`m`$ and $`\xi `$ are small, we numerically observe a more complicated behavior. The energy-momentum tensor quickly approaches the AF value and grows then linearly with comoving time. Our analytic proof implies that it must then slowly decay to the BD value. It is worth noting that for $`\mathrm{}(\nu )>0`$ the redshift of the state dependent terms in the quantum expectation value of $`T_{ab}`$ is slower than one might guess from the redshifting of classical matter or radiation, i.e. $`a^3`$ or $`a^4`$, respectively.
For the case $`\nu =\frac{3}{2}`$ we have to distinguish two different possibilities. If the field is massless and minimally coupled then we have proven that, for all fourth order adiabatic states that are infrared finite, the energy-momentum tensor approaches the AF value at late times. This is true for both vacuum and initially populated states. For any other values of $`m`$ and $`\xi `$ the renormalized energy-momentum tensor grows linearly with comoving time, indicating that back-reaction effects need to be taken into account. The sign of the linear growth depends on the sign of $`\xi `$.
For the tachyonic cases $`\nu >\frac{3}{2}`$ there is no attractor state. Instead, for most values of $`m^2`$ and $`\xi `$ the renormalized energy-momentum tensor grows like $`a^{2\nu 3}`$ at late times for all physically admissable states. Thus back-reaction effects again need to be taken into account.
For the cases in which either the Bunch-Davies or Allen-Folacci states serve as attractors in the above sense and the mass, $`m`$ of the field is zero, we have shown how these results are connected to the appearance of a certain non-local term in the quantum effective action for the scalar field. This term gives rise to the local geometrical tensor $`{}_{}{}^{(3)}H_{ab}^{}`$ in the asymptotic form of $`T_{ab}`$ at late times, and also determines the global scaling behavior of the effective action for massless fields. Determining this scaling behavior and relating it to the asymptotic $`T_{ab}`$ in de Sitter space has allowed us to generalize the notion of the trace anomaly to massless, non-conformally coupled scalar fields, in the sense that the coefficient of this non-local term in the effective action is well defined even for $`\xi \frac{1}{6}`$.
The interplay between the UV and IR properties of the state and the mode sums contributing to the energy-momentum tensor is a theme running through all of these results. We had to be careful to remove the UV divergences from the unrenormalized $`T_{ab}`$ in order to analyze the late time limit. Since we have found state independent de Sitter invariant results in both the $`\nu <\frac{3}{2}`$ and massless, minimally coupled cases, and since all the state dependence resides in the finite $`k`$ modes, their form at high $`k`$ being restricted by the requirement of matching the adiabatic order four vacuum, it is clear that state independent results for $`T_{ab}_{ren}`$ are possible only because the contribution to the BD or AF expectation value comes from arbitrarily large $`k`$ at very late times. In fact, inspection of the renormalized expectation value of $`T_{ab}`$ expressed as a mode sum, after the fourth order adiabatic subtraction has been made, shows that the finite contribution comes from $`ka`$, i.e. when the physical wavelength of the mode is of order of the de Sitter horizon. At very late times, this corresponds to arbitrarily large values of the coordinate wave number $`k`$.
The finite difference between the BD value and the AF value in the $`m=\xi =0`$ case comes entirely from the $`k=1`$ mode in closed spatial sections, which is a purely IR effect. This leads to a finite discontinuity in the infrared scaling properties of a massless field since the value of the energy-momentum tensor at $`m=\xi =0`$ is different from its value in the limit $`m=0`$ and $`\xi 0`$. The appearance of the $`{}_{}{}^{(3)}H_{ab}^{}`$ tensor and the corresponding non-local action is quite generic for massless fields; its coefficient $`Q^2`$ vanishes only if the mass is non-zero. Hence we should expect that, although they are certainly not conformal, gravitons will also contribute to this same infrared effective action with a finite value of $`Q^2`$, which can be determined in the same way by a background de Sitter calculation of their quantum $`T_{ab}`$ at late times. We plan to present the results of this calculation in a future publication.
It is also interesting to note that the coefficient of the generalized trace anomaly, $`Q^2`$, is not generically positive, in contrast to all the previously known examples of massless conformal fields . It appears that the reason for this is that a positive $`Q^2`$ comes from the ultraviolet behavior of $`T_{ab}`$ for conformal fields, while the infrared behavior of non-conformal fields can contribute a negative value. For any $`Q^20`$ the new term in the effective action leads to dynamics for the RW scale factor which is quite different from the Einstein theory, and remains to be investigated in a full dynamical back-reaction calculation.
###### Acknowledgements.
Several of us would like to thank V. Sahni for helpful discussions, and the Institute for Nuclear Theory, University of Washington, where some of this work was completed. P. R. A. would like to thank T-8, Los Alamos National Laboratory for its hospitality. This work was supported in part by grant numbers DMR-9403009 and PHY-9800971 from the National Science Foundation. It was also supported in part by contract number W-7405-ENG-36 from the Department of Energy.
## A The cases of imaginary and integer values of $`\nu `$
In this appendix we show that the proof in Section III works for imaginary values of $`\nu `$ as well as for $`\nu =0,1`$. For imaginary values of $`\nu `$ it is useful to write $`\nu =i\gamma `$ with $`\gamma `$ a positive real number. Then the formulas in Eq. (52) and (55) are the same as before. However the values of the $`A_i`$ are different. They are given by
$`A_1`$ $`=`$ $`{\displaystyle \frac{\pi }{2k}}e^{\gamma \pi }[\mathrm{coth}(\gamma \pi )\mathrm{csch}(\gamma \pi )((1+2n_k)|c_2|^2+n_k){\displaystyle \frac{1}{2}}\mathrm{coth}(\gamma \pi )\mathrm{csch}(\gamma \pi )(1+2n_k)(c_1c_2^{}+c_1^{}c_2)`$ (A2)
$`{\displaystyle \frac{1}{2}}\mathrm{csch}(\gamma \pi )(1+2n_k)(c_1c_2^{}c_1^{}c_2)],`$
$`A_2`$ $`=`$ $`{\displaystyle \frac{\pi }{k}}e^{\gamma \pi }[(1+\mathrm{csch}^2(\gamma \pi )((1+2n_k)|c_2|^2+n_k){\displaystyle \frac{1}{2}}\mathrm{csch}^2(\gamma \pi )(1+2n_k)(c_1c_2^{}+c_1^{}c_2)],`$ (A3)
$`A_3`$ $`=`$ $`{\displaystyle \frac{\pi }{2k}}e^{\gamma \pi }[\mathrm{coth}(\gamma \pi )\mathrm{csch}(\gamma \pi )((1+2n_k)|c_2|^2+n_k){\displaystyle \frac{1}{2}}\mathrm{coth}(\gamma \pi )\mathrm{csch}(\gamma \pi )(1+2n_k)(c_1c_2^{}+c_1^{}c_2)`$ (A5)
$`+{\displaystyle \frac{1}{2}}\mathrm{csch}(\gamma \pi )(1+2n_k)(c_1c_2^{}c_1^{}c_2)].`$
The last terms in the expressions for $`A_1`$ and $`A_3`$ are imaginary. Substitution into Eqs. (52) and (55) shows that the contributions of these terms to the energy-momentum tensor cancel.
Noting that $`\nu `$ is imaginary, one sees that the argument for the vanishing of the first integral in Eq. (64) is unchanged. Since the terms in the integrand are being bounded for the second integral one must take the real part of $`\beta _i`$ and substitute that for $`\beta _i`$ in Eqs. (67) and (79). After making this substitution it is still the case that each term in (79) vanishes in the limit $`\eta 0^{}`$. The argument for the vanishing of the third integral in Eq.(64) is unchanged when $`\nu `$ is imaginary. Thus, the proof is valid for imaginary values of $`\nu `$.
For $`\nu =n=0,1`$ we use the identity
$`H_n^{(1)}(z)=J_n(z)+iN_n(z),`$ (A6)
$`H_n^{(2)}(z)=J_n(z)iN_n(z),`$ (A7)
and then use the well known power series solutions for $`N_0`$ and $`N_1`$ to write
$$N_n(z)=\frac{2}{\pi }J_n(z)\mathrm{log}(z)+P_n(z),$$
(A8)
where $`P_n(z)`$ is a series of the form
$$P_n(z)=\underset{j=0}{\overset{\mathrm{}}{}}b_{nj}z^{2jn}.$$
(A9)
Because of the $`\mathrm{log}(z)`$ term it is useful to write
$$\frac{\mathrm{d}N_n(z)}{\mathrm{d}z}=\frac{2}{\pi }\frac{\mathrm{d}J_n(z)}{\mathrm{d}z}\mathrm{log}(z)+Q_n(z),$$
(A10)
with $`Q_n(z)`$ a series of the form
$$Q_n(z)=\underset{j=0}{\overset{\mathrm{}}{}}ϵ_{nj}z^{2jn1}.$$
(A11)
We can conclude that Eqs. (52) and (55) remain the same but with different expressions for $`A_i`$, $`S_i`$, and $`\beta _i`$. The new expressions for the $`A_i`$ are given by
$`A_1`$ $`=`$ $`{\displaystyle \frac{\pi }{2k}}[(1+{\displaystyle \frac{4}{\pi ^2}}(\mathrm{log}z)^2)((1+2n_k)|c_2|^2+n_k)+{\displaystyle \frac{1}{2}}(1{\displaystyle \frac{4}{\pi ^2}}(\mathrm{log}z)^2)(1+2n_k)(c_1c_2^{}+c_1^{}c_2)`$ (A13)
$`+{\displaystyle \frac{2i}{\pi }}\mathrm{log}z(1+2n_k)(c_1c_2^{}c_1^{}c_2)],`$
$`A_2`$ $`=`$ $`{\displaystyle \frac{\pi }{k}}\left[{\displaystyle \frac{2}{\pi }}\mathrm{log}z((1+2n_k)|c_2|^2+n_k){\displaystyle \frac{1}{\pi }}\mathrm{log}z(1+2n_k)(c_1c_2^{}+c_1^{}c_2)+{\displaystyle \frac{i}{2}}(1+2n_k)(c_1c_2^{}c_1^{}c_2)\right],`$ (A14)
$`A_3`$ $`=`$ $`{\displaystyle \frac{\pi }{2k}}\left[((1+2n_k)|c_2|^2+n_k){\displaystyle \frac{1}{2}}(1+2n_k)(c_1c_2^{}+c_1^{}c_2)\right].`$ (A15)
Note that the $`A_i`$ are now functions of both $`k`$ and $`\mathrm{log}(z)`$. The new expressions for the $`\beta _i`$ and $`S_i`$ are given in Table 2.
| $`i`$ | $`\beta _i`$ | $`S_i`$ |
| --- | --- | --- |
| 1 | $`5+2n`$ | $`z^5J_n^2(z)`$ |
| 2 | $`5`$ | $`z^5J_n(z)P_n(z)`$ |
| 3 | $`52n`$ | $`z^5P_n^2(z)`$ |
| 4 | $`3+2n`$ | $`z^3J_n^2(z)`$ |
| 5 | $`3`$ | $`z^3J_n(z)P_n(z)`$ |
| 6 | $`32n`$ | $`z^3P_n^2(z)`$ |
| 7 | $`3+2n`$ | $`z^4\frac{\mathrm{d}}{\mathrm{d}z}J_n^2(z)`$ |
| 8 | $`3`$ | $`z^4\left(\frac{\mathrm{d}}{\mathrm{d}z}J_n(z)\right)P_n(z)+J_n(z)Q_n(z))`$ |
| 9 | $`32n`$ | $`2z^4P_n(z)Q_n(z)`$ |
| 10 | $`3+2n`$ | $`z^5\left(\frac{\mathrm{d}}{\mathrm{d}z}J_n(z)\right)^2`$ |
| 11 | $`3`$ | $`z^5\left(\frac{\mathrm{d}}{\mathrm{d}z}J_n(z)\right)Q_n(z)`$ |
| 12 | $`32n`$ | $`z^5Q_n^2(z)`$ |
Table 2
Substitution these expressions into Eqs. (52) and (55) shows that for the first integral in (64) the expressions are of the same form, except that some terms have factors of $`\mathrm{log}(z)`$. However these do not prevent the terms in the first integral on the right hand side of (64) from vanishing asymptotically. For the second integral there are still terms of the form given in Eq. (67). However, some of them also have factors of $`\mathrm{log}(z)`$. Inserting factors of $`\mathrm{log}(z)`$ into Eq. (79) and computing the integrals, one finds that the terms all vanish in the limit $`\eta 0^{}`$. The factors of $`\mathrm{log}(z)`$ also do not affect the asymptotic vanishing of the third integral on the right in Eq. (64).
Therefore, in all cases where $`\mathrm{}(\nu )<\frac{3}{2}`$ the quantity $`T_{ab}_{SD}`$ vanishes in the limit $`\eta 0^{}`$.
## B The Harmonic Oscillator with $`\omega 0`$
In this Appendix we give a detailed discussion of the harmonic oscillator with vanishing frequency, pointing out the analogy with the $`\nu \frac{3}{2}`$ limit in de Sitter space. Consider a simple harmonic oscillator with Hamiltonian
$$H=\frac{1}{2}\dot{\varphi }^2+\frac{1}{2}\omega ^2\varphi ^2.$$
(B1)
The single degree of freedom $`\varphi `$ can be quantized by introducing the operator representation
$$\varphi (t)=\varphi (t)+a\psi (t)+a^{}\psi ^{}(t),$$
(B2)
where $`a^{}`$ and $`a`$ are creation and destruction operators obeying
$$[a,a^{}]=1.$$
(B3)
The canonical commutation relations $`[\varphi ,\dot{\varphi }]=i`$ are satisfied provided the mode function $`\psi `$ obeys the Wronskian condition (12). The equation of motion for the Heisenberg operator $`\varphi (t)`$ implies that the mode functions also satisfy
$$\ddot{\psi }+\omega ^2\psi =0.$$
(B4)
In order to satisfy the Wronskian condition we choose the fundamental normalized complex solution to this equation to be
$$f(t)\frac{1}{\sqrt{2\omega }}e^{i\omega t}.$$
(B5)
The general solution satisfying both the Wronskian condition and (B4) can be written as the linear combination
$$\psi (t)=\alpha f(t)+\beta f^{}(t),$$
(B6)
where
$$|\alpha |^2|\beta |^2=1.$$
(B7)
Thus, up to an irrelevant overall phase, the Bogoliubov coefficients can be parameterized in terms of two real parameters in the form
$`\alpha `$ $`=`$ $`\mathrm{cosh}\theta ,`$ (B8)
$`\beta `$ $`=`$ $`\mathrm{sinh}\theta e^{i\delta }.`$ (B9)
We could keep the c-number expectation value $`\varphi (t)`$ non-zero in general, which would lead to the general Gaussian state with displaced origin that we considered in some earlier papers . Since we do not require it in our free field de Sitter calculation, we specialize to the case where $`\varphi (t)=\dot{\varphi }(t)=0`$ in the following. With this restriction the field operator can be written simply as
$$\varphi (t)=a\psi (t)+a^{}\psi ^{}(t).$$
(B10)
We make use of the freedom in the $`\alpha `$ and $`\beta `$ coefficients to require
$`aa`$ $`=`$ $`a^{}a^{}=0,`$ (B11)
$`aa^{}`$ $`=`$ $`a^{}a+1={\displaystyle \frac{\sigma +1}{2}},`$ (B12)
with no loss of generality. The positive real parameter $`\sigma 1`$ has the interpretation $`\sigma =1+2n`$, with $`n`$ the average number of “particles” in the $`a`$ basis. Thus we need the three parameters $`\theta `$, $`\delta `$, and $`\sigma `$ to specify the general semi-classical (coherent) state of the field.
Using the definitions above it now follows that the average energy in this state is
$$E=H=\frac{\omega }{2}\sigma (1+2|\beta |^2).$$
(B13)
Note that we have not insisted that the state be an eigenstate of the number operator $`a^{}a`$ so $`n`$ is only the average number of particles in the state, which can take on any non-negative value and need not be an integer. If $`n>0`$ and the state is a Gaussian, consistent with all connected correlations vanishing except the two-point function, then it turns out that the state is necessarily a mixed state .
We are interested in the nature of the “vacuum” state in the limit $`\omega 0`$, i.e. when our single degree of freedom becomes that of a free particle. In this limit we can no longer retain the complex oscillating mode function basis since their normalization diverges due to the $`(2\omega )^{\frac{1}{2}}`$ factor in (B5). However, the mode equation $`\ddot{\psi }=0`$ clearly possesses the regular real solutions $`u=1`$ and $`v=t`$, which we can obtain from the complex oscillatory solutions by taking appropriate linear combinations in the limit of vanishing frequency, i.e.
$`u(t)`$ $``$ $`\underset{\omega 0}{lim}\left\{\sqrt{{\displaystyle \frac{\omega }{2}}}(f+f^{})\right\}=1,`$ (B14)
$`v(t)`$ $``$ $`i\underset{\omega 0}{lim}\left\{{\displaystyle \frac{1}{\sqrt{2\omega }}}(ff^{})\right\}=t.`$ (B15)
These definitions are analogous to Eq. (96) of the text, with $`\omega `$ replacing $`32\nu `$.
The general linear combination of modes can then be rewritten in the form
$$\psi (t)=\frac{(\alpha +\beta )}{\sqrt{2\omega }}ui\frac{\sqrt{\omega }}{2}(\alpha \beta )v=Bu+AvAt+B,$$
(B16)
where the quantities
$`A`$ $``$ $`i\underset{\omega 0}{lim}\left\{\sqrt{{\displaystyle \frac{\omega }{2}}}(\alpha \beta )\right\},`$ (B17)
$`B`$ $``$ $`\underset{\omega 0}{lim}\left\{{\displaystyle \frac{1}{\sqrt{2\omega }}}(\alpha +\beta )\right\},`$ (B18)
are analogous to those defined by (98) of the text. They also satisfy
$$A^{}BB^{}A=i(|\alpha |^2|\beta |^2)=i.$$
(B19)
We can also define the time-independent Hermitian operators
$`Q`$ $``$ $`{\displaystyle \frac{1}{\sqrt{2\omega }}}\left[(\alpha +\beta )a+(\alpha +\beta )^{}a^{}\right]Ba+B^{}a^{},`$ (B20)
$`P`$ $``$ $`i\sqrt{{\displaystyle \frac{\omega }{2}}}\left[(\alpha \beta )a(\alpha \beta )^{}a^{}\right]Aa+A^{}a^{},`$ (B21)
analogous to (102) in the de Sitter case. They obey the canonical commutation relations
$$[Q,P]=i,$$
(B22)
so that in the limit $`\omega 0`$
$$\varphi (t)u(t)Q+v(t)P=Q+tP,$$
(B23)
and in the same limit
$$E=\frac{1}{2}P^2=\frac{\sigma }{4}\underset{\omega 0}{lim}\left\{\omega |\alpha \beta |^2\right\}=\frac{\sigma }{2}|A|^2.$$
(B24)
Since $`\sigma `$ is just an overall factor and we are interested in pure vacuum-like states, we can set $`\sigma =1`$.
We see that depending on exactly how we take the limit $`\omega 0`$, we can get many different values for the energy. If we tried to keep the pure positive frequency solution $`\psi =f`$, i.e. $`\alpha =1`$ and $`\beta =0`$, then the energy would vanish in the limit $`\omega 0`$. However, by computing the correlation function, $`\varphi (t)\varphi (t^{})`$, we would find that the correlator diverges in this limit. In our simple example it is clear why. The ordinary ground state vacuum wave function of the simple harmonic oscillator is a Gaussian with a width proportional to $`\omega ^{\frac{1}{2}}`$. Hence $`Q^2\omega ^1`$ in this state, which becomes non-normalizable in the limit, unless we were to put the particle in a box of finite volume. But this option is unavailable to us if the “particle” is a mode of a field and the range of the field variable $`\varphi `$ is $`(\mathrm{},\mathrm{})`$. Thus, we must reject the positive frequency harmonic oscillator vacuum in the limit $`\omega 0`$ since it is a non-normalizable state.
In order for the state to remain normalizable and the mode functions bounded, the quantities $`A`$ and $`B`$ must remain finite in the limit of vanishing $`\omega `$. Inspection of the definitions (B18) shows that this requires that the original Bogoliubov coefficients $`|\alpha |`$ and $`|\beta |`$ go to infinity. In fact, using (B9) in the definitions of $`A`$ and $`B`$ a simple calculation shows that the Bogoliubov parameters $`\theta `$ and $`\delta `$ must behave like
$`\theta `$ $``$ $`{\displaystyle \frac{1}{2}}\mathrm{log}\left({\displaystyle \frac{\omega }{2}}\right)+\mathrm{log}|A|,`$ (B25)
$`\delta `$ $``$ $`\pi +\omega T,`$ (B26)
as $`\omega 0`$, so that
$`A`$ $`=`$ $`i|A|,\mathrm{and}`$ (B27)
$`B`$ $`=`$ $`{\displaystyle \frac{1}{2|A|}}+{\displaystyle \frac{i}{2}}|A|T,`$ (B28)
remain finite in that limit. Thus, the two complex numbers $`A`$ and $`B`$ obeying (B19) and the state they characterize can be specified by the two real numbers $`|A|0`$ and $`T`$, up to an unobservable overall phase. Comparison of these expressions with the corresponding equations (4.7) and (4.13) in the paper by Allen and Folacci shows that our $`A`$ here is equivalent to Allen and Folacci’s $`2A`$. This is because of a factor of 2 difference in the normalization of their time dependent mode $`v(t)`$.
Since $`\theta \mathrm{}`$, in a sense the normalizable “vacuum” state with finite $`P^2=|A|^2`$ and finite $`Q^2=|B|^2`$ is infinitely far from the usual ground state vacuum of the harmonic oscillator, and is not an eigenstate of the free particle Hamiltonian at $`\omega =0`$. Instead, it is just a Gaussian state centered at $`Q=0`$ and $`P=0`$ with normalized wave function
$$\mathrm{\Psi }(q)=\frac{1}{(2\pi |B|^2)^{\frac{1}{4}}}\mathrm{exp}\left(\frac{|A|}{2B^{}}q^2\right),$$
(B29)
in the position representation where $`Q\mathrm{\Psi }=q\mathrm{\Psi }`$ and $`P\mathrm{\Psi }=i\frac{\mathrm{d}}{\mathrm{d}q}\mathrm{\Psi }`$. This state is time reversal invariant if and only if $`T=0`$, in which case $`B=\frac{1}{2|A|}`$ becomes real. On the other hand, an eigenstate of the free Hamiltonian $`P^2/2`$ is necessarily non-normalizable since diagonalizing $`P`$ requires infinite spread in $`Q`$. These remarks carry over equally well to the $`k=1`$ mode in the de Sitter case.
There is no reason to diagonalize $`H`$ in the field theory case as Allen and Folacci do in their Eqs. (4.26) to (4.28) , since the energy-momentum tensor is time dependent in general, and back-reaction has been neglected in our simple calculations. Diagonalizing just the matter Hamiltonian without taking into account the metric is no better than a leading order mean field approximation to the much more difficult full problem with metric fluctuations taken into account. Given that what we are doing is only a semi-classical approximation to this full problem in any case, what makes the most sense from our perspective is to take well-defined normalizable initial adiabatic states for the scalar field and allow them to evolve as they will. Our initial state trial wave functional (B29) then remains a bounded Gaussian functional, consistent with semi-classical mean field methods (although it may well spread over time depending on the dynamics).
As should be clear, the non-normalizable positive frequency vacuum with $`\alpha =1`$ and $`\beta =0`$ is analogous to the Bunch-Davies vacuum, while the normalizable state related to it by an infinite Bogoliubov transformation (but only a finite shift in zero point energy) is analogous to the Allen-Folacci states parameterized by $`A`$ and $`B`$ or $`|A|`$ and $`T`$. However, whereas the energy of the normalizable state is finite and arbitrary, depending on the arbitrary $`|A|^2`$ coefficient in (B24), and the non-normalizable positive frequency state has zero energy in this simple harmonic oscillator example, the converse is true at late times in de Sitter space. This is due to the different kinematics and form of the energy-momentum tensor of the scalar field in the de Sitter case. Comparison of the energy (B24) of the harmonic oscillator example and the corresponding $`k=1`$ late time energy density in the de Sitter case with $`m=0`$ and $`\xi `$ small, viz.
$$\epsilon _1\frac{1}{2a^2}\left(\frac{\mathrm{d}\varphi _1}{\mathrm{d}\eta }\right)^2+3\xi \varphi _1^2=\frac{1}{4\pi ^2a^2}\left(\frac{\mathrm{d}v}{\mathrm{d}\eta }\right)^2P^2+\frac{3\xi }{2\pi ^2}(Q+vP)^2,$$
(B30)
shows that the difference is the redshift factor of $`a^2`$ multiplying the infrared finite $`P^2`$ in the de Sitter case. Thus, whereas this can take on the finite state dependent value $`|A|^2`$ in the harmonic oscillator example, its asymptotic late time value vanishes in de Sitter space when multiplied by $`a^2`$, and becomes state independent in any infrared finite state. Further, in the simple harmonic oscillator the usual ground state becomes an eigenstate of $`P`$ and has zero energy as $`\omega 0`$, since $`\omega ^2(Q+vP)^2\omega 0`$ drops out of the energy in this limit, whereas in the de Sitter case, since $`\xi `$ goes to zero only linearly with $`32\nu `$, the Bunch-Davies expectation value receives a finite contribution from the $`k=1`$ mode from the infinite spread in $`\xi Q^2`$ which has no $`a^2`$ redshift factor in (B30).
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# Cooling of a small sample of Bose atoms with accidental degeneracy
## I Introduction
### A Bose–Einstein condensation
The observation of effects related to the quantum statistical properties of weakly interacting gases of atoms has become in the last decades a major challenge of atomic physics. Thus, a large part of the theoretical and experimental research has been focused during the $`90`$’s on cooling atoms confined in traps at relatively high densities . These studies have led to the remarkable experimental realisation of a Bose-Einstein condensation (BEC) in Rubidium , and sodium vapors . Evidence of BEC in a Lithium gas with attractive interactions has been also reported . These remarkable achievements, have opened fascinating possibilities and applications, such as the developement of a coherent source of atoms, or atom laser.
The theoretical description of a system of ultracold bosonic or fermionic atoms is particularly convenient in the framework of second quantization. Quantum field theories of cold atoms , originally developed in a condensed–matter context, have been used in the diagnostics of a Bose-Einstein condensate , and in nonlinear atomic optics.
Several theoretical works in the recent years were more directly aimed to the dynamics of the cooling processes, dynamics of the possible phase transitions and the formation of quantum collective states. Those works concern both collisional cooling mechanisms, such as evaporation or symphatetic cooling , or laser cooling mechanisms, such as sideband cooling , Raman cooling , and dark state cooling . The latter processes allow to reach temperatures below the photon recoil energy $`E_R`$ \[equal to $`(\mathrm{}k)^2/(2M)`$, where $`k`$ is the laser wavevector and $`M`$ the atomic mass\]. One expects that the system might then enter a collective quantum state (such as Bose-Einstein condensate, or some analogue of it). In particular, under such conditions one hopes to realize also a coherent source of atoms.
### B Quantum Master Equation
In general, the quantum dynamics of a system of cold atoms is a very complex many-body problem. Some of the above mentioned processes may be analyzed using quantum Boltzmann equations . Starting from $`1994`$ a more general method based on a quantum master equation (QME) description has been developed. In particular, a QME describing the dynamics of a small sample of laser cooled atoms in a harmonic microtrap has been proposed and analyzed . The quantum statistical nature of the atoms is reflected in the dynamics of the cooling process. In the case when the trapping potentials for the atoms in the ground and excited electronic states are different such QME might lead to multistability and generalized Bose-Einstein distributions . The QME has the advantage that it permits to study atom number fluctuations in each of the trap levels, and thus provides a more complete description of the cooling process. In particular, it may be used, in principle, to describe the dynamics of condensate formation. We have also derived a QME for symphatetic cooling and analyzed the possibility of achieving the condensation of a system of light particles which interact with a reservoir of heavier particles. In the context of nonlinear atomic optics a self-consistent Born-Markov-Hartree-Fock master equation has been derived and analyzed to study the spontaneous emission effects on atomic solitons .
Since that early works, the theory of QME for many–body systems has been very strongly developed. In particular, C. Gardiner, P. Zoller and co–workers developed in a series of papers the QME and Quantum Kinetic Boltzmann approach to describe the dynamics of the evaporative cooling. We have also extended the theory of the collective laser cooling to much more realistic situations, avoiding the microtrap assumption, and working beyond the so–called Lamb–Dicke limit, in which the trap is of the size of the laser wavelength . In such analysis, we avoided the reabsorption problem by working in the so–called Festina Lente regime, in which the spontaneous emission rate is smaller than the trap frequency.
One should stress that, apart from the area of quantum optics, master equations for quantum Bose or Fermi gases have been used in statistical physics . However, the master equation is usually postulated there starting from general statistical requirements. It describes the approach towards the thermal equilibrium described by the Bose-Einstein or, correspondingly, Fermi-Dirac distributions, it fullfills detailed balance conditions, and sometimes it conserves some order parameters. The dynamics that it generates might have some universal properties, but does not have a direct physical interpretation in terms of interaction with specific energy reservoirs.
In most of the quantum optical examples, the master equation is derived via the elimination of the “bath” degrees of freedom starting from a more general theory that describes a very specyfic physical situation. The eliminated “bath” has a direct physical interpretation - it consists of photons, colliding atoms etc. Each of the jumps between the states of the system described by the QME usually corresponds to a very well defined physical process of photon emission, absorption, atom–atom collision etc. .
The QME approach, although valid in general, can only be used for practical calculations when the QME can be reduced to a set of kinetic equations, while the density matrix to a diagonal form. Such reduction is not always possible, and very often requires first the choice of an appropriate basis in the Hilbert space of the states of the system. In particular, the bare states of the system, that are the eigenstates of the system Hamiltonian in isolation from the “bath”, are not necessarily the right ones. For instance, in the process of cooling of an ideal bosonic gas confined in a microtrap, the reduction in the basis of bare states is possible only if additional assumptions are made. These statements seem a little surprising, since they hold even for arbitrary weak system-“bath” interactions, i.e. a situation in which the quantum Boltzmann equations should be valid . Note, however, that the validity of the quantum Boltzmann equations usually requires assumptions concerning quantum ergodicity, which simply do not hold in the above mentioned situations. The point here is that an ideal bosonic gas in a harmonic trap has plenty of degenerated bare energy levels. The quantum coherences between those levels (i.e. non-diagonal elements of the density matrix) can survive for very long times, and contribute significantly to the dynamics.
### C Degeneracy in many–body systems
Let us enumerate by $`\stackrel{}{m}`$ the eigenstates of a single atom Hamiltonian in the rotationally symmetric harmonic trap of frequency $`\omega `$, where $`\stackrel{}{m}`$ is a natural number in one dimension, a pair of natural numbers in 2D, a triple in 3D etc. When we consider an ensemble of $`N`$ atoms, the states of such an ideal gas can be written in the Fock representation as $`|n_\stackrel{}{0},n_\stackrel{}{1},\mathrm{}`$, where $`n_\stackrel{}{m}`$ denote the occupation numbers of the corresponding $`\stackrel{}{m}`$-th eigenstate. For noninteracting atoms there are two kinds of degeneracies in such a system. First, a degeneracy of energy levels due to rotational invariance; that is for the states for which the sum of $`n_\stackrel{}{m}`$’s with a fixed sum of the components of $`\stackrel{}{m}`$, is fixed itself. Obviously, such degeneracies are not present in 1D. We shall not discuss them here. Second, there exists an accidental degeneracy, due to the particular symmetry of the harmonic potential . This degeneracy occurs even in the case of 1D: for instance for the states $`|0,2,0,\mathrm{}`$ and $`|1,0,1,0,\mathrm{}`$. Here, the state with two atoms in the first energy level has an energy $`2\times \mathrm{}\omega `$, which is equal to the energy of the state with one atom in the ground level and another atom in the second excited level ($`1\times 0\mathrm{}\omega +1\times 2\mathrm{}\omega `$). Both kinds of degeneracies are lifted up if one considers anisotropic trap with anharmonic energy levels. If one then assumes that the resulting energy level shifts are larger than cooling rates, one can evoke standard secular arguments to reduce the master equation to a diagonal form in the basis of the bare ideal gas states . It is precisely the subject of the present paper to study the situation in which such reduction is not possible, i.e. when the effects of accidental degeneracy dominate the dynamics of the system.
One could argue that such problem is purely academic since in real physical system atom-atom interaction will always lift the accidental degeneracy. The whole point is, however, that as long as the interactions are not too strong, the system will still exhibit effects of quasi-degenracy. This will be the case when energy differences between the quasi-degenerated levels will be small in comparison to $`\mathrm{}\omega `$. Such condition is realised in not too dense systems, i.e. the system containing not too many atoms. The idealized theory presented in this paper is formulated for arbitrary number of atoms $`N`$, but in practice it applies to the situation when $`N`$ is not too large (see Section VII). Nevertheless, in view of the complexity of the problem, it is in our opinion reasonable to treat the ideal case of non-interacting atoms in order to get insight into more realistic cases. That said, let us note at this point that in the last years the external modification of the s–wave scattering length (which dominates the atom–atom collisions at low energies), has been theoretically investigated in different scenarios , and also experimentaly demonstrated by employing the so–called Feshbach resonances . Remarkably, very recent experimental results show that a modification of the scattering length to very small values is experimentally feasible , opening the fascinating possibility to acheive a quasi–ideal bosonic gas, as that studied in the present paper.
### D Content of the paper
The paper concerns thus the problem of cooling of an ideal Bose gas, or more precisely speaking a sample consisting of $`N`$ atoms, in a perfectly harmonic microtrap. It should be noted that this problem is quite general. Atomic traps, although frequently anisotropic (see for example Ref. ), can be designed to be, with a very good accuracy, harmonic. Moreover, even though in the small atomic samples atom-atom interactions will lift the exact degeneracy of energy levels, the system will remain quasi-degenerated. We expect that a cooled atomic sample in such a harmonic microtrap will exhibit the effects of accidental quasi-degeneracy regardless of the method used for its cooling In order to stress this general aspect of our study, we adopt here partially the statistical physics approach, and derive a master equation using a phenomenological model of the “bath”. In particular, our “bath” may represent one of the following two reservoirs: B1) a photon reservoir in the case of laser cooling ; B2) an atomic reservoir consisting of other atoms in the case of symphatetic cooling . In the case B2) the bath atoms are trapped in a larger trap than the system atoms. Such situation (proposed and discussed in ) could be realised if a small, say far-off-resonance dipole trap for system atoms was located inside a larger magnetic trap for the bath atoms .
Given that the resulting master equations for all these reservoirs have the same structure, the qualitative behavior for the cooling dynamics given by our specific model is quite general. The reason why we have chosen such a model for the bath is that it has the adventage that one can derive analytical formulas for the transition rates between different levels. We also stress that the mathematical treatment developed in this paper can be easily generalized to other physical situations, in particular for ultra–cold trapped polarised Fermi gases, which due to the suppression of the $`s`$–wave scattering induced by the Pauli principle, can be in an excelent approximation as ideal gases .
This paper is organized as follows: In Section II we introduce the model, describing separately the system in a trap, the atomic bath, and the system–bath interactions. This Section is very much analogous to Section II of Ref. , but is formulated for a different model of the bath. In Section III we derive the master equation governing the dynamics of the system, under Born and Markov approximations. This equation is further analyzed in Section IV in the microtrap limit in terms of the Lamb-Dicke (LD) expansion, i.e. a systematic expansion in a small parameter $`\eta =ka`$, where $`a`$ is the size of the trap ground state wavefunction, whereas $`k`$ is a typical wavevector of the bath quanta that is relevant for the cooling process. In Section V we discuss explicitely the breakdown of ergodicity due to the accidental degeneracy that occurs on a fast time scale in the lowest order of the LD expansion. Section VI is devoted to the discussion of the restoration of ergodicity on a slower time scale due to higher orders of LD expansion, whereas Section VII contains our conclusions. The paper also contains three appendices. Appendix A contains some useful formulas of the operator algebra used in the paper, and describes the structure of the Hilbert space. Appendix B describes the details of the construction of some of the multiple vacua that appear due to the accidental deneracy, whereas Appendix C presents matrix elements of the second quantized operators used in the paper. These elements are used for some calculations regarding the cooling rates. Finally, in Appendix D we present a proof of the decay of coherences on the slow time scale.
## II Description of the model
We consider a system “A” of bosonic particles that are confined in a trap, and interact with a bath bosonic particles “B”. We assume that the particles “B” are practically not affected by their interactions with the system “A”, so that “B” can be regarded as a phenomenological reservoir for “A”. The coupling to the bath represents the influence of some externally controlled cooling mechanism (laser cooling, symphatetic cooling, etc.) on the system “A”. The reservoir is assumed to be in thermal equilibrium at some given temperature. In this Section we introduce the Hamiltonian for the system “A”, the bath “B”, and for their mutual interactions. The formalism is developed for the case of $`d`$ spatial dimensions.
### A Description of the system
The system “A” is an ideal gas of $`N`$ bosonic atoms of mass $`M_A`$ confined by an isotropic harmonic potential in $`d`$ dimensions. In a second quantized form, and in the Fock representation the Hamiltonian describing the system can be written in the form
$$H_A=\underset{\stackrel{}{n}}{}\mathrm{}\omega (n_x+n_y+\mathrm{})a_\stackrel{}{n}^{}a_\stackrel{}{n},$$
(1)
where $`\omega `$ is the trap frequency, $`\stackrel{}{n}=(n_x,n_y,\mathrm{})`$ with $`n_x,n_y,\mathrm{}=0,1,2,\mathrm{}`$, and $`a_\stackrel{}{n}^{}`$ and $`a_\stackrel{}{n}`$ are creation and annihilation operators of particles in the $`\stackrel{}{n}`$–th level of the harmonic potential, respectively. Since in the present paper we are interested in the effects of the accidental degeneracy, we neglect the contribution of the atom–atom interactions to the total Hamiltonian (for the discussion of its role see , and Section VII).
### B Bath
Similarly as in Ref. , we assume that the system “B” contains a practically infinite number of bosons of mass $`M_B`$ embedded in a practically infinite volume, with finite density $`n_B`$. The free Hamiltonian for the bath “B” of particles in a second quantized form is
$$H_B=𝑑\stackrel{}{k}ϵ(\stackrel{}{k})b(\stackrel{}{k})^{}b(\stackrel{}{k}).$$
(2)
Here, $`\stackrel{}{k}`$ is a wavevector in a $`d`$–dimensional space. The function $`ϵ(\stackrel{}{k})`$ represents the dispersion relation of the bath particles. For instance, for quasi-free particles it reads
$$ϵ(\stackrel{}{k})=\frac{(\mathrm{}k)^2}{2M_B},$$
(3)
where $`M_B`$ is the effective mass of the bath “quanta”. In the following we shall use Eq. (3) but the theory is easily generalized to arbitrary shapes of the dispersion relation. In particular, we shall also consider massless bath quanta with a photonic-like linear dispersion relation. Note that such dispersion relation is in fact appropriate for both types of bath (B1), and (B2) mentioned above. This is obvious in the case of laser cooling, but is also true for collisional cooling schemes, since no mass is produced or lost in the collision processes. The operators $`b(\stackrel{}{k})^{}`$ and $`b(\stackrel{}{k})`$ are creation and annihilation operators of bath particles corresponding to plane wave states with momentum $`\stackrel{}{k}`$. In the case of laser cooling (B1) $`\stackrel{}{k}`$ corresponds to a photon momentum, whereas in the case (B2) $`\stackrel{}{k}`$ is rather a momentum transfer associated with the system atom-bath atom collision. The operators $`b(\stackrel{}{k})^{}`$ and $`b(\stackrel{}{k})`$ fulfill the usual commutation relations
$`[b(\stackrel{}{k}),b(\stackrel{}{k}^{})]`$ $`=`$ $`[b(\stackrel{}{k})^{},b(\stackrel{}{k}^{})^{}]=0,`$ (4)
$`[b(\stackrel{}{k}),b(\stackrel{}{k}^{})^{}]`$ $`=`$ $`\delta ^{(d)}(\stackrel{}{k}\stackrel{}{k}^{}).`$ (5)
In thermal equilibrium, the density operator describing the state of the bath $`\rho _B`$ corresponds to the usual Bose–Einstein distribution (BED) . In this situation, we have
$`b(\stackrel{}{k})b(\stackrel{}{k}^{})`$ $`=`$ $`b(\stackrel{}{k})^{}b(\stackrel{}{k}^{})^{}=0,`$ (7)
$`b(\stackrel{}{k})^{}b(\stackrel{}{k}^{})`$ $`=`$ $`n(\stackrel{}{k})\delta ^{(d)}(\stackrel{}{k}\stackrel{}{k}^{}),`$ (8)
where $`n(\stackrel{}{k})`$ is related to the number of particles with wavevector $`\stackrel{}{k}`$, and is given by
$$n(\stackrel{}{k})=\frac{ze^{\beta ϵ(\stackrel{}{k})}}{1ze^{\beta ϵ(\stackrel{}{k})}}.$$
(9)
In the above expression, $`\beta =1/(k_BT)`$ is the inverse temperature, and $`z=\mathrm{exp}(\beta \mu )`$ is the fugacity, while $`\mu `$ denotes the chemical potential. Note that $`\mu =0`$, $`z=1`$ for massless quanta, and for massive particles below the condensation point. Note also that both $`n(\stackrel{}{k})`$ and $`ϵ(\stackrel{}{k})`$ only depend on $`k|\stackrel{}{k}|`$. Particle and energy densities are connected with these quantities by the relations
$`n_B`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^d}}{\displaystyle 𝑑\stackrel{}{k}n(\stackrel{}{k})},`$ (11)
$`ϵ_B`$ $`=`$ $`{\displaystyle \frac{1}{(2\pi )^d}}{\displaystyle 𝑑\stackrel{}{k}n(\stackrel{}{k})ϵ(\stackrel{}{k})},`$ (12)
respectively.
### C Interactions
Within the present model the interactions between the particles and the bath describe the annihilation of an atom followed by its inmediate recreation, accompanied by absorption or emission of a single bath quantum. We have chosen such a model for the interactions since it is the simplest one that retains the effect produced by the accidental degeneracy in the cooling process. In any case, it may be regarded as a phenomenological interaction that may be due to various physical mechanisms (system atom-bath atom collisions, laser-atom interactions, etc.). Similarly to the case of atom-atom collisions we employ here an analogue of the shape–independent approximation to write down the corresponding atom-bath quantum interactions. We assume that these interactions are local (i.e. have a zero range) in the spatial representation. Mathematically, this approximation means that the wave functions of both kinds of particles do not change significantly over distances characterizing interparticle potential in the relevant energy range. In the Fock representation the interaction Hamiltonian is given by
$$H_{AB}=\underset{\stackrel{}{n},\stackrel{}{n}^{}}{}𝑑\stackrel{}{k}\gamma _{\stackrel{}{n},\stackrel{}{n}^{}}(\stackrel{}{k})a_\stackrel{}{n}^{}a_\stackrel{}{n}^{}b^{}(\stackrel{}{k}),$$
(13)
where
$$\gamma _{\stackrel{}{n},\stackrel{}{n}^{}}(\stackrel{}{k})=\frac{\kappa }{(2\pi )^d}𝑑\stackrel{}{x}\psi _\stackrel{}{n}(\stackrel{}{x})^{}\psi _\stackrel{}{n}^{}(\stackrel{}{x})e^{i\stackrel{}{k}\stackrel{}{x}},$$
(14)
$`\psi _\stackrel{}{n}(\stackrel{}{x})`$ is the wavefunction corresponding to the $`\stackrel{}{n}`$–th level of the harmonic oscillator, and $`\kappa `$ is a coupling constant. We have chosen the form of the interaction Hamiltonian (13), because of its simplicity. We stress, however, that the qualitative (and to some extend quantitative) results of the paper do not depend on the particular choice of $`H_{AB}`$.
Without loss of generality, we can exclude from the integration over $`\stackrel{}{k}`$ in (13) the value $`\stackrel{}{k}=0`$. This is clear since
$`{\displaystyle \underset{\stackrel{}{n},\stackrel{}{n}^{}}{}}\gamma _{\stackrel{}{n},\stackrel{}{n}^{}}(0)a_\stackrel{}{n}^{}a_\stackrel{}{n}^{}b^{}(0)`$ (15)
$`\left({\displaystyle \underset{\stackrel{}{n}}{}}a_\stackrel{}{n}^{}a_\stackrel{}{n}\right)\left(b^{}(0)\right),`$ (16)
is a constant shift operator of the zero momentum mode, proportional to the number of particles in the system “A”. One can always perform a unitary shift transformation of $`b(0)`$, and $`b^{}(0)`$ that cancels the term (16) and its hermitian conjugate in the Hamiltonian. Obviously, such transformation modifies the BED for the bath quanta with zero momentum, but the latter modification has no relevance for the transformed system-bath interactions in which the coupling to the bath zero mode is absent. Hence, from now on, in the integrals over $`\stackrel{}{k}`$ it will be implicitely assumed that $`\stackrel{}{k}0`$. On the other hand, since in the next Section we are going to make a rotating wave approximation (RWA) in the master equation derived from the Hamiltonian (13), we reduce Eq. (13) (as in Ref. ) to the form
$$H_{AB}=H_0+\underset{\alpha =1}{\overset{\mathrm{}}{}}(H_\alpha +H_\alpha ^{}).$$
(17)
Here, $`H_0`$ contains the part of $`H_{AB}`$ given in (13) in which the sum is extended over values with $`_{s=x,y,\mathrm{}}(n_sn_s^{})=0`$. $`H_\alpha `$ contains the part of $`H_{AB}`$ proportional to the bath annihilation operators $`b(\stackrel{}{k})`$, with the sum over $`\stackrel{}{n}`$, $`\stackrel{}{n}^{}`$ extended over the values for which $`_{s=x,y,\mathrm{}}(n_sn_s^{})=\alpha `$.
The QME that we shall derive in the next Section will contain a Hamiltonian part (describing energy level shifts), and non-Hamiltonian part describing dissipative decay processes. It is worth mentioning that due to the RWA the Hamiltonian part of the master equation will not be generally correct . However, as in Ref. , we will be only interested in the dissipative part of the master equation, which is correctly described under the mentioned RWA, provided the trap frequency $`\omega `$ is larger than the cooling rates. The latter assumption will be made all over the present paper.
## III Derivation of the master equation
The master equation for the above defined model can be derived following well–stablished procedures in the field of quantum optics , analogously to those discussed in Ref. . We first move to an interaction picture defined by the unitary operator $`\mathrm{exp}[i(H_A+H_B)t]`$. In this picture, the density operator $`\stackrel{~}{\rho }`$ describing system–plus–bath degrees of freeedom fulfills the following equation:
$$\frac{d\stackrel{~}{\rho }(t)}{dt}=\frac{i}{\mathrm{}}[\stackrel{~}{H}_{AB}(t),\stackrel{~}{\rho }(t)],$$
(18)
where the tilde indicates that the operators are expressed in the interaction picture. We integrate formally this equation, and substitute back into (18). Next, we define the reduced density operator for system “A”, $`\rho _A=\mathrm{Tr}_B(\rho )`$, where Tr<sub>B</sub> stands for the trace over the bath states, and make use of the fact that Tr$`{}_{B}{}^{}\{[H_{AB}(0),\rho (0)]\}=0`$, since we assume that the density operator for the bath $`\rho _B(0)`$ is diagonal in the Fock basis (with respect to $`H_B`$), whereas $`H_{AB}`$ does not contain any diagonal matrix elements \[the reader should recall that we have extracted the terms with $`\stackrel{}{k}=0`$ in $`H_{AB}`$, see Eq. (13)\]. We obtain the following equation:
$`{\displaystyle \frac{d\stackrel{~}{\rho }_A(t)}{dt}}=`$ (19)
$`{\displaystyle \frac{1}{\mathrm{}^2}}{\displaystyle _0^t}𝑑\tau \mathrm{Tr}_B\{[\stackrel{~}{H}_{AB}(t),[\stackrel{~}{H}_{AB}(t\tau ),\stackrel{~}{\rho }(t\tau )]]\}.`$ (20)
In the following step we perform Born–Markov approximation. For this approximation we have to assume that the correlation time $`\tau _c`$ of the reservoir is much shorter than the typical time over which $`\stackrel{~}{\rho }_A(t)`$ changes, i.e. the cooling time . The Born–Markov approximation is also related to the fact that the bath quanta are practically not affected by their interactions with the system; this allows us to write $`\stackrel{~}{\rho }(t\tau )=\stackrel{~}{\rho }_A(t\tau )\rho _B(0)`$. From the technical point of view, the correlation time $`\tau _c`$ can be defined as a time for which the integrand of (20) practically vanishes. For specific physical models of the bath it can be directly evaluated (see, for instance, Ref. ). The cooling time, on the other hand, depends on the physics of the interactions between atoms and bath quanta. It is controllable, and thus is assumed to be the longest time scale of the problem. In this case, we can safely substitute $`\stackrel{~}{\rho }_A(t\tau )`$ by $`\stackrel{~}{\rho }_A(t)`$ in the integral (20), and extend the upper limit of the integral to infinity.
In the next steps we make use of Eq. (17) and perform the RWA, i.e. neglect terms rotating at multiples of the trap frequency. Again, this approximation is based on the assumption that trap frequencies are large in comparison to the cooling rates.
Finally, taking into account the bath properties (II B), and coming back to the Schrödinger picture we obtain the following master equation:
$$\frac{d\rho _A}{dt}=\frac{i}{\mathrm{}}[H_A+H_{AA}^{},\rho _A]+\rho _A,$$
(21)
where $`H_{AA}^{}`$ is the Hamiltonian term produced by the elimination of the bath in the master equation. Physically, this term accounts for the energy level shifts due to the effective interaction between system particles via their interactions with the bath quanta. $`H_{AA}^{}`$ may also include the original atom-atom interactions (provided they were present in the original Hamiltonian) in the spirit of independent rates approximation . Similarly as in the case of the atom–atom collisions, all of the shifts caused by $`H_{AA}^{}`$ may be neglected in some situations depending on the specyfic model of the bath, the size of the trap, and the number of atoms in the system (for details see ). We shall omit them in the following, and come back to the discussion of their role in Section VII. We shall therefore postulate a QME restricting our attention to the Liouvillian $``$ that describes the cooling process,
$$=\underset{\alpha 0}{\overset{\mathrm{}}{}}_\alpha ,$$
(22)
where
$`_\alpha \rho _A`$ $`=`$ $`{\displaystyle \underset{\stackrel{}{n},\stackrel{}{n}^{},\stackrel{}{m},\stackrel{}{m}^{}}{}}\mathrm{\Gamma }_{\stackrel{}{n},\stackrel{}{n}^{}}^{\stackrel{}{m},\stackrel{}{m}^{}}(2a_\stackrel{}{m}^{}a_\stackrel{}{m}^{}\rho _Aa_\stackrel{}{n}^{}a_\stackrel{}{n}^{}`$ (23)
$``$ $`a_\stackrel{}{n}^{}a_\stackrel{}{n}^{}a_\stackrel{}{m}^{}a_\stackrel{}{m}^{}\rho _A\rho _Aa_\stackrel{}{n}^{}a_\stackrel{}{n}^{}a_\stackrel{}{m}^{}a_\stackrel{}{m}^{}).`$ (24)
The sum in this expression is extended to $`\stackrel{}{n},\stackrel{}{n}^{},\stackrel{}{m},\stackrel{}{m}^{}`$ fulfilling
$$\underset{s=x,y\mathrm{}}{}(n_sn_s^{})=\alpha ,\underset{s=x,y\mathrm{}}{}(m_sm_s^{})=\alpha .$$
(25)
Liouvillian (23) acounts for transitions of particles from one level of the harmonic oscillator to another, experiencing a change in the energy of $`\alpha \mathrm{}\omega `$, and a corresponding absorption or emission of a bath quantum. Thus, the term with $`\alpha =0`$ that conserves the energy does not enter (22) since we have excluded the bath quanta with zero energy from the interaction (see Sec. II). The terms with $`\alpha >0`$ ($`\alpha <0`$), on the other hand, describe processes increasing (decreasing) the energy. These transitions are characterized by
$`\mathrm{\Gamma }_{\stackrel{}{n},\stackrel{}{n}^{}}^{\stackrel{}{m},\stackrel{}{m}^{}}`$ $`=`$ $`{\displaystyle \frac{\pi }{\mathrm{}}}{\displaystyle 𝑑\stackrel{}{k}\gamma _{\stackrel{}{n},\stackrel{}{n}^{}}(\stackrel{}{k})\gamma _{\stackrel{}{m},\stackrel{}{m}^{}}^{}(\stackrel{}{k})}`$ (27)
$`\times [n(\stackrel{}{k})+1]\delta [ϵ(k)\mathrm{}\omega \alpha ],`$
for $`\alpha `$ positive, and by
$`\mathrm{\Gamma }_{\stackrel{}{n},\stackrel{}{n}^{}}^{\stackrel{}{m},\stackrel{}{m}^{}}`$ $`=`$ $`{\displaystyle \frac{\pi }{\mathrm{}}}{\displaystyle 𝑑\stackrel{}{k}\gamma _{\stackrel{}{n},\stackrel{}{n}^{}}(\stackrel{}{k})\gamma _{\stackrel{}{m},\stackrel{}{m}^{}}^{}(\stackrel{}{k})}`$ (29)
$`\times n(\stackrel{}{k})\delta [ϵ(k)\mathrm{}\omega |\alpha |],`$
for $`\alpha `$ negative.
We stress here the fact that for any kind of interactions between the atoms and the bath, the corresponding Liouvillian has the same form as in (22) and (23), provided in each interaction act one atom is annihilated and another created. This is the reason why our results of the following sections can be extended to other kind of interactions. Note, however, that in such cases the coefficients (27) and (29) may have a much more complicated form (see, for example, Ref. ).
## IV Lamb-Dicke expansion in 1D
In this Section we perform the Lamb-Dicke (LD) expansion of the master equation (22). From now on we shall concentrate on the one dimensional case and skip the vector notation.
The LD expansion is valid in the situation when the bath quanta relevant for the cooling process have momenta $`\stackrel{}{k}`$ much smaller than the inverse of the size of the trap, $`a`$. Their corresponding wave functions \[$`\mathrm{exp}(\pm ikx)`$\] vary slowly on a scale of $`a`$, and can be then expanded in Taylor series around $`x0`$. Since, according to Eqs. (27) and (29) the relevant bath quanta have energies $`ϵ(k)=\mathrm{}\omega \alpha `$, their corresponding momenta may be determined from the dispersion relation. For example, for the case of massive free particles characterized by the dispersion relation (3), the validity of LD at relatively low temperatures requires that
$$\left(\frac{2M_B\omega |\alpha |}{\mathrm{}}\right)^{1/2}\left(\frac{\mathrm{}}{2M_A\omega }\right)^{1/2}=\left(\frac{M_B|\alpha |}{M_A}\right)^{1/2}<1,$$
(30)
i.e. the effective mass of the bath particles must be much smaller than that of the system atoms . It is easy to find analogous conditions corresponding to other forms of the dispertion relation, including the case of massless bath quanta. In case of the laser cooling (B1) at low temperatures the condition is simply that $`ka=2\pi a/\lambda <1`$, where $`k`$ ($`\lambda `$) is the laser photon momentum (wavelength).
In the case of symphatetic cooling (B2) the condition is $`\mathrm{\Delta }ka<1`$, where $`\mathrm{\Delta }k`$ is the typical momentum transfer in a collision act. In such collision the kinetic energy of the incoming bath atom is of the order of $`1/\beta `$, and its correponding momentum of the order of $`2\pi /\lambda _B`$, where $`\lambda _B`$ denotes the thermal de Broglie wavelength. If the bath was cool, $`\beta \mathrm{}\omega 1`$, the momentum transfer would typically be equal to the final momentum of the bath atom. The validity of LD expansion would then require that the same condition as (30) is fulfilled, with $`M_B`$ denoting the real mass of the bath atoms, and that $`2\pi a/\lambda _B<(M_B/M_A)^{1/2}<1`$.
In the LD regime we can expand the rates (27), (29). Denoting by
$`k_\alpha `$ $`=\left({\displaystyle \frac{2M_B\omega |\alpha |}{\mathrm{}}}\right)^{1/2},`$ (31)
$`\eta _\alpha `$ $`=k_\alpha a\sqrt{\alpha }\eta ,`$ (32)
we obtain for $`n`$, $`n^{}`$, $`m`$, $`m^{}`$ fulfilling Eq. (25)
$`\mathrm{\Gamma }_{n,n^{}}^{m,m^{}}`$ $`=`$ $`{\displaystyle \frac{2\pi }{\mathrm{}}}\gamma _{n,n^{}}(k_\alpha )\gamma _{m,m^{}}^{}(k_\alpha )`$ (33)
$`\times `$ $`[n(k_\alpha )+(1+\mathrm{sign}(\alpha ))/2]{\displaystyle \frac{M_B}{\mathrm{}^2k_\alpha }},`$ (34)
and
$`\gamma _{n,n^{}}(k_\alpha )`$ $`=`$ $`{\displaystyle \frac{\kappa }{2\pi }}[\delta _{n,n^{}}i\eta _\alpha (\sqrt{n+1}\delta _{n,n^{}1}+\sqrt{n}\delta _{n,n^{}+1})`$ (35)
$``$ $`{\displaystyle \frac{\eta _\alpha ^2}{2}}(\sqrt{(n+2)(n+1)}\delta _{n,n^{}2}+(2n+1)\delta _{n,n^{}}`$ (36)
$`+`$ $`\sqrt{n(n1)}\delta _{n,n^{}+2})+\mathrm{}]`$ (37)
The master equation can be now rewritten using the above expressions in the interaction picture with respect to $`H_A`$, and upon neglection of $`H_{AA}^{}`$. It takes then the following form
$$\dot{\rho }=[^{(0)}+_1^{(1)}+_2^{(1)}+O(\eta ^6)+\mathrm{}]\rho .$$
(38)
The “zeroth” order term is actually of the order of $`\eta ^2`$ and has the form
$`^{(0)}\rho `$ $`=\mathrm{\Gamma }_{}\left[2A\rho A^{}A^{}A\rho \rho A^{}A\right]`$ (40)
$`+\mathrm{\Gamma }_{}\left[2A^{}\rho AAA^{}\rho \rho AA^{}\right];`$
where
$`\mathrm{\Gamma }_{}`$ $`=`$ $`N{\displaystyle \frac{\kappa ^2}{(2\pi )^2}}{\displaystyle \frac{2\pi }{\mathrm{}}}\eta ^2\left[n(k_1)+1\right]\left({\displaystyle \frac{M_B}{\mathrm{}^2k_1}}\right),`$ (42)
$`\mathrm{\Gamma }_{}`$ $`=`$ $`N{\displaystyle \frac{\kappa ^2}{(2\pi )^2}}{\displaystyle \frac{2\pi }{\mathrm{}}}\eta ^2\left[n(k_1)\right]\left({\displaystyle \frac{M_B}{\mathrm{}^2k_1}}\right),`$ (43)
whereas
$`A`$ $`={\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\sqrt{n+1}a_n^{}a_{n+1},`$ (45)
$`A^{}`$ $`={\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\sqrt{n+1}a_{n+1}^{}a_n.`$ (46)
Note that quite generally: a) both cooling rates $`\mathrm{\Gamma }_{}`$, $`\mathrm{\Gamma }_{}`$ are collective (i.e. proportional to $`N`$); b) their ratio is
$$\frac{\mathrm{\Gamma }_{}}{\mathrm{\Gamma }_{}}=\frac{n(k_1)}{n(k_1)+1}=ze^{\beta \mathrm{}\omega }.$$
(47)
Similarly, the higher order terms (of order $`\eta ^4`$) are
$`_1^{(1)}\rho `$ $`=\mathrm{\Gamma }_{}\eta ^2[2(A\rho C^{}+C\rho A^{})`$ (48)
$``$ $`(C^{}A\rho +A^{}C\rho )(\rho C^{}A+\rho A^{}C)]`$ (49)
$``$ $`\mathrm{\Gamma }_{}\eta ^2[2(A^{}\rho C+C^{}\rho A)`$ (50)
$``$ $`(CA^{}\rho +AC^{}\rho )(\rho CA^{}+\rho AC^{})],`$ (51)
with
$`C=`$ $`{\displaystyle \frac{1}{2\sqrt{N}}}{\displaystyle \underset{n=0}{}}(n+1)^{3/2}a_n^{}a_{n+1},`$ (53)
$`C^{}=`$ $`{\displaystyle \frac{1}{2\sqrt{N}}}{\displaystyle \underset{n=0}{}}(n+1)^{3/2}a_{n+1}^{}a_n,`$ (54)
and
$`_2^{(1)}\rho `$ $`=\mathrm{\Gamma }_{}^{(1)}\left[2B\rho B^{}B^{}B\rho \rho B^{}B\right]`$ (56)
$`+\mathrm{\Gamma }_{}^{(1)}\left[2B^{}\rho BBB^{}\rho \rho BB^{}\right];`$
where
$`\mathrm{\Gamma }_{}^{(1)}=N{\displaystyle \frac{\kappa ^2}{(2\pi )^2}}{\displaystyle \frac{2\pi }{\mathrm{}}}{\displaystyle \frac{\eta _2^4}{4}}\left[n(k_2)+1\right]\left({\displaystyle \frac{M_B}{\mathrm{}^2k_2}}\right),`$ (58)
$`\mathrm{\Gamma }_{}^{(1)}=N{\displaystyle \frac{\kappa ^2}{(2\pi )^2}}{\displaystyle \frac{2\pi }{\mathrm{}}}{\displaystyle \frac{\eta _2^4}{4}}\left[n(k_2)\right]\left({\displaystyle \frac{M_B}{\mathrm{}^2k_2}}\right),`$ (59)
whereas
$`B`$ $`={\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\sqrt{(n+1)(n+2)}a_n^{}a_{n+2},`$ (61)
$`B^{}`$ $`={\displaystyle \frac{1}{\sqrt{N}}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\sqrt{(n+1)(n+2)}a_{n+2}^{}a_n.`$ (62)
Note that similarly as in Eq. (47) the rates fulfill
$$\frac{\mathrm{\Gamma }_{}^{(1)}}{\mathrm{\Gamma }_{}^{(1)}}=\frac{n(k_2)}{n(k_2)+1}=ze^{2\beta \mathrm{}\omega }.$$
(63)
Finally, Eq. (38) contains terms of higher orders $`\eta ^6,\eta ^8`$, etc.
## V Accidental degeneracy and the breakdown of ergodicity
In this Section we discuss the dynamics governed by the lowest order term in the Lamb-Dicke expansion, i.e. by the equation
$$\dot{\rho }=^{(0)}\rho .$$
(64)
First we observe that
$$[A,A^{}]=1,$$
(65)
i.e. the operators $`A`$, and $`A^{}`$ represent an abstract harmonic oscillator Weyl–Heisenberg algebra, whereas the Liouville-von Neuman superoperator (40) describes interaction of that harmonic oscillator with an effective thermal bath (see, for instance ). The inverse temperature of the bath is given by
$$\mathrm{}\omega \beta _e=\mathrm{log}\left(\frac{\mathrm{\Gamma }_{}}{\mathrm{\Gamma }_{}}\right)=\mathrm{}\omega \beta \beta \mu .$$
(66)
Note that the effective temperature is never greater than the temperature of the bath, $`\beta _e\beta `$, since $`\mu 0`$. Alternatively, one may view the above equation, as if the effective temperature of the system was constant, but the frequency would change $`\mathrm{}\omega _e=\mathrm{}\omega \mu `$. This effect is a result of our definition of the system–bath interactions. Those interactions consist of absorption or emission of the bath quanta, and thus do not conserve the number of the bath particles. The Boltzmann exponent (66) must account for energy gain or loss due to bath particles creation/annihilation. The remaining question is how the abstract operator algebra is represented in the Fock-Hilbert space of our multiparticle system. To answer this question we first introduce the operator
$$D=BA^2/\sqrt{N},$$
(67)
which commutes with $`A`$, $`A^{}`$, $`B`$ and $`B^{}`$. We then observe that (see appendix A):
* For each value of $`l=0`$, or $`l=2,3,\mathrm{}`$ there exist $`n_N(l)`$, so called, vacuum states $`|0_{lsv}`$ that are annihilated by the “jump” operator $`A=_{n=0}^{\mathrm{}}\sqrt{n+1}g_n^{}g_{n+1}`$,
$`A|0_{lsv}=0,`$ (68)
where the index $`l`$ indicates the bare energy of the corresponding states (in units of $`\mathrm{}\omega `$),
$`{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}na_n^{}a_n|0_{l,s}=l|0_{l,s},`$ (69)
The state corresponding to $`l=0`$ contains all $`N`$ atoms in the ground state. The states with higher energy are constructed as linear combinations of degenerated Fock states. There is only one state of $`l=1`$ ($`N1`$ atoms in the ground state, and one atom in the first excited state), which is not annihilated by $`A`$; that is why there is no vacuum state with $`l=1`$. The index $`s`$ measures the number of excitations of $`D^{}`$ in the state $`|0_{lsv}(D^{})^s|0_{l2s0v}`$, and runs from 0 to $`E(l/2)`$, i.e. integer part of $`l/2`$. Each of the states $`|0_{lsv}`$ is an eigenstate of $`D^{}D`$ with an eigenvalue $`(4ls4s^22s+2Ns)/N`$. The index $`v`$ enumerates finally the states of energy $`l2s`$ annihilated by $`D`$. Denoting by $`m_N(l)`$ the number of states of energy $`l`$ annihilated by $`D`$, i.e. the states
$$D|0_{l0v}=0,$$
(70)
we obtain
$$n_N(l)=\underset{s=0}{\overset{E(l/2)}{}}m_N(l2s).$$
(71)
The construction of the vacuum states is descibed in the Appendix A. Some other examples of the vacua are constructed in Appendix B. Each of the vacuum states is a linear combination of the accidentally degenerated energy eigenstates. The number of accidentally degenerated states of the energy $`l`$ in the $`N`$ atom system, $`p_N(l)`$, is given by a solution of the partition problem of the number theory and is extravagantly large (c.f. $`p_N(l)O(\mathrm{exp}(\pi \sqrt{2l/3}))`$ for $`lN`$). The number of the vacua is given by $`n_N(l)=p_N(l)p_N(l1)`$. The very existence of these multiple vacuum states is thus a direct consequence and, at the same time, a signature of the accidental degeneracy .
* The vacuum states are orthonormal, $`0_{lsv}|0_{l^{}s^{}v^{}}=\delta _{ll^{}}\delta _{ss^{}}\delta _{vv^{}}`$.
* The Fock-Hilbert space of the system splits into an infinite number of Fock subspaces corresponding to each of the vacuum states. The Fock states in the $`(l,s,v)`$-th subspace are constructed as
$`|k_{lsv}={\displaystyle \frac{(A^{})^k}{\sqrt{k!}}}|0_{lsv},`$ (72)
with $`k=0,1,\mathrm{}`$. They are also mutually orthonormal, and are also eigenstates of the energy operator with the corresponding eigenvalue $`(l+k)`$. They are also highly degenerated (for $`k+l=k^{}+l^{}`$). In the following we shall use the notation $`w=(l,s,v)`$.
The dynamics exhibits in the Lamb-Dicke limit multiple timescales. On the fastest time scale of the order of $`\eta ^2`$ it is nonergodic, i.e. it does not mix the different $`w`$-subspaces. After a short time of the order of the inverse of $`\mathrm{\Gamma }_{}`$, $`\mathrm{\Gamma }_{}`$ all coherences between the $`|k_w`$ and $`|k_w^{}^{}`$ for $`kk^{}`$ vanish. Within each $`w`$-subspace the system approaches the ”thermal” equilibrium characterized by the density matrix diagonal in $`k`$, with zero off-diagonal elements, and the inverse temperature $`\beta _e`$. The dynamics, however, cannot be reduced to a Poisson jump process, i.e. a sequence of random jumps between the various $`|k_w`$ states with the transition probabilities governed by the detailed balance conditions characteristic for the thermal equilibrium. The reason is that coherences corresponding to $`ww^{}`$, but $`k=k^{}`$ do not vanish.
The quantum mechanical density matrix in this (formally stationary) limit becomes a sum of diagonal canonical ensembles in each of the subspaces, accompanied by a sum of non-diagonal terms connecting different $`w`$ and $`w^{}`$ for the same $`k`$’s,
$`\rho =\left[1\mathrm{exp}(\beta _e\mathrm{}\omega )\right]{\displaystyle \underset{w}{}}{\displaystyle \underset{k}{}}n_w|k_wk_w|e^{\beta _e\mathrm{}\omega k}`$ (73)
$`+\left[1\mathrm{exp}(\beta _e\mathrm{}\omega )\right]{\displaystyle \underset{ww^{}}{}}{\displaystyle \underset{k}{}}r_{ww^{}}|k_wk_w^{}|e^{\beta _e\mathrm{}\omega k}.`$ (74)
where the coefficients describe the populations of the corresponding subspaces, and are defined as
$$n_w=\mathrm{Tr}\left(P_w\rho (0)P_w\right),$$
(75)
and the cumulative coherences,
$$r_{ww^{}}=\underset{k}{}k_w|\rho (0)|k_w^{}.$$
(76)
Evidently, $`\rho `$ exhibits nonergodic effects and depends explicitely on the initial density operator. In the above expression $`P_w=_k|k_wk_w|`$ denotes a projection operator onto the $`w`$-th subspace. Note that the definition (75) implies that $`_wn_w=1`$.
We stress that in the Fock basis spanned by the states $`|k_w`$ the matrix $`\rho `$ is, in principle, not diagonal; moreover, in general, depending on the initial condition, it does contain coherences when represented in the Fock-Hilbert space corresponding to noninteracting atoms. That is the reason why the master equation in the latter basis cannot be reduced to the diagonal form in the Lamb-Dicke limit. That is why in order to arrive at such reduction additional assumptions have to be evoked that lift up the accidental degeneracy, as anharmonicity of energy levels caused by anharmonicity of the trap potential or interatomic interactions.
## VI Restoration of ergodicity
From the previous Section it is clear that the breakdown of ergodicity has only an approximate character since it is related to the lowest order dynamics in the LD expansion. It is natural to expect that Eq. (74) describes only a quasi-stationatry solution which is indeed reached on a time scale $`1/\mathrm{\Gamma }_{},1/\mathrm{\Gamma }_{}O(1/\eta ^2)`$, but undergoes further slow evolution on a time scale $`1/\mathrm{\Gamma }_{}^{(1)},1/\mathrm{\Gamma }_{}^{(1)}O(1/\eta ^4)`$. On this slower time scale the ergodicity should be (at least partially) restored, and the transitions between the different $`w`$-th subspaces should become possible. As we shall see below, that is indeed the case.
To this aim we consider the higher order corrections to the master equation (38) and treat them perturbatively within the standard adiabatic elimination method . We introduce the projection operator
$`𝒫\rho (t)={\displaystyle \underset{w}{}}n_w(t){\displaystyle \underset{k}{}}|k_wk_w|e^{\beta _e\mathrm{}\omega k}(1e^{\beta _e\mathrm{}\omega })`$ (77)
$`+{\displaystyle \underset{ww^{}}{}}r_{ww^{}}(t){\displaystyle \underset{k}{}}|k_wk_w^{}|e^{\beta _e\mathrm{}\omega k}(1e^{\beta _e\mathrm{}\omega }),`$ (78)
where the populations of the corresponding subspaces
$$n_w(t)=\mathrm{Tr}\left(P_w\rho (t)P_w\right),$$
(79)
and the cumulative coherences
$$r_{ww^{}}(t)=\underset{k}{}k_w|\rho (t)|k_w^{},$$
(80)
become now slowly varying functions of time.
Introducing the complementary projector $`𝒬=1𝒫`$ obtain
$`\dot{𝒫}\rho `$ $`=𝒫\left(_1^{(1)}+_2^{(1)}\right)𝒫\rho +𝒫\left(_1^{(1)}+_2^{(1)}\right)𝒬\rho ,`$ (81)
$`\dot{𝒬}\rho `$ $`=𝒬^{(0)}𝒬\rho +𝒬\left(_1^{(1)}+_2^{(1)}\right)𝒫\rho +\mathrm{}`$ (82)
In the latter Eq. (82) we have already employed the fact that $`\mathrm{\Gamma }^{(1)}/\mathrm{\Gamma }O(\eta ^2)`$, and neglected the higher order terms in $`\eta ^2`$. Moreover, adiabatic elimination of $`𝒬\rho `$ from Eq. (82) introduces corrections of the order $`\eta ^6`$ to Eq. (81) for $`𝒫\rho `$. Thus it may also be neglected in comparison to the leading term on the RHS of Eq. (81).
We obtain thus a relatively simple master equation
$$\dot{𝒫}\rho =𝒫\left(_1^{(1)}+_2^{(1)}\right)𝒫\rho $$
(83)
From this master equation, after elementary algebra we obtain a set of closed rate equations for the populations $`n_w`$ of the $`w`$-th subspaces, and the coherences $`r_{ww^{}}`$. These equations can be enormously simplified using the properties of the operators $`A`$, $`B`$, $`C`$, and their hermitian conjugates that are discussed in Appendix A. Amazingly, it is only the term $`_2^{(1)}`$ which contributes in this order of the LD expansion; moreover, the equations for populations and coherences decouple (see Appendix C for details),
$`\dot{n}_w`$ $`=`$ $`2\mathrm{\Gamma }_{}^{(1)}(1e^{\beta _e\mathrm{}\omega }){\displaystyle \underset{k}{}}\left({\displaystyle \underset{k^{},w^{}}{}}|k_w|B|k_w^{}^{}|^2e^{\beta _e\mathrm{}\omega k^{}}n_w^{}k_w|B^{}B|k_we^{\beta _e\mathrm{}\omega k}n_w\right)`$ (84)
$`+`$ $`2\mathrm{\Gamma }_{}^{(1)}(1e^{\beta _e\mathrm{}\omega }){\displaystyle \underset{k}{}}({\displaystyle \underset{k^{},w^{}}{}}|k_w|B^{}|k_w^{}^{}|^2e^{\beta _e\mathrm{}\omega k^{}}n_w^{}k_w|BB^{}|k_we^{\beta _e\mathrm{}\omega k}n_w),`$ (85)
$`\dot{r}_{ww^{}}`$ $`=`$ $`2\mathrm{\Gamma }_{}^{(1)}(1e^{\beta _e\mathrm{}\omega }){\displaystyle \underset{k}{}}({\displaystyle \underset{k^{},w_1,w_2}{}}k_w|B|k_{w_1}^{}k_{w_2}^{}|B^{}|k_w^{}e^{\beta _e\mathrm{}\omega k^{}}r_{w_1w_2}`$ (86)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{w_1}{}}[k_w|B^{}B|k_{w_1}r_{w_1w^{}}+k_{w_1}|B^{}B|k_w^{}r_{ww_1}]e^{\beta _e\mathrm{}\omega k})`$ (87)
$`+`$ $`2\mathrm{\Gamma }_{}^{(1)}(1e^{\beta _e\mathrm{}\omega }){\displaystyle \underset{k}{}}({\displaystyle \underset{k^{},w_1,w_2}{}}k_w|B^{}|k_{w_1}^{}k_{w_2}^{}|B|k_w^{}e^{\beta _e\mathrm{}\omega k^{}}r_{w_1w_2}`$ (88)
$``$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{w_1}{}}[k_w|BB^{}|k_{w_1}r_{w_1w^{}}+k_{w_1}|BB^{}|k_w^{}r_{ww_1}]e^{\beta _e\mathrm{}\omega k}),`$ (89)
The above equations are further reduced inserting the unity between $`B`$ and $`B^{}`$, and using the fact that $`A`$, and $`A^{}`$ by definition do not couple the different $`w`$-subspaces, whereas $`B`$ ($`B^{}`$) has the only relevant matrix elements between the different $`lsv`$-subspaces for the ($`l\pm 2`$, $`s\pm 1`$, $`v`$-th) and $`lsv`$-th subspace, for $`l=2,3,\mathrm{}`$ ($`l=0,2,3,\mathrm{}`$), $`s=0,\mathrm{},E(l/2)`$, and $`k=k^{}=0,1,\mathrm{}`$ (see Appendix C).
From the above considerations we obtain the final form of the equations
$`\dot{n}_w`$ $`=`$ $`2\mathrm{\Gamma }_{}^{(1)}\left[f_{w+2}^2n_{w+2}f_w^2n_w\right]`$ (90)
$`+`$ $`2\mathrm{\Gamma }_{}^{(1)}\left[f_w^2n_{w2}f_{w+2}^2n_w\right].`$ (91)
$`\dot{r}_{ww^{}}`$ $`=`$ $`2\mathrm{\Gamma }_{}^{(1)}\left[f_{w+2}f_{w^{}+2}r_{w+2w^{}+2}{\displaystyle \frac{1}{2}}\left(f_w^2+f_w^{}^2\right)r_{ww^{}}\right]`$ (92)
$`+`$ $`2\mathrm{\Gamma }_{}^{(1)}\left[f_wf_w^{}r_{w2w^{}2}{\displaystyle \frac{1}{2}}\left(f_{w+2}^2+f_{w^{}+2}\right)r_{ww^{}}\right],`$ (93)
where we have denoted $`w=(lsv)`$, $`w\pm 2=(l\pm 2,s\pm 1,v)`$, with $`f_w=0_{w2}|B|0_w`$. Note that $`f_w`$ can be assumed to be real and positive without loss of generality. As we see, both the populations $`n_w`$, and the coherences $`r_{ww^{}}`$ fulfill a closed set of rate equations. The explicit expressions for the matrix elements involved in the above formulae are derived in Appendix C.
The above equations provide the basic result of this paper. It shows that due to the presence of accidental degeneracy in the LDL the dynamics occurs essentially on several time scales. On the fastest time scale (governed by $`^{(0)}`$) the dynamics is nonergodic and consist in approach toward the canonical equilibrium states in each of the $`l`$-subspaces, accompanied by creation of quasi-equilibrium coherences between the states corresponding to the same $`k`$’s, but different $`w`$ and $`w^{}`$. On this scale the populations of each of the $`w`$-subspace, as well as the cumulative coherences $`r_{ww^{}}`$ may be regarded as constant, and therefore the values of these coherences, as well as the populations in each of the subspaces depend on the initial conditions. In the higher order of expansion (on the time scale $`1/\eta ^2`$ times longer), the mixing between different $`w`$-subspaces becomes possible. This mixing leads to a partial restoration of ergodicity. In the example considered here this restoration is not full, however, because as easily seen from Eq. (91) the $`w`$-subspaces with different $`v`$ still do not mix. The reader can easily check that the further mixing of the odd and even subspaces does take place in the next order of the Lamb-Dicke expansion (for instance due to term containing bi-products of the operators $`C`$, and $`C^{}`$).
The stationary state that is reached on the slower time scale is easy to find since Eqs. (91) fulfill the detailed balance conditions, whereas the cumulative coherences are damped to zero, as demonstrated in Appendix D. We obtain that
$$n_{l+2,s+1,v}=ze^{2\beta \mathrm{}\omega }n_{lsv},$$
(95)
so that
$`n_{2l+m,l,v}=e^{2\beta _e^{}\mathrm{}\omega l}n_{m0v},`$ (96)
with
$$\beta _e^{}=\beta \beta \mu /2\mathrm{}\omega .$$
(97)
The ratio of $`n_{m0v}`$, and $`n_{m^{}0v^{}}`$ remain undetermined in this order. Note that the reason why $`\beta _e^{}\beta _e`$ is that both temperatures correspond to the processes that involve single bath quantum absorbtion or emission, but different energy changes (by $`\mathrm{}\omega `$, or $`2\mathrm{}\omega `$, respectively). Indeed, it is elementary to check that the stationary diagonal matrix elements of the desity matrix are proportional to the corresponding Boltzmann factors,
$`k_{2l+m,l,v}|\rho |k_{2l+m,l,v}e^{\beta _e\mathrm{}\omega k2\beta _e^{}\mathrm{}\omega l}n_{m0v}`$ (98)
$`=k_{2l+m,l,v}|e^{\beta _e^{}\mathrm{}\omega \widehat{E}m+\beta \mu A^{}A/2}n_{m0v}|k_{2l+m,l,v}.`$ (99)
Finally, it is easy to check by substitution in the Liouvillian (22) that for $`\mu =0`$ the steady state to all orders in the LD expansion is precisely $`\rho e^{\beta H_A}`$, which is diagonal in the original basis. Obviously, the steady state solutions obtained in the first, and the second order of our expansion are compatible with such a steady state.
## VII Conclusions
In a series of papers we have studied in detail the quantum dynamics of bosonic and fermionic gases of cold atoms in traps in the absence of accidental degeneracy. We studied various cooling mechanism, and various limiting cases. In this paper we have presented the solution of the corresponding problem accounting for accidental degeneracy effects. We have studied interactions of a gas of trapped atoms with a heat bath in the Lamb-Dicke limit using the master equation approach. We have demonstrated that the system approaches an equilibrium on two (or more) distinguished time scales, and that the dynamics has the corresponding number of stages. At each stage a quasi-equilibrium state within appropriately determined subspaces of the Hilbert space is reached. This quasi-equilibrium corresponds to a canonical ensemble resticted to the appropriate subspace, and characterized by an effective temperature determined by the temperature and the chemical potential of the heat bath. In the next stage thermalization between the groups of subspaces occurs leading to another quasi-equilibrium in the larger subspaces, and so on.
We would like to stress once more that the problem considered in this paper is quite general. Atomic traps, although frequently anisotropic (see for example Ref. ), can be designed to be harmonic with a very good accuracy. A cooled atomic sample in such a harmonic microtrap will necessarily exhibit the effects of accidental degeneracy regardless of the method used for its cooling!
One may question the generality of our results, since we have used the Lamb-Dicke expansion, and at the same time neglected in this paper atom-atom collisions, as well as atom-atom interactions mediated by the coupling with the bath. Such processes (described by the Hamiltonian $`H_{AA}^{}`$, see Section III) lead evidently to shifts of the atomic energy levels, and will, in principle, lift up the accidental degeneracies. As we argued in Refs. , as long as the number of atoms in the microtrap is not too large, those shifts remain small and can be treated perturbatively. The system will then still exhibit the effects of accidental quasi-degeneracy. We stress that the theory developed in this article is valid for arbitrary numbers of atoms, and in particular it is for two atoms. Using far-off-resonance dipole traps , or loading atoms to a single minimum of a dark optical lattice (see Ref. (b)) it should be accesible experimentaly to confine several atoms in the trap of the size $`a0.10.05\mu `$m. That implies validity of the LD expansion for the laser colling case (see ). Similarly, one can use a cooled atomic gas close to, or below the condensation point as a bath in the symphatetic cooling case. In the conditions of the experiments of Refs. and that implies de Broglie wavelength of the order of $`\mu `$m, and thus validates the LD expansion. Using Bogoliubov-Hartree theory it is possible to estimate perturbatively that the energy level splittings in a ”band” of the quasi-degenerated states due to atom-atom collisions will be in such a case of the order of $`N(a_{sc}/a)/\sqrt{N_D}`$, where $`a_{sc}`$ is the scattering length of the system atoms, and $`N_D`$ is the number of levels in the band. Note that $`N_D`$ increases dramatically for higher excited levels. We see with $`a_{sc}=5`$ nm, our theory should work for $`N`$ up to $`20`$ even in the worst case when $`N_D=2`$. Note that the cooling of $`20`$ atoms to the ground state of a harmonic trap might be a very interesting task for the rapidly–developing research field of quantum informacion processing. Additionaly, we want to recall at this point that as pointed above, the external modification of the $`s`$–wave scattering length via Feshbach resonances has been demonstrated, been experimentally feasible the achievement of a quasi–ideal gas.
The main physical results of the paper are thus the following. We have been able to treat analytically the quantum dynamics of an ideal gas of $`N`$ atom in the LD limit. We have shown that the dynamics naturally splits into two parts: a fast part, during which coherences between the degenerated states are preserved, and a slower part, during which thermal equilibrium is achieved. Even though, the ideal case considered is not realistic (at least without external modification of the scattering length), we think i) that it provides a lot of insight into more realistic situations; ii) it is, to our knowledge, one of the extremely rare examples of soluble quantum dynamical problems in the area of statistical physics. The method that we developed, and results can be carried over to more realistic situations concerning cooling of small atomic samples ($``$20 atoms) in microtraps. Such situations are not far from the reach of present experiments. The calculations for such a case should follow exactly the lines described in this paper, with the only difference that the parts of the Liouville equations describing the atom-atom interactions that lift up the exact degeneracies should be included into the corrections to the $`^{(0)}`$. In another words, they should be treated just like the corrections to the lowest order term in the LD expansion have been treated in this paper. It is obvious, that as long as the splittings of the quasi-degenerated levels will remain small relative to $`\mathrm{}\omega `$, such realistic system will exhibit basic effects presented in this paper, i.e. step-wise dynamics on the two time scales.
The quantum dynamics of samples of cold atoms exhibits, in our opinion, an enormous reachness of interesting physical and mathematical phenomena, such as multistable, exotic stationary states, multistage dynamics etc. The present paper is another example of this reachness. Further studies are, however, required to get more understanding of this new physics, including for instance developement of other statistical physics tools (c), such as diffusion equations, hydrodynamic limits etc.
After this work was finished we have learned from T. Fischer, K. Vogel and W. Schleich that similar algebra to the one used by us appears in the problem of the cooling of a sample of bosons with a simple particle reservoir . We thank Yvan Castin, Jean Dalibard, Ralph Dum, T. Fischer, K. Vogel, and P. Zoller for enlighting discussions. We acknowledge the support of the Deutsche Forschungsgemeinschaft (SFB 407), ESF PESC Proposal BEC2000+, and TMR ERBXTCT96-0002.
## A The structure of the Fock-Hilbert space
The matrix elements of the operators $`B`$, and $`C`$ can be calculated using elegant algebraic methods. To this aim we first observe that the operators in question fulfill the commutation relations
$`[A,B]=`$ $`0,`$ (A2)
$`[A,C]=`$ $`B/(2\sqrt{N}),`$ (A3)
$`[A,B^{}]=`$ $`2A^{}/\sqrt{N},`$ (A4)
$`[A,C^{}]=`$ $`\widehat{E}/N+1/2,`$ (A5)
$`[B,B^{}]=`$ $`4\widehat{E}/N+2,`$ (A6)
$`[B,\widehat{E}]=`$ $`2B,`$ (A7)
$`[C,\widehat{E}]=`$ $`C,`$ (A8)
$`[A,\widehat{E}]=`$ $`A,`$ (A9)
with $`\widehat{E}=_{n=0}^{\mathrm{}}na_n^{}a_n`$ denoting the normalized energy operator.
It proves to be very useful to introduce the operator
$$D=BA^2/\sqrt{N}.$$
(A10)
This operator fulfills
$`[D,A]=`$ $`0,`$ (A12)
$`[D,A^{}]=`$ $`0,`$ (A13)
$`[D,D^{}]=`$ $`4(\widehat{E}A^{}A)/N+2(11/N),`$ (A14)
Since the operators $`A^{}A`$, and $`D^{}D`$ commute, it is useful to characterize the multiple vacua in terms of the eigenvalues of these two hermitian operators.
Let $`|0_{l0v}`$ denote the states that fulfill
$`A|0_{l0v}=0,`$ (A16)
$`D|0_{l0v}=0.`$ (A17)
There are $`m_N(l)`$ such states, and the index $`v`$ enumerates them. Note that the states
$$|0_{lsv}=(D^{})^s|0_{l2s,0,v}/(D^{})^s|0_{l2s,0,v},$$
(A18)
have energy $`l`$, are annihilated by $`A`$, and are eigenstates of $`D^{}D`$ with the eigenvalue
$$\underset{s^{}=0}{\overset{s1}{}}\left[4(l2s+2s^{})/N+2(11/N)\right].$$
(A19)
In the subsequent energy sectors we have thus the vacua: $`|0_{001}`$, $`|0_{211}`$, $`|0_{301}`$, $`|0_{421}`$, $`|0_{401}`$, $`|0_{511}`$, $`|0_{501}`$, $`|0_{631}`$, $`|0_{611}`$, $`|0_{601}`$, $`|0_{602}`$, etc.
The Fock-Hilbert space is spanned by the vectors
$$|k_{lsv}=\frac{(A^{})^k}{\sqrt{k!}}|0_{lsv}.$$
(A20)
## B Construction of vacuum states
We have seen in Appendix A that the vacuum states can be constructed by applying the operator $`D^{}`$ consecutively to the states $`|0_{l,0,v}`$. In this Appendix we present explicit construction of another family of the vacuum states that are annihilated by the operator $`A`$. In fact we consider a more general case with
$$A=\underset{n=0}{}A_na_n^{}a_{n+1}.$$
(B1)
Such defined operator reduces to the one given by Eq. (45) if we put $`A_n=\sqrt{n+1}`$. The vaccum states fulfill
$$A|0_l=0$$
(B2)
There is one obvious solution of the above equation which describes the global ground state
$$|0_0=|N,0,0,\mathrm{}.$$
(B3)
Apart from that, for $`l=2,3,\mathrm{}`$ we define
$`|0_l`$ $`={\displaystyle \underset{m=1}{\overset{l1}{}}}\alpha _m^la_{l+1m}^{}|Nm,m1,0,\mathrm{}`$ (B5)
$`+\alpha _l^l|Nl,l,0,\mathrm{}.`$
From the above definition it is clear, that different vacua are orthogonal. Applying $`A`$ to the above expression after elementary algebra we derive the recurrence formulas for the coefficients
$$\alpha _m^l=\frac{A_0}{A_{lm}}\sqrt{m}\sqrt{Nm}\alpha _{m+1}^l$$
(B6)
valid for $`m=1,l2`$, and
$$\alpha _{l1}^l=\frac{A_0}{A_1}\frac{\sqrt{l}\sqrt{Nl+1}}{\sqrt{l1}}\alpha _l^l.$$
(B7)
From the above expression it is easy to construct explicitely the corresponding vacuum states. The value of $`\alpha _l^l`$ can be conveniently chosen for normalisation of the states.
In the case considered in this paper ($`A_n=\sqrt{n+1}`$) the first few normalized vacuum states are:
$`|0_2={\displaystyle \frac{1}{\sqrt{N}}}(|N2,2,0,\mathrm{}`$ (B8)
$`\sqrt{N1}|N1,0,1,0,\mathrm{}),`$ (B9)
$`|0_3={\displaystyle \frac{\sqrt{8}}{\sqrt{N^2+3N2}}}(|N3,3,0,\mathrm{}`$ (B10)
$`{\displaystyle \frac{\sqrt{3(N2)}}{2}}|N2,1,1,0,\mathrm{}`$ (B11)
$`+{\displaystyle \frac{\sqrt{(N1)(N2)}}{2\sqrt{2}}}|N1,0,0,1,0,\mathrm{})`$ (B12)
$`|0_4={\displaystyle \frac{3}{\sqrt{N^35N^23N+21}}}(|N4,4,0,\mathrm{}`$ (B13)
$`{\displaystyle \frac{\sqrt{2(N3)}}{3}}|N3,2,1,0,\mathrm{}`$ (B14)
$`+{\displaystyle \frac{2\sqrt{(N2)(N3)}}{3}}|N2,1,0,1,0,\mathrm{}`$ (B15)
$`{\displaystyle \frac{\sqrt{(N1)(N2)(N3)}}{3}}|N1,0,0,0,1,0,\mathrm{})`$ (B16)
etc.
## C Calculation of the matrix elements
Let us first consider the operator B, and derive the explicit expressions for the martix elements
$$f_{lsv}=0_{l2s1v}|B|0_{lsv}=0_{l2s1v}|D|0_{lsv}$$
(C1)
that enter Eq. (91). Note that the coefficients $`f_{lsv}`$ are real, since the matrix elements of the operator $`A`$ in the non-interacting atom basis are real (see Appendix A). Moreover, without any loss of generality $`f_{lsv}`$’s may be assumed to be non-negative. In the following I will skip the index $`v`$ which is not affected by the dynamics.
Since we know that $`B|0_{ls}|0_{l2s1}`$ for $`l2s`$, we can write
$$B|0_{ls}=f_{ls}|0_{l2s1}.$$
(C2)
Similarly, using the commutation relation (A4) we can write
$$B^{}|0_{l2s1v}=f_{ls}|0_{ls}+\sqrt{2/N}|2_{l2,s},$$
(C3)
or
$$f_{ls}^2=0_{l2,s1}|BB^{}|0_{l2,s1}2/N.$$
(C4)
From the above expressions using the commutation relation (A6) we obtain
$$f_{l+2,+1}^2=f_{ls}^2+4l/N+22/N.$$
(C5)
The above recurrence can be easily solved yielding
$`f_{ls}^2`$ $`=(22/N)s+4((l2s)s+s(s1))N.`$ (C6)
since $`f_{l2s,0}=0`$.
In general, we may write
$$B|k_{ls}=f_{ls}|k_{l2,s1}+\sqrt{k(k1]/N}|(k2)_{ls}.$$
(C7)
The above formulae provide a very efficient method of calculating all of the matrix elements of the operators $`B`$ and $`B^{}`$.
It is a little more tedious to derive corresponding formulae for the operator $`C`$. From Eq. (A2) we obtain
$$AC|0_{ls}=\frac{1}{2\sqrt{N}}f_{ls}|0_{l2,s1},$$
(C8)
so that
$$C|0_{ls}=\frac{1}{2\sqrt{N}}f_{ls}|1_{l2,s1}+\underset{s^{}}{\overset{E(l1/2)}{}}w_{ls^{}}|0_{l1,s^{}}.$$
(C9)
The coefficients $`w_{ls^{}}=0_{l1,s1}|C|0_{ls}`$ may be also regarded to be real, and can be, for instance, determined directly from the definitions of the vacuum states.
In general, we can write
$`C|k_{ls}`$ $`=`$ $`{\displaystyle \frac{\sqrt{k+1}}{2\sqrt{N}}}f_{ls}|k+1_{l2,s1}`$ (C10)
$`+`$ $`{\displaystyle \underset{s^{}}{\overset{E(l1/2)}{}}}w_{ls^{}}|k_{l1,s^{}}+v_{ls}(k)|k1_{ls},`$ (C11)
with
$$v_{ls}(k+1)=\sqrt{\frac{k}{k+1}}v_{ls}(k)+\frac{1}{\sqrt{k+1}}\left(\frac{k+l}{N}+\frac{1}{2}\right),$$
(C12)
and $`v_{ls}(0)=0`$. The above formulae allow for very efficient calculations of the matrix element of the operators $`C`$ and $`C^{}`$, provided the states $`|0_{m0v}`$ are known.
The expression (C11) implies immediately that $`_1^{(1)}`$ does not contribute at all to the final equations (86) and (89). Since the matrix $`\rho `$ is diagonal in the $`k`$ index, whereas the operators $`A`$ and $`A^{}`$ change $`k`$ to $`k1`$, and $`k+1`$ respectively, only those parts of the operators $`C`$ and $`C^{}`$ that change $`k\pm 1`$ back to $`k`$ could contribute. It is evident from Eq. (C11), however, that these parts of $`C`$ and $`C^{}`$ do not change $`lsv`$. Therefore, their contributions to Eqs. (86) and (89) vanish identically, due to the trace-like sums over $`k`$ appearing on the right hand side.
Similar considerations show that there is no mixing of the populations and the cumulative coherences in Eqs. (86) and (89). Let us, for example, consider Eq. (86) for the populations $`n_{lsv}`$. As in the previous case, the contributions of the parts of the operators $`B`$ and $`B^{}`$ that do not change $`lsv`$ vanish identically, due to the trace-like sums over $`k`$ appearing on the right hand side. The parts of $`B`$ and $`B^{}`$ that change $`l`$ by $`2`$, and $`s`$ by $`1`$ do contribute, but they can only transform the parts of the density matrix proportional to $`|k_{lsv}k_{l^{}s^{}v^{}}|`$ into $`|k_{lsv}k_{l^{}s^{}v^{}}|`$, or $`|k_{l\pm 2,s\pm 1,v}k_{l^{}\pm 2,s^{}\pm 1,v^{}}|`$, i.e. they can only lead to couplings between the populations $`n_{lsv}`$ and $`n_{l\pm 2,s\pm 1,v}`$, or between the cumulative coherences $`r_{lsv,l^{}s^{}v^{}}`$ and $`r_{l\pm 2,s\pm 1,v,l^{}\pm 2,s^{}\pm 1,v^{}}`$.
## D Decay of coherences
In this Appendix we keep a single index $`w=(l,s,v)`$, and denote $`w\pm 2=(l\pm 2,s\pm 1,v)`$. In order to prove that the cumulative coherences decay to zero on the slow time scale, we rewrite Eq. (LABEL:glowna1) in the form
$$\dot{r}_{ww^{}}=\dot{r}_{ww^{}}^{DB}+\dot{r}_{ww^{}}^{NEG},$$
(D1)
where the first term
$`\dot{r}_{ww^{}}^{DB}=2\mathrm{\Gamma }_{}^{(1)}\left[f_{w+2}f_{w^{}+2}r_{w+2w^{}+2}f_wf_w^{}r_{ww^{}}\right]`$ (D2)
$`+2\mathrm{\Gamma }_{}^{(1)}\left[f_wf_w^{}r_{w2w^{}2}f_{w+2}f_{w^{}+2}r_{ww^{}}\right],`$ (D3)
corresponds to a set of kinetic equations with (positive) rates that fulfill detailed balance conditions. The matrix that enters the right hand side and generates the evolution is therefore evidently non-positively defined, and has exactly one eigenvector with zero eigenvalue, corresponding to the stationary solution of the Boltzmann-Gibbs form.
The second term in Eq. (D1)
$`\dot{r}_{ww^{}}^{NEG}=2\mathrm{\Gamma }_{}^{(1)}\left[{\displaystyle \frac{1}{2}}\left(f_w^2+f_w^{}^2f_wf_w^{}\right)r_{ww^{}}\right]`$ (D4)
$`2\mathrm{\Gamma }_{}^{(1)}\left[{\displaystyle \frac{1}{2}}\left(f_{w+2}^2+f_{w^{}+2}f_{w+2}f_{w^{}+2}\right)r_{ww^{}}\right],`$ (D5)
corresponds to a set of simple decay equations, with the rates which are evidently positive, since $`\frac{1}{2}(f_w^2+f_w^{}^2)f_wf_w^{}`$ is strictly greater than zero for $`ww^{}`$.
The full dynamics of the cumulative coherences is thus generated by the sum of the two martices, one of which is non-positively defined, and the other being strictly negatively defined. The sum itself must therefore be negatively defined, ergo it generates the decay to zero.
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# Drifting Pattern Domains in a Reaction-Diffusion System with Nonlocal Coupling
## Abstract
Drifting pattern domains (DPDs), i.e. moving localized patches of traveling waves embedded in a stationary (Turing) pattern background and vice versa, are observed in simulations of a reaction-diffusion model with nonlocal coupling. Within this model, a region of bistability between Turing patterns and traveling waves arises from a codimension-2 Turing-wave bifurcation (TWB). DPDs are found within that region in a substantial distance from the TWB. We investigated the dynamics of single interfaces between Turing and wave patterns. It is found that DPDs exist due to a locking of the interface velocities, which is imposed by the absence of space-time defects near these interfaces.
\- Introduction - Pattern forming processes in nonequilibrium systems can be classified according to the primary instability of the spatially homogeneous state. Ref. distinguishes three basic types of instabilities in unbounded systems: (i) spatially periodic and stationary in time, (ii) spatially periodic and oscillatory in time and (iii) spatially homogeneous and oscillatory in time. Within the reaction-diffusion literature, these instabilities are known as Turing, wave and Hopf bifurcation, respectively.
Many chemical and biological patterns are well captured by so called activator-inhibitor models describing the dynamics of two reacting and diffusing substances with two coupled partial differential equations. In such two component reaction-diffusion models only Turing and Hopf instabilities are possible. Recently, numerical investigations of chemical reaction-diffusion systems with three components and nonlocal coupling have yielded the occurrence of wave instabilities and the corresponding patterns. A universal description of patterns near these instabilities is achieved within the framework of amplitude equations .
Here, we study a simple FitzHugh-Nagumo model with inhibitory nonlocal coupling that is obtained as a limiting case of a three component reaction-diffusion system. It describes the interaction of an activator species with an inhibitor. For slow inhibitor diffusion (compared to the activator diffusion), the model exhibits wave instabilities, while, for fast inhibitor diffusion, Turing instabilities are found. The two instabilities occur simultaneously at a codimension-2 Turing-wave bifurcation (TWB). Such a situation has been found earlier within a model for binary convection and is a generalization of the well investigated Turing-Hopf instability in reaction-diffusion systems . Basic properties of a TWB have been studied theoretically in amplitude equations as well as experimentally in a one-dimensional gas-discharge system . In our model, we find a pattern previously unknown in
reaction-diffusion systems: drifting pattern domains (DPDs), i.e. localized patches of traveling waves embedded in a Turing background and vice versa (see Fig. 1). These patches have constant width and move (drift) with constant speed. As they drift along, maxima of the concentrations of activator and inhibitor are conserved beyond both boundaries of the DPD, where Turing and wave patterns are joined together. Consequently, formation of space-time defects by coalescence of maxima and minima is prevented. Similar patterns have been reported in a variety of hydrodynamical experimental systems, see e.g. and have been related to secondary instabilities (parity breaking) of stationary patterns . DPDs exist in a broad region of the para-
meter space, but appear only in a substantial distance to the onset of pattern formation. Their existence region is characterized by bistability between waves and Turing patterns. Near the boundary of DPD existence small DPDs containing only a single wave or Turing ,,cell” inside a background of the respective other state are encountered (see Figs. 1b,c). Outside the DPD existence region, the width of pattern domains shrinks or expands; these transient domains exhibit defects at the interfaces. Large DPDs are composed of two interfaces separating wave and Turing patterns. In this Letter we will study the dynamics of such interfaces. These interfaces typically select the wavenumber of the invading domain and form one parameter families characterized by the wavenumber of the invaded domain. Far away from the TWB, interfaces can exhibit a locking mechanism of their velocities due to the absence of defects. This locking implies that the interface velocity is fixed by the wavenumbers and frequency of the patterns that they separate. If both interfaces in a DPD are locked, they have to travel with equal speed. This mechanism is responsible for the existence of DPDs with constant width.
\- Model equations and linear stability - We study a variant of the FitzHugh-Nagumo equation supplemented by an inhibitory nonlocal coupling in the dynamics of the activator $`u`$
$`_tu`$ $`=`$ $`au+\beta u^2\alpha u^3bv+_x^2u`$ (2)
$`\mu {\displaystyle _{\mathrm{}}^+\mathrm{}}e^{\sigma |xx^{^{}}|}u(x^{^{}},t)𝑑x^{^{}}`$
$`_tv`$ $`=`$ $`cudv+\delta _x^2v.`$ (3)
Eqs. (1) represent a limiting case of a three variable model involving the activator $`u`$, the inhibitor $`v`$ and an additional fast inhibitor . Related three variable models have been introduced previously to describe pattern formation on sea shells and in cell biology as well as spot dynamics in gas discharges and concentration patterns in heterogeneous catalysis . Here, the emphasis is on the onset of pattern formation resulting from destabilization of a single homogeneous steady state. Eqs. (1) possess the trivial homogeneous fixed point $`𝐮_0\stackrel{\mathrm{def}}{=}(u_0,v_0)^T=(0,0)^T`$ for all parameter values. Here, we consider the regime where this fixed point is the only one present, i.e. $`a<bc/d+2\mu /\sigma `$ and consider perturbations $`e^{ikx\lambda (k)t}`$, where $`\lambda (k)=\chi (k)+i\omega (k)`$. The growth rates $`\lambda (k)`$ are given by the eigenvalues of the Jacobian. Linear stability analysis reveals that Eqs. (1) exhibit wave instabilities if the nonlocal coupling is of sufficiently long range $`\sigma <\sigma _c=(2\mu /(1+\delta ))^{1/3}`$.
In the following we vary the control parameters $`a`$ and $`\delta `$; the ,,driving force” $`a`$ represents the kinetics, whereas the ratio of diffusion coefficients $`\delta `$ describes the spatial coupling in the medium. All other parameters of Eqs. (1) have been fixed . For the wave bifurcation, the critical wavenumber $`k_W^c`$ and parameters $`a_W,\delta _W`$ are obtained from the condition $`\lambda (k_W^c)=\pm i\omega _0`$ where the perturbation with $`k_W^c`$ is the fastest growing mode with $`(k_W^c)^2=\sqrt{\frac{2\mu \sigma }{1+\delta }}\sigma ^2`$. Note, that for both $`\sigma =\sigma _C`$ and $`\sigma =0`$ (global coupling limit) the critical wavenumber is $`k_W^c=0`$, see Fig. 2a. Similarly, a competing Turing instability appears for a critical parameter $`a_T`$ with a wavenumber $`k_T^c`$, where the leading eigenvalue $`\lambda (k_T^c)=0`$. For large enough driving $`a`$, the wave instability appears for small $`\delta `$, while for large $`\delta `$ the Turing instability destabilizes the homogeneous state. For the chosen parameter values, the system exhibits a TWB point (see Fig. 3a and ). For the corresponding $`\lambda (k)`$, see Fig. 2b.
\- Weakly nonlinear analysis - Near the TWB, we can write $`𝐮\stackrel{\mathrm{def}}{=}(u,v)^T`$ as a perturbative expansion around $`𝐮_\mathrm{𝟎}`$ using a small parameter $`\epsilon `$, indicating the distance to the instability threshold: $`𝐮=𝐮_0+\epsilon 𝐮_1+\epsilon ^2𝐮_2+\epsilon ^3𝐮_3+\mathrm{}`$ and use the following multiple scale ansatz:
$`𝐮_1`$ $`=`$ $`[A(X,T__1,T__2)𝐔_Ae^{i(\omega _0t+k_W^cx)}+B(X,T__1,T__2)𝐔_Be^{i(\omega _0tk_W^cx)}`$
$`+(X,T__1,T__2)𝐔_{}e^{ik_T^cx}+c.c.]/2.`$
This leads to a set of coupled equations for the amplitudes $`A`$, $`B`$ and $``$ for left-, right-going waves and Turing pattern that depend on slow time and space variables. After reestablishing the original time and space variables and performing further $`\epsilon `$-independent scaling, one obtains:
$`_t`$ $`=`$ $`\eta ||^2+\xi _x^2\zeta (|A|^2+|B|^2)`$ (4)
$`_tA+c_g_xA`$ $`=`$ $`\rho A+(1+ic_1)_x^2A(1ic_3)|A|^2A`$ (6)
$`g(1ic_2)|B|^2A\nu (1i\kappa )||^2A`$
$`_tBc_g_xB`$ $`=`$ $`\rho B+(1+ic_1)_x^2B(1ic_3)|B|^2B`$ (8)
$`g(1ic_2)|A|^2B\nu (1i\kappa )||^2B.`$
For the detailed values of all coefficients, see . Note, that the nonlocal term of Eqs. (1) only enters into the diffusion coefficients of Eqs. (2) and does not give rise to a nonlocal term in Eqs. (2). Knowledge of the coefficients of Eqs. (2) allows analytical predictions of the pattern dynamics. Here, traveling waves are always preferred over standing waves ($`g>1`$, see ) and bistability between wave and Turing patterns is found ($`\nu \zeta >1`$). In this bistability region in parameter space (see Fig. 3a), a family of stable Turing patterns and two families of stable left- and right-traveling waves parametrised by their corresponding wavenumbers coexist. To get further insight, we take a closer look at the dynamics of single interfaces separating domains of Turing and wave patterns.
\- Interface dynamics. - With suitable initial conditions, a moving interface between Turing and wave patterns will be formed in simulations of Eqs. (1). We can distinguish two types of interfaces depending on whether the phase velocity of the waves points towards the interface or away from it. This classification is independent from the direction in which the interface is moving. In the following, we will call the first type inward-interfaces and the second outward-interfaces. Fig. 3b and 3c show examples of the latter type.
Near the TWB, we have studied general properties of such interfaces in amplitude equations (3) by counting arguments as well as by direct numerical simulations. Counting arguments are applied to ordinary differential equations obtained from a coherent structure ansatz in a comoving frame. We observe that, typically, the wavenumber of the invaded domain remains constant, while it adapts in the invading domain. In other words, the interface selects a particular wavenumber for the invading state, while the initial wavenumber of the invaded state is a free parameter. The velocity of the interface is a function of this parameter. For Turing patterns the selected wavenumber is always the critical one, i. e. $`k_T^{sel}=k_T^c`$, while for waves typically $`k_W^{sel}k_W^c`$ and therefore $`\omega ^{sel}\omega _0`$. This is valid for both inward- and outward-interfaces. Thus, we typically have two one-parameter families of interfaces for a given point in parameter space. These results are confirmed by numerical simulations of the nonlocal model (1) near the TWB.
Simulations of interfaces in Eqs. (1) far away from the TWB, show qualitatively similar behavior with respect to the selected wavenumbers. In addition, interfaces far away from the TWB may exhibit a locking mechanism of their velocities, which are fixed by the wavenumbers and frequencies of the Turing and wave domains. More specifically, the selected velocity is determined by the absence of defects at the interface. For geometrical reasons an interface without defects, that connects a wave state with wavenumber $`k_W`$ and frequency $`\omega `$ and a Turing state with $`k_T`$, has a speed $`|v_{lock}|=\omega /(k_Tk_W)`$. This velocity locking mechanism is found for both types of interfaces. Two examples of areas where locking occurs (locking tongues) in parameter space for inward- and outward-interfaces are given in Fig. 3a. A locked outward-interface is displayed in Fig. 3b. Outside the tongue the outward-interfaces display phase slips (see Fig. 3c). The area of the locking tongues depends only weakly on the front parameter $`k_W`$ for inward-interfaces and $`k_T`$ for outward-interfaces. Note, that the locking tongues open at a sub-
stantial distance from the TWB. Since the rapidly varying space and time scales have been factored out in the amplitude equations (2), locking tongues cannot be found therein. However on a line in the parameter space (see Fig. 3a), the velocity of interfaces in Eqs. (2) coincides with the velocity prescribed by the locking mechanism. The locking mechanism arises when the characteristic width of the interfaces is of the same order than the characteristic length scale of the patterns.
\- DPDs and their Phase Diagram. - Consider that a large DPD is composed of an inward-interface and an outward-interface, that practically do not interact (see e.g. Fig. 1a). If both interfaces exhibit no defects and are locked, their velocities have equal magnitude $`|v_{lock}|`$ but opposite signs. This ensures constant width and allows construction of DPDs of arbitrary size. Indeed, the region of existence of large DPDs starts to open where the locking tongues for both interface types begin to overlap (see Fig. 3a). This is the case for $`a5.7`$. Above that
value, DPDs spontaneously form from a variety of initial conditions. We have determined the parameter region, where they propagate with constant width and drift speed, from extensive simulations in systems with sizes $`L>400`$ and periodic boundary conditions. The results are shown in the phase diagram of Fig. 4.
We can distinguish three different subregions. In region B, DPDs of any size, with two locked interfaces traveling at the same speed, are found (see Fig. 1a). In region A, the inward-interface is no longer locked and its speed is smaller than $`|v_{lock}|`$. Therefore large domains of Turing (wave) patterns contract (expand) in size until only a stable DPDs containing a single Turing cell is left (see Fig. 1b). In region C, the outward-interface selects a $`k_W^{sel}`$ which would be unstable against Turing patterns in an infinite domain. Therefore, the wave domain forming the DPD is mostly replaced by a Turing pattern. However, small DPDs with a few wavelength of wave pattern are still encountered. At the outer boundary of region C, only DPDs with a single wave cell are found to be stable (see Fig. 1c).
\- Conclusion - We found a large variety of drifting pattern domains in a reaction-diffusion model with nonlocal coupling. Their ingredients include a bistability between wave and Turing patterns near a codimension-2 point as well as absence of defects at the interface. They exist as robust patterns only in a finite distance to the onset of pattern formation. Our results are not limited to the reaction-diffusion model studied here and should carry over to other physical systems with similar pattern forming instabilities. Altogether, DPDs and their constituting interfaces represent a generalization of simpler structures such as fronts and pulses in bistable reaction-diffusion systems, which do not simply combine two homogeneous states, but, instead, select their constituents from whole families of possible traveling or stationary periodic patterns.
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# On the origin of the difference between the runaway velocities of the OB-supergiant X-ray Binaries and the Be/X-ray Binaries
## 1 Introduction
A high-mass X-ray binary (HMXB) consists of a massive OB-type star and a compact X-ray source, a neutron star or a black hole. The X-ray source is powered by accretion of wind material, though in some systems mass transfer takes place through Roche-lobe overflow; the compact stars in the latter systems are surrounded by an accretion disk. Since wind accretion plays an important role, in practice only an OB supergiant or a Be-star companion have a strong enough stellar wind to result in observable X-ray emission. In a Be/X-ray binary the X-ray source is only observed when the neutron star moves through the dense Be-star disk at periatron passage. About 75% of the known HMXBs are Be/X-ray binaries, although this is a lower limit given their transient character.
Chevalier & Ilovaisky (1998) derived the proper motions for a sample of HMXBs from Hipparcos measurements. The four OB-supergiant HMXBs for which proper motions are available (0114+65, 0900-40 \[Vela X-1\], 1700-37 and Cyg X-1) have relatively large peculiar tangential velocities. Some corrections to the values given by these authors are needed (cf. Steele et al. 1998, Kaper et al. 1999). Taking these into account (Table 1) the mean peculiar velocity of these systems is $`42\pm 14`$$`\mathrm{kms}^1`$. It was already known that the OB-supergiant system of 1538-52 (QV Nor) has a peculiar radial velocity of about 90 $`\mathrm{kms}^1`$ with respect to its local standard of rest (Crampton et al. 1978; Gies & Bolton 1986; van Oijen 1989). For the 13 Be/X-ray binaries with measured proper motions Chevalier & Ilovaisky found peculiar tangential velocities ranging from $`v_{\mathrm{tr}}=3.3\pm 0.7`$ to $`21\pm 7.4`$$`\mathrm{kms}^1`$, with an average of $`v_{\mathrm{tr}}=11.4\pm 6.6`$$`\mathrm{kms}^1`$. Again, after corrections (see Sect. 2) and excluding the Oe systems XPer (0352+309) and V725Tau (0535+262), one finds for the genuine Be/X-ray binary a slightly higher value of $`v_{\mathrm{tr}}=15\pm 6`$$`\mathrm{kms}^1`$.
We would like to point out here that these mean values are in good agreement with the runaway velocities of these two types of systems predicted on the basis of simple “conservative” evolutionary models (van den Heuvel 1983, 1985, 1994; Habets 1985; van den Heuvel & Rappaport 1987) and even better agreement is obtained when mass is not conserved in the transfer process (Portegies Zwart 2000). The effect of sudden mass loss during the supernova explosion is taken into account and in a massive binary this is the dominant contribution to the runaway velocity; a random kick velocity of a few hundred $`\mathrm{kms}^1`$ imparted to the neutron star at birth (see e.g. Hartman 1997) has only a small effect, as the kick’s impulse has to be distributed over the entire massive ($`\stackrel{>}{}15`$$`\mathrm{M}_{}`$) system. (See Portegies Zwart & van den Heuvel 1999, for arguments in favor of kicks). Therefore, in first-order approximation, these kicks can be neglected in calculating the runaway velocities of HMXBs, but not in calculating their orbital eccentricities (see Sects. 3.4 and 3.5).
The aim of the present paper is to give a quantitative assessment of the above-mentioned conjectures. It should be noted here that five Be-star systems in the Be/X-ray binary sample studied by Chevalier & Ilovaisky (1998) are of spectral type B4Ve or later (masses $`6M_{}`$). The companions of these stars might be white dwarfs instead of neutron stars. Therefore, a supernova explosion is not necessarily the reason for their (excess) space velocity, which, in any case, is relatively small. It may be due to the typical random velocities observed in young stellar systems. Leaving these late-type Be/X-ray binaries out does not result in a significant change in the observed mean peculiar velocity of the Be-systems. Furthermore, there is some doubt concerning the use of the distances based on Hipparcos parallaxes of several of the other Be-systems, as these distances differ very much from the distances determined in other ways, e.g. by using reddening etc. (Steele et al. 1998). In Sect. 2 we therefore critically examine the distances and proper motions of all the systems with Be companions.
In Sects. 3.1 and 3.2 we present an analytical calculation of the expected runaway velocities and orbital eccentricities of typical OB-supergiant and Be HMXBs, on the basis of the standard evolutionary models for these systems, adopting conservative mass transfer during phases of mass exchange, and including the effects of stellar-wind mass loss for the OB-supergiant systems. In Sect. 4 we discuss the effect of non-conservative mass transfer on the runaway velocity and in Sect. 5.1 for the Be/X-ray binaries with known orbital eccentricities. We calculate which kick velocities should be imparted to the neutron stars of Be/X-ray binaries in order to produce their, on average, large orbital eccentricities (since the mass-loss effects alone cannot produce these). In Sect. 5.2, as an alternative, we compare the observed runaway velocities and orbital eccentricities of the Be/X-ray binaries with those expected on the basis of symmetric mass ejection and show that without kicks their combination of high orbital eccentricities and low space velocities cannot be explained. Our conclusions are summarized in Sect. 6.
## 2 The observed peculiar tangential velocities of HMXBs
The 4 OB-supergiant systems in the Hipparcos sample of Chevalier & Ilovaisky (1998) have distances larger than 1 kpc, which is too remote for a reliable parallax determination. For these systems they estimated the distances based on the spectral type, visual magnitude and reddening, and eventually the strength and velocity of interstellar absorption features, etc. After correcting for the peculiar solar motion and differential galactic rotation (see also Moffat et al. 1998) the Hipparcos proper motions result in the peculiar tangential velocities listed in Table1. Chevalier & Ilovaisky give a mean peculiar tangential velocity of $`v_{\mathrm{tr}}=41.5\pm 15`$$`\mathrm{kms}^1`$. We derive a similar value of $`42\pm 14`$$`\mathrm{kms}^1`$.
For the Be-systems, Chevalier & Ilovaisky use the Hipparcos parallaxes to determine the distances. For some systems this leads to very surprising results. In particular, Steele et al. (1998) point out that for the system of 0236+610 (LSA + $`61^{}303`$) the Hipparcos parallax leads to a ten times smaller distance than the distance estimated from the spectral type and reddening. These authors convincingly show that for this system the distance estimate based on the Hipparcos parallax cannot be correct; the distance of the system must be of order 1.8kpc instead of the 177pc determined from the Hipparcos parallax. Similarly, from a variety of criteria they find that for A0535+262 the distance must be $`>1.3`$kpc, instead of the 300pc determined from the Hipparcos parallax. Steele et al. point out that for both systems the Hipparcos parallaxes are smaller than 3 times their probable (measurement) error, and are therefore not reliable. In such a case one cannot reliably use the Hipparcos parallax to determine the distance. With the Hipparcos distances the OB-star companions of 0236+610 and 0535+262 would become highly underluminous for their spectral types, and would be very peculiar stars, as was already noticed by Chevalier & Ilovaisky (1998). On the other hand, using alternative distance criteria, their absolute luminosities become perfectly normal for their spectral types. This gives confidence that the latter distances are more reliable.
The systems including a Be star with spectral type later than B4V (mass $`6M_{}`$) may well have white dwarfs instead of neutron stars as companions (Portegies Zwart 1995). Therefore, their space velocity is not necessarily caused by a supernova explosion, which is the scenario we exploit in this paper. Excluding these systems, the observed mean peculiar velocity hardly changes ($`v_{\mathrm{tr}}=14.4\pm 6.6`$$`\mathrm{kms}^1`$ in stead of $`15\pm 6`$, excluding XPer and V725Tau), and since the nature of their compact companions is not known anyway (e.g. no X-ray pulsations observed which would identify the compact star as a neutron star), we decided to leave them in the calculation of the mean peculiar velocity. The peculiar tangential velocities of XPer (0352+309, O9III-IVe, 27$`\mathrm{kms}^1`$) and V725Tau (0535+262, O9.7IIe, 97$`\mathrm{kms}^1`$) are relatively high; their early spectral types suggest that they have masses comparable to those of the OB supergiants, so that, like the OB-supergiant systems, they would also originate from relatively massive binary systems. In Table1 we list the peculiar tangential velocity for each individual system and calculate the average for different subsamples. We left out the Bestar $`\gamma `$Cas, because its X-ray binary nature is not clear; furthermore, its X-ray spectrum is consistent with that of a white dwarf (Haberl 1995).
For the systems 0236+610, 0535+262, 1036-565 and 1145-619 the Hipparcos parallaxes yield absolute visual magnitudes very different from those expected on the basis of the OB-spectral types of the stars. In these cases, the Hipparcos parallax measurements are less than three times their probable errors and thus not reliable. For these stars we therefore used the distances determined from spectral type and reddening, which yield absolute visual magnitudes consistent with their spectral types.
We rederived the peculiar tangential velocities relative to the local restframe from the Hipparcos proper motions (cf. Kaper et al. 1999). Table1 lists the peculiar tangential velocity corrected for the peculiar solar motion and differential galactic rotation for three different distances ($`d/1.4`$, $`d`$, $`1.4d`$, following Gies & Bolton 1986). The uncertainty in distance (and thus in peculiar motion) is difficult to estimate; therefore, we calculated the space velocity for different values of the distance. The peculiar tangential velocities for the HMXBs discussed in Clark & Dolan (1999) are identical to ours for the OB-supergiant systems, though they find different values for the Be/X-ray binaries XPer ($`15\pm 3`$$`\mathrm{kms}^1`$, $`d=700`$pc), V725Tau ($`57\pm 14`$$`\mathrm{kms}^1`$, $`d=2`$kpc), and 1145-619 ($`17\pm 7`$$`\mathrm{kms}^1`$, $`d=510`$pc). Obviously, the precise values for the peculiar motion depend on the adopted model for the galactic rotation; we used the formalism employed in Comerón et al. (1998). For the OB-supergiant systems in the sample of Chevalier & Ilovaisky (1998) also the radial velocities are available from literature. This is not the case for the Be/X-ray binary systems. Therefore, we only consider the two components of the tangential velocity for the comparison of the kinematic properties of the two groups. The table shows that, leaving the two O-emission systems out, the Be/X-ray binaries have low space velocities: $`15\pm 6`$$`\mathrm{kms}^1`$.
## 3 Runaway velocities expected on the basis of models with conservative mass transfer and symmetric mass ejection.
### 3.1 Change of orbital period due to mass transfer
We only consider here so-called case B mass transfer since for the evolution of massive close binaries this is the dominant mode of mass transfer (cf. Paczynski 1971; van den Heuvel 1994, but see Wellstein & Langer 1999). In case B the mass transfer starts after the primary has terminated core-hydrogen burning, and before core-helium ignition. After the mass transfer in this case the remnant of the primary star is its helium core, while its entire hydrogen-rich envelope has been transferred to the secondary, which due to this became the more massive component of the system. There is a simple relation between the mass of the helium core $`M_{\mathrm{He}}`$ and that of its progenitor $`M_{}`$(see for example van der Linden 1982; Iben & Tutukov 1985). We adopt here the relation given by Iben & Tutukov (1985):
$$M_{\mathrm{He}}=0.058M_{}^{1.57},$$
(1)
which results in a fractional helium core mass $`p`$ given by:
$$p=M_{\mathrm{He}}/M_{}=0.058M_{}^{0.57}.$$
(2)
The change in orbital period of the system in case of conservative mass transfer (i.e.: conservation of total system mass $`M_{\mathrm{tot}}`$ and orbital angular momentum $`J`$) and initially circular orbits is (Paczynski 1971; van den Heuvel 1994):
$$\frac{P_f}{P_{}}=\left(\frac{M_{}m_{}}{M_fm_f}\right)^3,$$
(3)
where $`P_{}`$, $`M_{}`$ and $`m_{}=M_{\mathrm{tot}}M_{}`$ denote the orbital period and component masses before the mass transfer, and $`P_f`$, $`M_f`$ and $`m_f=M_{\mathrm{tot}}M_f`$ are the orbital period and component masses after the transfer. The transformation between orbital separation and the orbital period is given by Kepler’s third law.
Introducing the initial mass ratio $`q_{}=m_{}/M_{}`$ and using equation2, equation3 can be written as
$$\frac{P_f}{P_{}}=\left(\frac{q_{}}{p(q_o+1p)}\right)^3.$$
(4)
Since according to equation2, $`p`$ increases for increasing stellar mass, one observes that, due to the third power in equation4, for the same $`q_{}`$ the orbital period of a very massive system increases much less as a result of the mass transfer, than for systems of lower mass. (The term $`(q_o+1p)`$ changes much less than $`p`$ itself for increasing stellar mass, so for a given $`q_o`$ this term has only a modest effect.) This is the main reason for the systematically longer orbital periods of the Be/X-ray binaries (always $`>16`$ days) relative to those of the OB-supergiant HMXBs (in all but one case: between 1.4 days and 11 days, cf. van den Heuvel, 1983, 1985, 1994). This is illustrated in Table 2 where we list the relative post-mass-transfer periods $`P_f/P_{}`$ for typical Be/X-ray binary progenitor systems, with $`M_{}=10`$$`\mathrm{M}_{}`$ and 12 $`\mathrm{M}_{}`$, respectively, and for two typical OB-supergiant HMXB progenitors with $`M_{}=25`$$`\mathrm{M}_{}`$ and 35 $`\mathrm{M}_{}`$, respectively, for $`q_{}`$ values ranging from 0.4 through 0.8.
### 3.2 Possible effects of further mass transfer and stellar winds on the orbits
#### 3.2.1 Case BB mass transfer
The helium cores left by the 10 $`\mathrm{M}_{}`$ and the 12 $`\mathrm{M}_{}`$ stars have masses of 2.15 $`\mathrm{M}_{}`$ and 2.87 $`\mathrm{M}_{}`$, respectively. During helium-shell burning, when these stars have CO-cores, their outer layers may expand to dimensions of a few to several tens of solar radii, and a second, so-called Case BB, mass transfer may ensue before their cores collapse to neutron stars (Habets 1985,1986ab). However, since the radius of the 2.87 $`\mathrm{M}_{}`$ helium star will not exceed 5$`R_{}`$, a second mass transfer phase is unlikely to occur here. In the case of the 2.15 $`\mathrm{M}_{}`$ helium star, which does attain a large radius, the amount of mass that is in the extended envelope is not more than 0.1 $`\mathrm{M}_{}`$. For these reasons, we will neglect here the effects of Case BB mass transfer, and will assume that these helium stars do not lose any mass before their final supernova explosion. This means, that we will slightly overestimate the imparted runaway velocities (as the orbits at the time of the explosion will be slightly wider than we assume, and the ejected amounts of mass will be somewhat smaller than we assume).
#### 3.2.2 Stellar-wind mass loss in massive stars
Since wind mass-loss rates from Wolf-Rayet (WR) stars –massive helium stars– are much larger (viz.: $`\stackrel{>}{}10^5`$$`\mathrm{M}_{}\mathrm{yr}^1`$) than the mass-loss rates of lower-mass main-sequence stars, for the sake of argument (in order to include only the largest effects) we take into account the effects of the wind mass loss during the WR phase. The effects of these winds are: (1) to widen the orbits, and (2) to considerably decrease the mass of the helium ($``$ WR) star before its core collapses. In the cases of the 25 $`\mathrm{M}_{}`$ and 35 $`\mathrm{M}_{}`$ primary stars, the masses of the helium cores are 9.1 $`\mathrm{M}_{}`$ and 15.4 $`\mathrm{M}_{}`$, respectively. Such stars live $`9\times 10^5`$ years and $`7.5\times 10^5`$ years, respectively, and are expected to lose about 4.0 $`\mathrm{M}_{}`$ and 7.4 $`\mathrm{M}_{}`$ through their wind during this phase of their evolution, respectively<sup>2</sup><sup>2</sup>2We assumed here wind mass-loss rates of $`0.5\times 10^5`$$`\mathrm{M}_{}\mathrm{yr}^1`$ for the 9.1 $`\mathrm{M}_{}`$ star and $`10^5`$$`\mathrm{M}_{}\mathrm{yr}^1`$ for the 15.4 $`\mathrm{M}_{}`$ star, respectively. These rates are in good agreement with observed WR-wind mass-loss rates (cf. Leitherer et al. 1995), but are lower than the rates adopted by Woosley et al. (1995), which may overestimate the real mass-loss rates, since they give for all initial helium star masses, final masses before core collapse of only about 4 $`\mathrm{M}_{}`$.
Thus, at the moment of the supernova explosion the collapsing cores of these stars will have masses of 5 $`\mathrm{M}_{}`$ and 8.0 $`\mathrm{M}_{}`$, respectively. To keep the same notation we will express the relative mass loss in the stellar wind with $`\delta =\mathrm{\Delta }M_{\mathrm{wind}}/M_{}`$. The value for $`\delta `$ is 0.16 for a primary with a mass of 25 $`\mathrm{M}_{}`$ and 0.21 for a 35 $`\mathrm{M}_{}`$ primary star. In the cases of no wind mass loss (in lower mass primaries): $`\delta =0`$.
The wind mass loss will change the post-mass-transfer orbits as follows (van den Heuvel 1994):
$$d\mathrm{log}a=d\mathrm{log}M_{\mathrm{tot}},$$
(5)
and
$$d\mathrm{log}P=2d\mathrm{log}M_{\mathrm{tot}},$$
(6)
where $`a`$ is the orbital separation and $`M_{\mathrm{tot}}`$ the total system mass.
Eq.(6) results in:
$$P/P^{}=\left(\frac{M_{\mathrm{tot}}^{}}{M_{\mathrm{tot}}}\right)^2,$$
(7)
where $`P`$ and $`P^{}`$ correspond to the system masses $`M_{\mathrm{tot}}`$ and $`M_{\mathrm{tot}}^{}`$, respectively. $`M_{\mathrm{tot}}`$ is the total mass at the beginning of the WR phase and $`M_{\mathrm{tot}}^{}`$ the total system mass at the end of this phase, just prior to the supernova explosion of the WR Star. The orbital separation after mass transfer and additional WR mass loss phase is expressed as:
$$\frac{a^{}}{a_{}}=\frac{a_f}{a_{}}\frac{a^{}}{a_f}\left(\frac{q_{}}{p(1+q_{}p)}\right)^2\left(1\frac{\delta }{1+q_{}}\right)^1$$
(8)
### 3.3 Runaway velocities induced by symmetric supernova mass ejection
The runaway velocity imparted to the system by the supernova mass loss is calculated from the loss of momentum of the system during the explosion: $`V_{\mathrm{orb},1}\mathrm{\Delta }M_{\mathrm{sn}}`$, where $`V_{\mathrm{orb},1}`$ is the orbital velocity of the helium star prior to the explosion and $`\mathrm{\Delta }M_{\mathrm{sn}}`$ is the amount of mass ejected in the supernova.
We assume all compact remnants to be a $`1.4\mathrm{M}_{}`$ neutron star. Then $`\mathrm{\Delta }M_{\mathrm{sn}}`$ is given by $`\mathrm{\Delta }M_{\mathrm{sn}}=(p\delta )M_o1.4M_{}`$. The remaining mass of the system is:
$$M_{\mathrm{tot}}^{\prime \prime }=m_f+1.4\mathrm{M}_{}=(q_o+1p)M_o+1.4M_{}.$$
(9)
This yields a recoil velocity (or runaway velocity) of the system of:
$`V_{\mathrm{rec}}`$ $`=`$ $`V_{\mathrm{orb},1}{\displaystyle \frac{\mathrm{\Delta }M_{\mathrm{sn}}}{(q_o+1p)M_o+1.4M_{}}}`$ (10)
Here the second term in the right argument is simply the post supernova eccentricity and we may write simply
$$V_{\mathrm{rec}}=eV_{\mathrm{orb},1}$$
(11)
The relative orbital velocity before the explosion is $`\sqrt{GM_{\mathrm{tot}}^{}/a^{}}`$. One therefore has:
$$V_{\mathrm{orb},1}=\left(\frac{GM_o}{a_o}\right)^{1/2}\frac{p(1+q_op)^2}{q_o(1+q_o)^{1/2}}$$
(12)
Substitution of Eq.(11) into Eq.(10) results now in
$`\text{V}_{\mathrm{rec}}`$ $`=`$ $`\left({\displaystyle \frac{GM_o}{a_o}}\right)^{1/2}`$ (13)
$`\times `$ $`{\displaystyle \frac{p(1+q_op)^2}{q_o(1+q_o)^{1/2}}}{\displaystyle \frac{(p\delta 1.4M_{}/M_o)}{(q_o+1p+1.4M_{}/M_o)}}`$
which in numerical form becomes:
$`V_{\mathrm{rec}}`$ $`=`$ $`212.9[\mathrm{kms}^1]\left({\displaystyle \frac{M_o}{[M_{}]}}{\displaystyle \frac{[\mathrm{days}]}{P_o}}\right)^{1/3}`$ (14)
$`\times `$ $`{\displaystyle \frac{p(q_o+1p)^2(p\delta 1.4M_{}/M_o)(1+q_o\delta )}{q_o(1+q_o)^{2/3}(q_o+1p+1.4M_{}/M_o)}}`$
Fig.1 shows for $`P_{}=5`$ days the values of $`V_{\mathrm{rec}}`$ as a function of $`q_{}`$ for the four primary masses of Table2, using the $`\mathrm{\Delta }M_{\mathrm{sn}}`$ and $`\mathrm{\Delta }M_{\mathrm{wind}}`$ as given above. The figure shows that for the “Be-systems” (initial primary masses 10 $`\mathrm{M}_{}`$ and 12 $`\mathrm{M}_{}`$ yielding Be-star masses ranging from 11.85 $`\mathrm{M}_{}`$ to 18.7 $`\mathrm{M}_{}`$) the expected recoil velocities range from 5 to 21 $`\mathrm{kms}^1`$, whereas for the “OB-supergiant systems” (with OB-companions between 25 $`\mathrm{M}_{}`$ and 40 $`\mathrm{M}_{}`$) they range between 21 and $`>80`$$`\mathrm{kms}^1`$, respectively. These velocities correspond to transverse velocities that are $`\pi /4`$ times these values, i.e.: 3.9 to 17 $`\mathrm{kms}^1`$ for the Be-systems, and 16.5 to $`>71`$ $`\mathrm{kms}^1`$ for the OB-supergiant systems with neutron stars. Thus one expects average transverse velocities of order 10.5$`\mathrm{kms}^1`$ and 45$`\mathrm{kms}^1`$ for the Be/X-ray binaries and OB-supergiant systems, respectively. For both the Be/X-ray binaries and the OB-supergiant systems Be/X-ray the predicted and observed mean transverse runaway velocities agree well: $`15\pm 6`$$`\mathrm{kms}^1`$and $`42\pm 14`$, respectively.
As Eq.(14) shows, the dependence of the recoil velocity on $`P_{}`$ is rather weak, so for initial orbital periods between a few days and 10 days these results don’t change by more than a factor 1.5. Therefore, certainly qualitatively, Fig.1 is representative for the two types of systems. Eq.(14) further shows that the large difference in runaway velocity between the two types of systems is due to a combination of two factors, as follows: (1) the larger fractional helium core masses ($`p`$) in the more massive systems, which cause their pre-supernova orbital periods to be shorter and thus their pre-supernova orbital velocities to be larger than those of the lower-mass systems; and (2) the much lower amounts of mass ejected ($`\mathrm{\Delta }M_{\mathrm{sn}}`$) in the lower mass systems compared to the systems of higher mass, which leads to a lower recoil effect.
Relaxing the assumption that mass is conserved during the phase of mass transfer changes little, which we will discuss now.
## 4 The effects of non-conservative mass transfer
In the above it was assumed that the case B mass transfer was conservative in all systems. For the Be-systems this “conservative” assumption seems confirmed quite straightforwardly as the Be nature is interpreted by the accretion of angular momentum and thus of mass. On the other hand, for the OB-supergiant X-ray binaries several authors (starting with Flannery and Ulrich 1977 for the Cen X-3 system) have pointed out that certainly in part of the systems the mass transfer has been non-conservative, and there is a considerable evolution for massive close binaries altogether (De Loore & De Greeve 1992). Indeed, close Wolf-Rayet binaries with high mass ratios $`q=M_{\mathrm{WR}}/M_{OB}`$ such as CQ Cep ($`P=1.64`$ days, $`q=1.19`$) and CX Cep ($`P=2.22`$ days, $`q=.44`$) cannot have been produced by conservative evolution, and just these systems are the progenitors of the OB-supergiant X-ray binaries (cf. van den Heuvel 1994).
The amount of mass lost from the system during the transfer will depend on the initial mass ratio $`q_{}`$ of the system. For small $`q_{}`$ the companion will accrete little and most of the envelope mass of the primary will be lost from the system. On the other hand, for large $`q_{}`$ little mass will be lost from the system. Therefore, in order to study the effect of mass and angular momentum loss on the runaway velocity, we assume as a first approximation that the fraction $`f`$ of the primary’s envelope which is accreted by the companion star is proportional to the initial mass ratio $`q_{}`$ (Portegies Zwart 1995)
$$f=q_{}.$$
(15)
After mass transfer the secondary mass then becomes
$$m_f=M_{}+fM_{}(1p)M_{}q_{}(2p).$$
(16)
The gas lost by the donor leaves with low velocity but gains angular momentum via the interaction with the companion star. It finally leaves the binary system via the second Lagrangian point $`L_2`$, carrying specific angular momentum with it (de Loore & De Greve 1992). The specific angular momentum of this lost matter is considerably larger than what is lost in the specific amount of angular momentum in the stellar wind (given by Eq. 5), see for example Soberman et al. (1997).
We assume that the mass that leaves the system carries a fraction $`\beta `$ of the specific angular momentum of the binary. We can then write the change in orbital separation due to mass transfer as
$$\frac{a^{}}{a_{}}=\left(\frac{M_fm_f}{M_{}m_{}}\right)^2\left(\frac{M_f+m_f}{M_{}+m_{}}\right)^{2\beta +1}.$$
(17)
and use it as an alternative for Eq.(7). Following Portegies Zwart (1995) we use $`\beta =3`$.
Eq. (8) then becomes:
$$\frac{a^{}}{a_{}}=\left(\frac{1}{q_{}p(2p)}\right)^2\left(\frac{1+q_{}}{p+q_{}(2p)}\right)^{2\beta 1}\left(1\frac{\delta }{1+q_{}}\right)^1$$
(18)
The result of this calculation is presented as the dotted lines in Fig.1. The small number near each $``$ indicates the mass of the visible component, which is smaller than if mass transfer would proceed conservatively. One observes that for the same mass of the visible component of the binary, the runaway velocity of the OB-system is between 50 and 100 per-cent larger than in the conservative case. The higher velocity of the binary is mainly caused by the smaller orbital separation at the moment of the supernova. We thus see, from this simple numerical experiment, that non-conservative mass transfer makes the difference in runaway velocity between the two types of high mass X-ray binaries considerably larger.
## 5 Predicted and observed orbital eccentricities of the Be/X-ray binaries: evidence for kicks
### 5.1 Orbital eccentricities of Be/X-ray binaries in case of symmetric ejection
In the case of spherically symmetric mass ejection the orbital eccentricity induced by the mass loss is (cf. Hills 1983):
$$e=\frac{\mathrm{\Delta }M_{\mathrm{sn}}}{M_{\mathrm{tot}}^{}\mathrm{\Delta }M_{\mathrm{sn}}}\frac{\mathrm{\Delta }M_{\mathrm{sn}}}{M_{\mathrm{tot}}^{\prime \prime }}.$$
(19)
One expects that because of the extensive mass transfer and the fact that before the mass transfer the primary was a (sub)giant, the orbits just prior to the explosions are circular. Hence, in case of spherically symmetric mass ejection, one expects the eccentricities of the Be/X-ray binaries simply to be given by Eq.(19).
Table 3 shows that for the Be/X-ray binaries resulting from systems with an initial primary mass of 10 $`\mathrm{M}_{}`$, the orbital eccentricities expected on the basis of Eq.(19) range from 0.045 to 0.060. For systems resulting from binaries with primaries of 12 $`\mathrm{M}_{}`$, the eccentricities range from 0.073 to 0.096. It should be noted that these are, in fact, overestimates, since we ignored case BB mass transfer, which would still have somewhat reduced these values.
Orbital eccentricities are known for only five Be/X-ray binaries, as is listed in Table3. They range from 0.3 to $`>0.7`$, with an average of about 0.5. For the two long-period systems, 1145-619 and 1258-613, the orbital eccentricities have not yet been measured, but a rough estimate of their values can be made as follows. Both systems are recurrent transients, with outbursts occurring once per orbit, when the Be star is active, presumably when the stars are near periastron. The same is true for the systems 0115+63, V0331+53 and EX02030+375 when their Be-stars are in an active phase. For the latter systems one calculates from their orbital periods and eccentricities that within 20 percent their periastron distances are the same. Apparently, this is the periastron distance required for triggering an outburst when the Be-star is in an active phase. It thus seems reasonable to assume that the same is true for 1145-619 and 1258-613. Using this, one finds the latter systems to have eccentricities of between 0.75 and 0.83, and between 0.70 and 0.80, respectively. To be conservative, we have indicated this in Table 3 as: $`e0.70`$.
Two more systems consisting of a B-star and a neutron star are known: the binary radio pulsars PSRJ 1259-63 and PSRJ 0045-7319. These have very eccentric orbits as indicated in Table3. So, in total we have nine B-star plus neutron star systems with measured or estimated orbital eccentricities. References to the orbital parameters of these systems are indicated in the table.
Observations show that all binaries –including detached ones– with orbital periods shorter than 10 days have circular orbits, whereas detached systems with longer orbital periods do not. This suggests that in systems with orbital periods shorter than 10 days tidal forces are effective in circularizing the orbits on a timescale considerably shorter than the lifetimes of the components of the binary, whereas in wider systems they apparently are not. Since the Be/X-ray binaries are detached systems (cf. van den Heuvel & Rappaport 1987; van den Heuvel 1994) and have orbital periods longer than 16 days, it is not surprising that their orbits have not yet been circularized.
The lifetime of a Be/X-ray binary is expected to be of the order of a few million years up to about 10 Myr, the lifetime of the Be companion of the neutron star. The timescale for tidal circularization for main-sequence binaries with orbital periods $`>16`$ days is at least a few tens of Myr (see Zahn 1977, Kochanek 1992). Therefore it is unlikely to catch the binary in the circularization process. Therefore, we expect that the eccentricities for the Be/X-ray binaries in Table3 are still close to those just after the supernova explosion. The orbits of the high-mass X-ray binaries with orbital periods $`<10`$ days are all practically circularized by tidal effects.
It should be noted that if the eccentricities of the Be-systems had resulted from spherically symmetric supernova mass ejection, the amounts of mass ejected in their supernovae should have been very large, of order 4 to over 7 solar masses (see for example Iben & Tutukov 1998). Since in the case of symmetric mass ejection the orbital eccentricity and runaway velocity are directly proportional to each other (see Eqs. and ), also the induced runaway velocities should have been much larger than observed. For example, induction of an eccentricity 0.5 with a symmetric explosion requires 1/3 of the system mass to be ejected in the explosion. With a Be star of 12 $`\mathrm{M}_{}`$, as is representative for a typical B0.5Ve star, and a neutron star mass of 1.4 $`\mathrm{M}_{}`$, the initial system mass in this case must have been 20.1 $`\mathrm{M}_{}`$, implying an ejected amount of mass of 6.7$`\mathrm{M}_{}`$. In order to obtain a post-supernova orbital period of about 30 days, as is typical for many Be/X-ray binaries, the initial orbital period in this example must have been around 11 days. With this initial period, and 6.7 $`\mathrm{M}_{}`$ explosively ejected, the induced runaway velocity $`V_{\mathrm{rec}}`$ would be $`87`$$`\mathrm{kms}^1`$(see the equations in Sect. 3.3), which is advariant with the observed velocities.
Similarly, if the induced eccentricity would be 0.3, one finds that for the same final system mass and orbital period, the runaway velocity induced by the explosion would have been about 45 $`\mathrm{kms}^1`$.
As these velocities are some 5, respectively 2.5 times larger than the mean excess space velocity of 19 $`\mathrm{kms}^1`$ \[$`(4/\pi )\times 15\mathrm{kms}^1`$\] of the Be/X-ray binaries, it is clear that the orbital eccentricities of the Be/X-ray binaries cannot be due purely to symmetric mass ejection in the supernova explosion.
The only way to obtain both a low runaway velocity of the system and the high orbital eccentricities listed in Table3 is by having a small amount of mass ejected in the supernova, in combination with a velocity kick of order 60 to 250 $`\mathrm{kms}^1`$ imparted to the neutron star at birth. We describe below how these required kick velocities were calculated. The randomly directed kick hardly changes the runaway velocity of the system, as the impulse of the kick imparted to the neutron star is shared by the entire system (with a mass of order 15 solar masses in the case of the Be/X-ray binaries), and thus the kick velocity is “diluted” to an extra velocity of the system of only 4 to 16 $`\mathrm{kms}^1`$, in a random direction. Adding this velocity quadratically (because of its random direction) to the velocity of between 5 and 21 $`\mathrm{kms}^1`$ imparted to the systems purely by the mass loss (Fig. 1), one obtains mean runaway velocities of between 6 and 21 $`\mathrm{kms}^1`$ for a 60 $`\mathrm{kms}^1`$ kick and between 17 and 26 $`\mathrm{kms}^1`$ for a 250 $`\mathrm{kms}^1`$ kick.
These values are in good agreement with the observed mean excess space velocities of Be/X-ray binaries of $`19\pm 8`$$`\mathrm{kms}^1`$($`\pi /4`$ times their average peculiar tangential velocities).
We calculated the minimum kick velocities that have to be imparted to the neutron star during the supernova explosion in order to obtain the presently observed orbital eccentricities of the Be-systems in Table3. We used the equations derived by Wijers et al. (1992). The minimum required kick velocitiy is the one that is imparted in the orbital plane in the direction of motion of the pre-supernova star (assuming the initial orbit was circular). We assumed in these calculations that the Bstars have a mass of $`15\mathrm{M}_{}`$, as corresponds to a B0-1 main-sequence star, and that the neutron star has a mass of $`1.4\mathrm{M}_{}`$. (For B-star masses in the range $`1020\mathrm{M}_{}`$ the required minimum runaway velocities do not differ by more than $`\pm `$ 10 per cent from the values for $`15\mathrm{M}_{}`$). The table shows that the required minimum kick velocities range from about 50 $`\mathrm{kms}^1`$ to about 200 $`\mathrm{kms}^1`$. Assuming the real kick velocities to be randomly distributed, the required kicks become $`\sqrt{3/2}`$ times larger, and range from about 60 to about 250 $`\mathrm{kms}^1`$.
We conclude from the above that the combination of low mean space velocity of the Be/X-ray binaries and large mean orbital eccentricity provides unequivocal evidence for the existence of velocity kicks imparted to neutron stars at their birth.
An alternative way to approach the problem of the orbital eccentricities is to calculate, from the measured mean runaway velocities of Be/X-ray systems, what orbital eccentricity these systems should have had, were this runaway velocity imparted by purely symmetric mass ejection. This is the topic of the next section.
### 5.2 Predicted relation between orbital eccentricity and runaway velocity expected in case of symmetric explosions - comparison with observations
Eq.(14) yields:
$$\frac{e}{1+e}=\frac{\mathrm{\Delta }M_{\mathrm{sn}}}{M_{\mathrm{tot}}^{}},$$
(20)
Combination of Eqs(10) and (20) yields:
$$V_{\mathrm{rec}}=\sqrt{\frac{GM_{\mathrm{tot}}^{\prime \prime }}{a^{}}}\frac{m_f}{M_{\mathrm{tot}}^{\prime \prime }}\frac{e}{1+e},$$
(21)
where $`a^{}`$ is the pre-supernova orbital radius. The semi-major axis after the supernova $`a^{\prime \prime }`$ follows from:
$$\frac{a^{}}{a^{\prime \prime }}=1\frac{\mathrm{\Delta }M_{\mathrm{sn}}}{M_{\mathrm{tot}}^{\prime \prime }},$$
(22)
and by writing
$$M_{\mathrm{tot}}^{}=M_{\mathrm{tot}}^{\prime \prime }+\mathrm{\Delta }M_{\mathrm{sn}}=M_{\mathrm{tot}}^{\prime \prime }(1+\frac{\mathrm{\Delta }M_{\mathrm{sn}}}{M_{\mathrm{tot}}^{\prime \prime }}).$$
(23)
one obtains after insertion of Eq.(18) in Eq.(16):
$$V_{\mathrm{rec}}^2=\frac{GM_{\mathrm{tot}}^{\prime \prime }}{a^{\prime \prime }}\frac{1+\mathrm{\Delta }M_{\mathrm{sn}}/M_{\mathrm{tot}}^{\prime \prime }}{1\mathrm{\Delta }M_{\mathrm{sn}}/M_{\mathrm{tot}}^{\prime \prime }}\left(\frac{m_f}{M_{\mathrm{tot}}^{\prime \prime }}\right)^2\left(\frac{e}{1+e}\right)^2.$$
(24)
Defining now the presently observed mean orbital velocity by
$$V_{\mathrm{orb}}^2=\frac{GM_{\mathrm{tot}}^{\prime \prime }}{a^{\prime \prime }},$$
(25)
and substituting $`\mathrm{\Delta }M_{\mathrm{sn}}/M_{\mathrm{tot}}^{\prime \prime }`$ from Eq.(19) one obtains:
$$\frac{V_{\mathrm{rec}}}{V_{\mathrm{orb}}}\frac{M_{\mathrm{tot}}^{\prime \prime }}{m_f}=\frac{e}{(1e^2)^{1/2}}.$$
(26)
This defines, in the case of symmetric supernova-mass ejection the relation that is expected to be found between the observed system runaway velocity $`V_{\mathrm{rec}}`$ and the observed orbital eccentricity $`e`$, for a system with a Be/X-ray star of mass $`m_f`$, and observed mean orbital velocity $`V_{\mathrm{orb}}`$. Since $`M_{\mathrm{tot}}^{\prime \prime }=m_f+1.4`$$`\mathrm{M}_{}`$, and since in general $`m_f>10`$$`\mathrm{M}_{}`$, the quantity $`M_{\mathrm{tot}}^{\prime \prime }/m_f`$ is close to unity. Defining:
$$f_v\frac{V_{\mathrm{rec}}}{V_{\mathrm{orb}}}\frac{M_{\mathrm{tot}}^{\prime \prime }}{m_f}=\frac{e}{(1e^2)^{1/2}},$$
(27)
one obtains a simple relation between $`f_v`$ and $`e`$, the plotted curve in Fig.2. In the case of symmetric supernova mass ejection, the observed value of $`f_v`$ of a Be/X-ray binary should be related to the observed orbital eccentricity according to this curve, which shows that large eccentricities correspond to large runaway velocities.
In Fig.2 we also plotted the values of $`f_v`$ and $`e`$ for the nine systems with observed orbital periods and eccentricities (see Table3), taking runaway velocities $`V_{\mathrm{rec}}`$ in the observed range $`19\pm 8`$ $`\mathrm{kms}^1`$ for the Be X-ray binaries. We assumed a Be-star mass of $`15\mathrm{M}_{}`$. The figure shows that all systems fall far below the curve expected for symmetric supernova mass ejection. This again shows that the combination of low runaway velocities and large orbital eccentricities as observed in the Be/X-ray binaries cannot be obtained by symmetric mass ejection in the supernovae, and that a velocity kick imparted to the neutron stars at birth is absolutely required.
## 6 Conclusions
The measured tangential velocities of the Be/X-ray binaries and OB-supergiant X-ray binaries by the Hipparcos satellite confirm the expectations from the evolution of massive close binaries in which little mass is lost from the binary systems during the first mass transfer phase. The much higher tangential velocities of supergiant X-ray binaries than those of the Be-systems follow from a combination of (1) the much larger fractional helium core masses in the progenitors of the OB-supergiant systems which cause their pre-supernova orbital periods to be shorter, and thus their pre-supernova orbital velocities to be much larger than those of the less massive Be-systems, and (2) the much lower amounts of mass ejected during the supernova explosion in the lower-mass Be-systems compared to the OB-supergiant systems.
The combination of a high orbital eccentricity with a low space velocity observed for the Be type X-ray binaries can only be understood if a kick with appreciable velocity –in the range 60 to 250$`\mathrm{kms}^1`$– is imparted to the newly born neutron star. Such a kick tends to only slightly affect the space velocity of the binary system since the neutron star has to drag along its massive companion. The orbital eccentricity, however, is strongly affected by such a asymmetric velocity kick. If the supernova explosions in these systems had been symmetric, the high orbital eccentricities observed in the class of Be X-ray binaries are impossible to reconcile with their on average low runaway velocities.
* This work was supported by Spinoza Grant 08-0 to E.P.J. van den Heuvel and by NASA through Hubble Fellowship grant HF-01112.01-98A awarded (to SPZ) by the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc., for NASA under contract NAS 5-26555. LK is supported by a fellowship of the Royal Netherlands Academy of Arts and Sciences. The first two authors thank the Institute for Theoretical Physics UCSB, where most of this research was carried out in the fall of 1997.
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# Optical conductivity in the normal state fullerene superconductors
## I Introduction
The optical spectra of fullerene superconductors in the normal state were found to exhibit some unusual features . The optical conductivity, $`\sigma (\omega )`$, deviates considerably from the simple Drude behavior expected for conventional metals: the spectral weight of the Drude peak is reduced by about an order of magnitude and transfered to a mid-infrared (MIR) region around $`0.06`$ eV. This suggests that the strong interaction effects due to the Coulomb and electron-phonon interactions should be important in the optical spectra of the fullerenes. Understanding this unusual behavior in the optical conductivty, therefore, could reveal important information about the fullerenes and contribute to understanding other physical properties of the material.
The optical conductivity $`\sigma (\omega )`$ represents the rate at which electrons absorb the incident photons at energy $`\omega `$, and is a useful probe in determining electronic characteristics of the material under study. For an ideal free electron gas, where the interactions between the electrons, and between the electrons and phonons are neglected, and the impurity scattering rate $`1/\tau 0`$, $`\sigma (\omega )`$ collapses to a delta-function, $`\sigma (\omega )=D_{tot}\delta (\omega )`$, where the coefficient $`D_{tot}`$ represents the total spectral weight. In this case, the optical conductivity sum rule, $`_0^{\mathrm{}}𝑑\omega \sigma (\omega )=D_{tot}=\pi e^2n/m`$, where $`n`$ is the density and $`m`$ is the mass of the electrons, is exhausted entirely by the delta function Drude contribution alone. When the material becomes dirtier, the Drude peak of the optical conductivity acquires the Lorentzian shape with the width of $`1/\tau `$. The conductivity sum rule is still exhausted by the Drude part alone when only the impurity scatterings are present in the system. The total weight $`D_{tot}`$, however, can change as the inpurity scatterings or other interactions are introduced when we consider a finite bandwidth, because the projection to a restricted basis set disregards all excitations to higher energy than the bandwidth. When other interactions are present, the free-carrier Drude weight is reduced by the quasiparticle renormalization factor $`Z`$ such that $`D=Z^1D_{tot}`$, and the missing spectral weight from the Drude part is transfered to a higher energy region of $`\sigma (\omega )`$ reflecting the excitation of incoherent scatterings.
The experimentally measured $`\sigma (\omega )`$ in the normal state $`\mathrm{A}_3\mathrm{C}_{60}`$ shows a remarkable reduction of Drude weight and, concomitantly, a pronounced MIR absorption below the inter-band absorption peak: DeGiorgi $`etal.`$ found a pronounced MIR peak around 0.06 eV and analysized that the Drude weight is reduced to about $`0.10.2`$ of the total intra-band spectral weight , while Iwasa $`etal.`$ observed the MIR absorption peak around 0.4 eV and determined that the Drude weight is reduced to about 0.6 of the total intra-band spectral weight . Although their results show somewhat different Drude weight and MIR absorption energy each other, the pronounced suppression of the Drude weight and the accompanying MIR absorption imply strong electron-phonon and/or electron-electron interactions in this material.
In order to understand this unusual feature in $`\sigma (\omega )`$ of doped fullerenes, Gunnarsson $`etal.`$ studied the effects of the electron-phonon interaction on $`\sigma (\omega )`$ assuming that the Migdal theorem is valid . They showed that the electron-phonon interaction leads to a narrowing of the Drude peak by the factor $`Z=1+\lambda `$, where $`\lambda `$ is the dimensionless electron-phonon coupling constant, and a transfer of the depleted Drude weight to a MIR region at somewhat larger energies than the phonon energy. Their results, however, are far from sufficient to describe experimental observations. Therefore, they hinted that the Coulomb interaction between conduction electrons, which is neglected in their study, could lead to futher reduction of the Drude weight and more pronounced MIR absorption. On the other hand, one of the present authors recently found, by studying the NMR coherence peak supression in the fullerene superconductors, that the Coulomb interaction between conduction electrons, characterized by $`UN_F0.30.4`$, where $`U`$ is the effective Coulomb interaction and $`N_F`$ is the density of states (DOS) at the Fermi level, should be included in addition to the electron-phonon interaction to understand the various experimental observations in fullerenes in a coherent way . We, therefore, included the electron-electron as well as electron-phonon interactions at the presence of the impurity scatterings in the present paper, to better understand the experimentally obsered unusual features in the optical spectra of the fullerene superconductors in the normal state.
For fullerene superconductors, the Fermi energy $`\epsilon _F=B/20.20.3`$ eV and the average phonon frequency $`\omega _{ph}0.050.15`$ eV, where $`B`$ is the bandwidth. Therefore, $`\omega _{ph}/\epsilon _F1`$ for fullerenes unlike conventional metals, where $`\omega _{ph}/\epsilon _F1`$. When $`\omega _{ph}/\epsilon _F1`$, the phonon vertex correction becomes important because the Migdal theorem does not hold , and the frequency dependence of the effective Coulomb interaction, $`V_{eff}(\omega )`$, should be considered because the frequency scale at which $`V_{eff}(\omega )`$ varies is comparable with that of electron-phonon interaction . In this present work, concerned with the effects of the Coulomb and electron-phonon interactions on the optical spectra in the narrow band fullerene superconductors, the vertex correction is incorporated in calculating the electron self-energy . The Coulomb interaction, modelled in terms of the onsite Hubbard repulsion, is included on an equal footing with the electron-phonon interaction, and considered fully self-consistently in calculating the effective electron-electron interaction . The effective electron-electron interaction becomes frequency dependent through the screening. The impurity effects are included with the $`t`$-matrix approximation.
Through the relation
$`\mathrm{\Sigma }(ip)=G^1(ip)G_0^1(ip),`$ (1)
one obtain the electron self-energy $`\mathrm{\Sigma }(ip)`$ in the Mastubara frequency, which gives $`\mathrm{\Sigma }(\omega )`$ in the real frequency after the analytic continuation. $`G_0`$ and $`G`$ are, respectively, the bare and renormalized electron Green’s functions. $`\mathrm{\Sigma }(\omega )`$ or $`Z(\omega )`$, where the renormalization function $`Z(\omega )`$ is given by $`\mathrm{\Sigma }(\omega )=\omega \omega Z(\omega )`$, defines the single-particle Green’s function of an interacting system as
$`G^1=\omega \xi _k\mathrm{\Sigma }(\omega )=\omega Z(\omega )\xi _k,`$ (2)
where $`\xi _k`$ is the electron energy measured from the chemical potential, $`\xi _k=\epsilon _k\mu `$. Then, the optical conductivity can be obtained by calculating the current-current correlation function, $`\mathrm{\Pi }(i\omega )`$, using the renormalized Green’s function obtained from solving Eq. (1) self-consistently. The calculated optical conductivity shows a strong reduction of Drude weight and a broad MIR absorption, although the MIR feature around 0.06 eV is less pronounced and broader compared with experimental observations.
This paper is organized as follows: In the following section, we present the Eliashberg-type formalism in the Matsubara frequency to calculate the renormalized Green’s function with the impurity, electron-phonon, and Coulomb interactions included self-consistently. We then describe the analytic continuation procedure to obtain the renormalization function $`Z(\omega )`$ in the real frequency. The optical conductivity calculated with the renormalized Green’s function is presented in Sec. III. We will discuss how the Drude part and the MIR absorption of $`\sigma (\omega )`$ are affected as the impurity scattering rate, the electron-phonon and electron-electron interactions are varied. These result will then be compare with the experimental observations. Finally, Sec. IV is for the summary and some concluding remarks.
## II Formalism
The optical conductivity is calculated from the current-current correlation function, $`\mathrm{\Pi }(\omega )`$, as $`\sigma (\omega )=\frac{i}{\omega }lim_{q0}\mathrm{\Pi }(q,\omega )`$ . We use the approximation where the electron self energy is momentum independent. In this case, it can be shown that the vertex correction in the current-current correlation function vanishes for $`q0`$ . This leads to
$$\mathrm{\Pi }(i\omega _m)=\frac{2e^2}{3m^2V}\underset{\stackrel{}{p}}{}\stackrel{}{p}^2\frac{1}{\beta }\underset{ip_n}{}G(\stackrel{}{p},ip_n+i\omega _m)G(\stackrel{}{p},ip_n),$$
(3)
where $`ip_n=\pi T(2n+1)`$ and $`i\omega _m=2\pi Tm`$ are, respectively, fermion and boson Mastubara frequencies, where $`T`$ is the temperature, $`m`$ and $`n`$ are the integers. $`\beta =1/k_BT`$, and $`V`$ is the volume. The evaluation of Eq. (3) using Eq. (2) produces
$`\mathrm{\Pi }(i\omega )={\displaystyle \frac{2\pi e^2n}{m}}{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{ip}{}}{\displaystyle \frac{\theta (ip+i\omega )\theta (ip)}{(p+\omega )Z(ip+i\omega )pZ(ip)}}`$ (4)
in the Mastubara frequency. After performing the analytic continuation of $`i\omega \omega +i\delta `$, to the real frequency, the optical conductivity is given by
$`\sigma (\omega )`$ $`=`$ $`{\displaystyle \frac{1}{\omega }}{\displaystyle \frac{e^2n}{m}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\epsilon [f_F(\epsilon )f_F(\epsilon +\omega )]`$ (5)
$`\times `$ $`Re[i{\displaystyle \frac{\theta (\epsilon +i\delta )\theta (\epsilon +\omega +i\delta )}{\epsilon Z(\epsilon +i\delta )(\epsilon +\omega )Z(\epsilon +\omega +i\delta )}}.`$ (7)
$`.i{\displaystyle \frac{\theta (\epsilon i\delta )\theta (\epsilon +\omega +i\delta )}{\epsilon Z(\epsilon i\delta )(\epsilon +\omega )Z(\epsilon +\omega +i\delta )}}],`$
where $`f_F(\epsilon )=1/(1+e^{\beta \epsilon })`$ is the Fermi distribution function, and $`\theta (\omega +i\delta )=\mathrm{tan}^1\left[\frac{i\epsilon _F}{\omega Z(\omega +i\delta )}\right]`$. The finite conduction bandwidth $`B`$ with a constant DOS is explicitly considered through the factor of $`\theta `$, which is $`\pi /2`$ for the usual case of infinite bandwidth metal. In order to calculate the optical conductivity from Eq. (7) we need $`Z(\omega )`$ which defines single-particle interacting Green’s function $`G(\omega )`$. This can be obtained by solving Eq. (1) self-consistently. The electron self-energy is obtained by calculating the exchange diagram of the renormalized electron Green’s function and the effective electron-phonon and Coulomb interactions with the vertex correction included via the method of Nambu. The Coulomb interaction, modelled in terms of the onsite Hubbard repulsion for simplicity, is included on an equal footing with the electron-phonon interaction. The impurity effects are included with the $`t`$-matrix approximation. The Eliashberg-type equation can be written in the Mastubara frequency as
$`Z_np_n`$ $`=`$ $`p_n+{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{m}{}}[\lambda _{ph}(nm)\lambda _{ch}(nm)`$ (8)
$`+`$ $`\lambda _{sp}(nm)]2\theta _m\mathrm{\Gamma }+{\displaystyle \frac{1}{\pi \tau }}\theta _n,`$ (9)
where $`\theta _n=\mathrm{tan}^1(\frac{B}{2p_nZ_n})`$, and $`\lambda _{ph}(nm)=_0^{\mathrm{}}𝑑\mathrm{\Omega }\frac{\alpha ^2F(\mathrm{\Omega })2\mathrm{\Omega }}{[\mathrm{\Omega }^2+(p_np_m)^2]}`$ is the electron-phonon interaction kernel. $`\lambda _{ch}(nm)`$ and $`\lambda _{sp}(nm)`$ are, respectively, the interactions in the charge and spin channels due to the Hubbard repulsion. They are determined self-consistently as
$`\lambda _{ch}(k)`$ $`=`$ $`UN_F\left\{{\displaystyle \frac{1}{2}}\chi _n+\chi _n^2\mathrm{ln}[1+1/\chi _n]\right\},`$ (10)
$`\lambda _{sp}(k)`$ $`=`$ $`UN_F\left\{{\displaystyle \frac{1}{2}}+\chi _n+\chi _n^2\mathrm{ln}[11/\chi _n]\right\},`$ (11)
where $`\chi _n(k)`$ is the dimensionless susceptibility given by
$`\chi _n(k)={\displaystyle \frac{N_FU}{\epsilon _F}}{\displaystyle \frac{1}{\beta }}{\displaystyle \underset{l}{}}\theta _l\theta _{l+k}.`$ (12)
The $`\mathrm{\Gamma }`$ on the right hand side of Eq. (9) represents the vertex correction satisfying the Ward-identity. When we neglect the vertex correction, $`\mathrm{\Gamma }=1`$. If we assume a weak frequency dependence of $`\mathrm{\Gamma }`$, the vertex function $`\mathrm{\Gamma }`$ reduces to $`Z(ip_m)`$. In this work, we treat the vertex correction exactly, and $`\mathrm{\Gamma }`$ is given by
$`\mathrm{\Gamma }=\left[{\displaystyle \frac{ip_nZ(ip_n)ip_mZ(ip_m)}{ip_nip_m}}\right].`$ (13)
Solving Eq. (9) self-consistently yields $`Z(i\omega )`$ in the Mastubara frequency. In order to calculate $`\sigma (\omega )`$, analytic continuation of $`i\omega \omega +i\delta `$ should be performed to get $`Z(\omega )`$ in real frequency. The numerically exact analytic continuation of standard Eliashberg equation is usually performed by the iterative method developed by Marsiglio, Shossmann, and Carbotte (MSC) using a mixed-representation . But when we include the vertex function exactly, the MSC method can not be applied because it needs a specific form of equation. Here, in order to consider vertex correction exactly, we do the alanytic continuation by employing the iterative method extended by Takada . In this case, Eq. (9) is transformed to a mixed representation as fallow:
$`Z(\omega )`$ $`=`$ $`\stackrel{~}{Z}(\omega )+{\displaystyle _0^{\mathrm{}}}d\mathrm{\Omega }P(\mathrm{\Omega })\{[n_B(\mathrm{\Omega })+n_F(\omega +\mathrm{\Omega })]`$ (14)
$`\times `$ $`G(\omega +\mathrm{\Omega })\left[{\displaystyle \frac{(\omega +\mathrm{\Omega })Z(\omega +\mathrm{\Omega })\omega Z(\omega )}{\mathrm{\Omega }}}\right]`$ (15)
$`+`$ $`[n_B(\mathrm{\Omega })+n_F(\mathrm{\Omega }\omega )]G(\omega \mathrm{\Omega })`$ (16)
$`\times `$ $`[{\displaystyle \frac{(\omega \mathrm{\Omega })Z(\omega \mathrm{\Omega })\omega Z(\omega )}{\mathrm{\Omega }}}]\},`$ (17)
where
$`\stackrel{~}{Z}(\omega )`$ $`=`$ $`1+{\displaystyle \frac{1}{\omega \beta }}{\displaystyle \underset{m}{}}{\displaystyle _0^{\mathrm{}}}𝑑\mathrm{\Omega }P(\mathrm{\Omega })\left({\displaystyle \frac{1}{ip_m\omega \mathrm{\Omega }}}{\displaystyle \frac{1}{ip_m\omega +\mathrm{\Omega }}}\right)`$ (18)
$`\times `$ $`G(ip_m)\left[{\displaystyle \frac{ip_mZ(ip_m)\omega Z(\omega )}{ip_m\omega }}\right]+{\displaystyle \frac{i}{\pi \tau }}{\displaystyle \frac{\theta (\omega )}{\omega }},`$ (19)
$`P(\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \frac{1}{\pi }}Im\mathrm{\Lambda }(\mathrm{\Omega })`$ (20)
$`\mathrm{\Lambda }(\mathrm{\Omega })`$ $`=`$ $`\lambda _{ch}(\mathrm{\Omega })\lambda _{ph}(\mathrm{\Omega })\lambda _{sp}(\mathrm{\Omega })`$ (21)
$`G(ip_m)`$ $`=`$ $`2\theta (ip_m),G^R(\omega )=2i\theta (\omega ).`$ (22)
$`\stackrel{~}{Z}(\omega )`$ of Eq. (17) represents the renormalization function obtained by substituting $`i\omega `$ to $`\omega +i\delta `$ $`before`$ the frequency summation. The second term is the correction to $`\stackrel{~}{Z}(\omega )`$ to yield the correct retarded renormalization function $`Z(\omega )`$ one would have obtained if the analytic continuation were performed $`after`$ the frequency summation. Putting the solutions of Eq. (9), $`Z(i\omega )`$, into the $`\stackrel{~}{Z}(i\omega )`$ Eq. (17) yields a self-consistent Eliashberg-type equation in the real frequency. Then, $`Z(\omega )`$ can be obtained by computing iteratively Eq. (17).
In order to model fullerene superconductors, three truncated-Lorentization functions were used to represent $`\alpha ^2F(\mathrm{\Omega })`$ as follow :
$`\alpha ^2F(\mathrm{\Omega })`$ $`=`$ $`{\displaystyle \underset{\nu =1}{\overset{3}{}}}\alpha _\nu ^2F_\nu (\mathrm{\Omega }),`$ (23)
$`F_\nu (\mathrm{\Omega })`$ $`=`$ $`\{\begin{array}{c}\frac{1}{R}\left[\frac{1}{(\mathrm{\Omega }\omega _\nu )^2+\mathrm{\Gamma }^2}\frac{1}{\mathrm{\Gamma }_c^2+\mathrm{\Gamma }^2}\right],for|\mathrm{\Omega }\omega _\nu |\mathrm{\Gamma }_c,\hfill \\ 0,otherwise,\hfill \end{array}`$ (26)
where $`F_\nu (\mathrm{\Omega })`$ is the truncated Lorentizian centered at $`\omega _\nu `$ with the width of $`\mathrm{\Gamma }=\omega _\nu /5`$, $`\mathrm{\Gamma }_c`$ is the cutoff frequency of $`\mathrm{\Gamma }_c=3\mathrm{\Gamma }`$, and $`R`$ is normalization constant such that $`_0^{\mathrm{}}𝑑\mathrm{\Omega }F_\nu (\mathrm{\Omega })=1`$. Various theoretical and experimental estimates do not agree well each other in terms of distribution of coupling strength $`\alpha _\nu ^2`$ among different modes. These estimates show, however, that the phonon frequency derived from intramolecular $`A_g`$ and $`H_g`$ modes are distributed over $`0.030.2`$ eV with the total $`\lambda `$ in the range of $`0.51`$ eV. In view of this, we represent the phonon modes with three groups centered around $`\omega _\nu =0.04,\mathrm{\hspace{0.33em}0.09},\mathrm{\hspace{0.33em}0.19}`$ eV, and $`2N_F\alpha _\nu ^2/\omega _\nu =0.3\lambda _s,\mathrm{\hspace{0.33em}0.2}\lambda _s,\mathrm{\hspace{0.33em}0.5}\lambda _s`$, respectively, for $`\nu =1,2,3`$. Note that $`_{\nu =1}^32N_F\alpha _\nu ^2/\omega _\nu =\lambda _s`$. The $`\lambda _s`$ sets the strength of $`\alpha ^2F(\mathrm{\Omega })`$ and $`N_F\alpha ^2F(\mathrm{\Omega })/\lambda _s`$ is independent of $`\lambda _s`$. For infinite bandwidth superconductors, $`\lambda `$ is equal to $`\lambda _s`$ in the limit $`\mathrm{\Gamma }0`$. For a finite bandwidth system, however, $`\lambda `$ is reduced from $`\lambda _s`$ because the available states to and from which quasiparticles can be sccattered are restricted as the bandwidth is reduced.
## III Results
The self-consistent equation of Eq. (1) is solved numerically as described in the previous section to obtain $`Z(\omega )`$. Then, the optical conductivity is calculated from the Eq. (7). Fig. 1 shows the optical conductivity $`\sigma (\omega )`$ as $`\lambda `$ is varied when the Coulomb interaction $`U`$ is set to 0 for a reference. Here, the Fermi energy $`\epsilon _F`$, temperature $`T`$ and impurity scattering rate $`1/\tau `$ are set to $`0.25,\mathrm{\hspace{0.33em}0.001},\mathrm{\hspace{0.33em}0.01}`$ eV, respectively. This result shows quite a similar behavior to van den Brink $`etal.`$ calculation. As $`\lambda `$ is increased, the width of Drude peak becomes narrower and it’s weight is transferred to a mid-infrared spectrum. However, the reduction of Drude weight is less than the factor of ($`1+\lambda `$), because of the finite bandwidth. The inset shows a MIR absorption spectra obtained by extracting Drude part from the total optical conductivity. In determining the Drude weight, fitting procedure was carefully employed and confirmed by examining zero frequency extrapolation in the Mastubara frequency which is proposed by Scalapino $`etal`$ . The three Lorentizian peaks of $`\alpha ^2F(\omega )`$ in the electron-phonon paring kernel are attributted to the development of these MIR peaks. But, the MIR peaks are broadened and move to slightly high frequencies. Fig. 2 shows the MIR absorption due to the Coulomb interaction. The MIR part is also extracted by fitting as shown in the inset. In order to focus on how $`U`$ affects the total optical conductivity, $`\lambda `$ is set to 0. $`\epsilon _F`$, $`T`$ and $`1/\tau `$ are same as in Fig. 1. The Coulomb interaction induces the strong $`\omega `$ dependence of renormalization function $`Z(\omega )`$, and the low frequency strong $`\omega `$ dependence of $`Z(\omega )`$ distorts the Drude part of optical conductivity and induces the MIR absorption in the fairly low frequency region. As the impurity effect is enhanced, the MIR absorption due to Coulomb interactions tends to shift to higher frequency and finally merge together with the MIR peaks developed by electron-phonon interaction, as shown in Fig. 3 for $`UN_F=0.3`$ and $`\lambda =0.7`$. Note that the position of this merged MIR peak in Fig. 3 is around and above $`0.06`$ eV which is experimentally observed value of DeGiorgi.
Fig. 4 is $`\sigma (\omega )`$ of doped fullerenes with $`T=0.005`$ eV, $`\epsilon _F=0.25`$ eV, $`1/\tau =0.1`$ eV, $`UN_F=0.3`$, and $`\lambda =0.7`$, which is to be compared with the experimental observations. The Drude weight is reduced to 0.467 of the total intra-band optical weight. The reduction factor of the Drude weight by electron phonon interaction is $`1+\lambda `$, and the finite bandwidth futher restricts the reduction factor. It therefore seems unlike that the Drude weight less than about 0.6 of the total intra-band spectral weight can be explained without the Coulomb interactions, when we take $`\lambda 0.70.8`$. The Coulomb interaction suppresses the Drude part substantially by inducing $`\omega `$ dependence of the renormalization function $`Z(\omega )`$ in the low frequency region. We think that the large reduction of Drude weight like the DeGiorgi experiment is a result of the strong Coulomb interaction between conduction electrons in addtion to the electron-phonon interaction. However, our results are still not sufficient to explain experimentally founded results: (a) The Drude weight is about 0.46 of the total intra-band optical weight with a resonable set of parameter values while DeGiorgi found $`0.10.2`$. (b) The MIR absorption is very broad which begins around 0.02 eV, has a peak around 0.07 eV and extends well over the Fermi energy.
## IV Summary and Conclusion
In this paper, we tried to give an explanation for the unusual behavior of optical conductivity in the normal state $`\mathrm{A}_3\mathrm{C}_{60}`$. It is generally accepted that the fullerene superconductor could be characterized by the phonon-mediated $`s`$-wave superconductor . However, a few experiments like optical conductivity still remain not understood by the electron-phonon scattering together with the disorder effects. Our motivation lies in that the fullerene superconductors have such a narrow bandwidth that the phonon frequency, the Coulomb interaction, and the Fermi energy are all comparable, $`\omega _{ph}V\epsilon _F`$. In order to consider properly the Coulomb interaction and the narrow bandwidth of fullerene superconductors, the self-consistent Eliashberg-type coupled equations are solved to obtain the renormalized Green’s function. The theory includes the frequency dependent screened Coulomb interaction together with the electron-phonon interaction and includes the vertex correction via Nambu’s method. In order to treat the vertex function exactly, analytic continuation is performed via the iterative method of mixed repersentation which is developed by Takada. Once we get renormalization function $`Z(\omega )`$ in real frequency, we can calculate optical conductivity in normal states. As we expected, the electron-phonon interaction is not suffficient to resolve the substantial reduction of Drude weight and pronounced MIR peak. The strong Coulomb interaction induces $`\omega `$ dependence in renormalization function $`Z(\omega )`$. As a result, the Drude form in optical spectra is distorted accompanying the reduction of Drude weight. When the impurity effect is enhanced, the MIR absopption induced by strong Coulomb interaction merge together with the MIR peaks due to electron-phonon scattering showing large reduction of Drude weight and MIR peak around $`0.06`$ eV. Although it is not sufficent to explain experimentally founded results, our result is close to DeGiorgi’s experiment. We improve Gunnarsonn’s calculation by considering the electron-electron interaction and finite bandwidth effects explicitly. In conclusion, the unusual behavior of optical conductivity of the normal state $`\mathrm{A}_3\mathrm{C}_{60}`$ reveals the fact that both the Coulomb interaction and electron-phonon interaction are important in examining dynamical properties of fullerene superconductors.
Figure Captions
Figure 1. The optical conductivity as a function of $`\omega `$ for various electron-phonon coupling constants $`\lambda `$ when $`T=0.001`$ eV, $`\epsilon _F=0.25`$ eV and $`1/\tau =0.01`$ eV. $`U`$ is set to 0 for a referrence. As $`\lambda `$ is increased, the width and the weight of the Drude peak are reduced. The inset shows a decomposition of the total conductivity into the Drude and MIR parts for $`\lambda =0.7`$.
Figure 2. The MIR spectra induced by the Coulomb interaction when $`T=0.001`$ eV, $`\epsilon _F=0.25`$ eV, $`1/\tau =0.01`$ eV and $`\lambda =0`$. The inset shows a decomposition of the total conductivity, as in Fig. 1, into the Drude and MIR parts for $`UN_F=0.5`$.
Figure 3. The MIR spectra as the impurity scattering rates $`1/\tau `$ are varied when $`T=0.005`$ eV, $`\epsilon _F=0.25`$ eV, $`\lambda =0.7`$ and $`UN_F=0.3`$. When $`1/\tau =0.01`$ eV the lower peak is mainly from the Coulomb interaction while the other peaks are from the electron-phonon interaction. As $`1/\tau `$ is increased, these peaks are merged altogether and finally evolve into a single broad peak around $`0.060.1`$ eV.
Figure 4. The total optical conductivity with it’s Drude and MIR parts for $`T=0.005`$ eV, $`\epsilon _F=0.25`$ eV, $`1/\tau =0.1`$ eV, $`\lambda =0.7`$ and $`UN_F=0.3`$. The Drude part in the low frequency region is substantially suppressed due to the strong Coulomb interaction. Consequently, the missing spectral weight is transfered to the broad MIR peak, which peaks around 0.07 eV and extends well into the higher energy region. The ratio of the MIR spectral weight to the total intra-band spectral weight is 0.533.
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# A Dehn surgery description of regular finite cyclic covering spaces of rational homology spheres
## 1 Introduction
It is well-known that any closed, oriented, connected 3-manifold may be obtained by Dehn surgery on a link in $`S^3`$. In recent years, 3-manifold theorists have exploited this fact, computing certain 3-manifold invariants by describing how the invariant changes under Dehn surgery. This has been an effective method for computing invariants of individual manifolds. However, until now, there was no known procedure for relating Dehn surgery descriptions of manifolds with those of their covering spaces. Therefore, although invariants for a manifold and a covering space of the manifold could each be computed using Dehn surgery formulas, no general statements could easily be made regarding the relationship between the invariants.
Let $`(X,\stackrel{~}{X})`$ be a 3-manifold pair, where $`X`$ is a rational homology sphere and $`\stackrel{~}{X}`$ is a regular finite cyclic covering space of $`X`$, say with $`\pi _1(X)/\varphi _{}(\pi _1(\stackrel{~}{X}))=/k`$, where $`\varphi :\stackrel{~}{X}X`$ is the projection map. Since this group is abelian, we may factor the quotient map $`\pi _1(X)\pi _1(X)/\varphi _{}(\pi _1(\stackrel{~}{X}))`$ through the first homology group $`\mathrm{H}_1(X;)`$.
###### Definition 1.1
We call the covering $`\stackrel{~}{X}X`$ torsion-split if there exists a homology decomposition $`\mathrm{H}_1(X;)=/kpH`$ (where possibly $`H=0`$) satisfying:
* the decomposition is a decomposition of the torsion linking pairing on $`\mathrm{H}_1(X)`$: i.e. if $`\alpha `$ is a generator of the $`/kp`$-summand and $`\beta _1,\beta _2,\mathrm{},\beta _j`$ is a torsion basis for $`H`$, then $`link(\alpha ,\alpha )=m/n`$ for some $`m`$ relatively prime to $`n`$ and $`link(\alpha ,\beta _i)=0`$ for all $`i`$
* any generator of the $`/kp`$-summand maps to a generator in $`\pi _1(X)/\varphi _{}(\pi _1(\stackrel{~}{X}))`$ under the quotient map $`\mathrm{H}_1(X)\pi _1(X)/\varphi _{}(\pi _1(\stackrel{~}{X}))`$
* $`H`$ maps to 0 in $`\pi _1(X)/\varphi _{}(\pi _1(\stackrel{~}{X}))`$.
In this paper, we provide a Dehn surgery description for torsion-split regular $`k`$-fold cyclic covering space pairs $`(X,\stackrel{~}{X})`$ with base space a rational homology sphere. Specifically, we prove the following:
Let $`K`$ be a knot in $`S^3`$, and let $`L=(L_1,L_2,\mathrm{},L_n)`$ be a link in $`S^3`$. Assume $`K`$ bounds a Seifert surface $`\mathrm{\Sigma }`$ with the property that there exist Seifert surfaces $`\mathrm{\Sigma }_1,\mathrm{\Sigma }_2,\mathrm{},\mathrm{\Sigma }_n`$ for $`L_1,L_2,\mathrm{},L_n`$ disjoint from a neighborhood of $`\mathrm{\Sigma }`$. Let $`p`$, $`q`$, and $`k`$ be integers with $`1q`$, $`1|p|`$, and $`k>1`$. Assume $`kp`$ and $`q`$ are relatively prime. Let $`M`$ be the 3-manifold obtained by $`kp/q`$-Dehn surgery on $`K`$ in $`S^3`$ followed by surgery on $`L`$ with surgery coefficients $`I=(i_1,i_2,\mathrm{},i_n)`$, where $`i_j=\pm 1`$ for each $`j`$.
Note that for each $`j=1,2,\mathrm{},n`$, the component $`L_j`$ of $`L`$ has $`k`$ disjoint lifts in the $`k`$-fold branched cyclic cover of $`S^3`$ branched along $`K`$, since $`(K,L_j)`$ is a boundary link. Let $`\stackrel{~}{M}`$ be the 3-manifold obtained by $`p/q`$-Dehn surgery on the lift of $`K`$ to the $`k`$-fold branched cyclic cover of $`S^3`$ branched along $`K`$, followed by surgery on the $`k`$ lifts of $`L`$ with surgery coefficient $`i_j`$ for every lift of $`L_j`$.
###### Theorem 1.2
$`\stackrel{~}{M}`$ is a regular $`k`$-fold cyclic covering space of $`M`$.
Call $`(K,L,k,p,q,I)`$ a pairwise Dehn surgery description for $`(M,\stackrel{~}{M})`$. Then we also have
###### Theorem 1.3
Let $`(X,\stackrel{~}{X})`$ be a torsion-split regular $`k`$-fold cyclic covering space pair with base space $`X`$ a rational homology sphere. Then $`(X,\stackrel{~}{X})`$ has a pairwise Dehn surgery description.
Thus, every torsion-split regular $`k`$-fold cyclic covering space pair over a rational homology sphere has a pairwise Dehn surgery description. This, then, may be used for computing 3-manifold invariants for the pair. This should be a necessary first step in drawing general conclusions about the relationships between invariants for $`X`$ and those for $`\stackrel{~}{X}`$.
The paper is outlined as follows. In Section 2, we show that the manifolds generated by the pairwise Dehn surgery description $`(K,L,k,p,q,I)`$ are a regular $`k`$-fold cyclic covering space pair. In Section 3, we show that all torsion-split regular $`k`$-fold cyclic covering space pairs over rational homology spheres arise in this way. In Section 4, as an application, we compute Casson invariants for the pair $`(X,\stackrel{~}{X})`$.
In Section 5, we discuss the implications of our work for generating examples of finite and cyclic Dehn fillings of 3-manifolds with toral boundary. Specifically, let $`K`$ be a knot homologous to 0 in a closed, connected, oriented 3-manifold $`Y`$, and let $`X`$ be the result of $`p/q`$-Dehn surgery on $`K`$ in $`Y`$. We show that if $`\pi _1(X)`$ is finite (resp. cyclic), then so too is the fundamental group of the 3-manifold obtained by $`r/q`$-Dehn surgery on the lift of $`K`$ to the $`p/r`$-fold cyclic branched cover of $`Y`$ branched along $`K`$, where $`r`$ is any positive integer dividing $`p`$. Thus, one example of a finite or cyclic Dehn filling may generate many more such examples. Moreover, this restricts the possible number of finite and cyclic surgery slopes in the cyclic covering spaces of $`Ynbhd(K)`$.
Finally, in the appendix, we demonstrate some relationships between certain Alexander polynomials which are required in Section 4.
We remark that the construction presented here may generalize to pairs $`(X,\stackrel{~}{X})`$ for which $`X`$ is not a rational homology sphere. This theory, together with applications, will be the topic of a future paper.
I am indebted to Andrew Clifford for his help throughtout the writing process and to Steve Boyer for his corrections to an early version of this paper. I also thank Nancy Hingston and Andy Nicas for their helpful comments.
## 2 Generating covering space pairs
Retain all notation from Section 1. Further, denote by $`\overline{S}_K^3`$ the $`k`$-fold cyclic covering space of $`S^3`$ branched over $`K`$. Let $`N`$ be the closed 3-manifold resulting from $`kp/q`$-Dehn surgery on $`K`$ in $`S^3`$, and let $`\stackrel{~}{N}`$ denote the 3-manifold which is the result of $`p/q`$-Dehn surgery on the lift of $`K`$ to $`\overline{S}_K^3`$. Finally, if $`Y`$ is any 3-manifold and $`K`$ is a knot in $`Y`$, we abuse notation and denote by $`YK`$ the compact 3-manifold formed by removing an open tubular neighborhood of $`K`$.
Before proving Theorem 1.2, we prove the following:
###### Lemma 2.1
$`\stackrel{~}{N}`$ is a regular $`k`$-fold cyclic covering space of $`N`$.
Proof of Lemma 2.1: Let $`\overline{K}`$ denote the lift of $`K`$ to $`\overline{S}_K^3`$, and let $`\overline{m}`$ denote its meridian. Let $`m`$ denote the meridian of $`K`$. Let $`l`$ denote the preferred longitude of $`K`$ on the boundary of $`S^3K`$, and let $`\overline{l}`$ be a longitude in the boundary of $`\overline{S}_K^3\overline{K}`$ which lifts $`l`$. Then $`\stackrel{~}{N}`$ is obtained by gluing a solid torus to $`\overline{S}^3\overline{K}`$, sending a meridian of the solid torus to the simple closed curve $`\overline{m}^p\overline{l}^q`$. Similarly, $`N`$ is obtained by gluing a solid torus to $`S^3K`$, sending a meridian of the solid torus to $`m^{kp}l^q`$.
Let $`\varphi `$ denote the regular $`k`$-fold cyclic covering map of the knot complements. We wish to show that $`\varphi `$ extends to a regular $`k`$-fold cyclic (unbranched) covering map $`\stackrel{~}{N}N`$.
Note that $`\varphi `$ restricts to a regular $`k`$-fold cyclic covering map from the boundary of $`\overline{S}^3\overline{K}`$ to $`S^3K`$. Thus, the map from the boundary of the solid torus in $`\stackrel{~}{N}`$ to the boundary of the solid torus in $`N`$ is a regular $`k`$-fold cyclic covering map.
Now a $`k`$-fold cyclic covering map from the boundary of one solid torus to the boundary of another solid torus extends to a regular $`k`$-fold cyclic covering map from the solid torus to the solid torus if and only if the meridian of the solid torus in the initial solid torus is taken to the meridian of the final solid torus. In our case, the meridian of the solid torus in $`\stackrel{~}{N}`$ is $`\overline{m}^p\overline{l}^q`$, and the meridian of the solid torus in $`N`$ is $`m^{kp}l^q`$. But $`\varphi `$ takes $`\overline{m}`$ to $`m^k`$ and $`\overline{l}`$ to $`l`$. Hence, $`\varphi (\overline{m}^p\overline{l}^q)=m^{kp}l^q`$. It follows that $`\varphi `$ takes a meridian of the solid torus in $`\stackrel{~}{N}`$ to a meridian of the solid torus in $`N`$, and therefore $`\varphi `$ extends to a covering map $`\stackrel{~}{N}N`$. $`\mathrm{}`$
We now prove Theorem 1.2.
Proof of Theorem 1.2: We know that $`\stackrel{~}{N}`$ is a regular $`k`$-fold cyclic covering space of $`N`$ by Lemma 2.1. Let $`\varphi `$ denote the covering map. We abuse notation and denote the image of $`L_j`$ in $`N`$ by $`L_j`$ for $`j=1,2,\mathrm{},n`$. Let $`\overline{L}_j`$ denote the inverse image under $`\varphi `$ of $`L_j`$ for each $`j`$. Since $`(K,L_j)`$ is a boundary link in $`S^3`$ for each $`j`$, we see that $`\overline{L}_j`$ consists of $`k`$ disjoint simple closed curves in $`\stackrel{~}{N}`$. Choose a Seifert surface $`\mathrm{\Sigma }`$ for $`K`$ and Seifert surfaces $`\mathrm{\Sigma }_j`$ for $`L_j`$ in $`S^3`$ such that $`\mathrm{\Sigma }\mathrm{\Sigma }_j=\mathrm{}`$. Then the image of $`\mathrm{\Sigma }_j`$ in $`N`$ is a Seifert surface for $`L_j`$ in $`N`$, and it lifts to $`k`$ Seifert surfaces $`\overline{\mathrm{\Sigma }}_{j,1},\overline{\mathrm{\Sigma }}_{j,2},\mathrm{},\overline{\mathrm{\Sigma }}_{j,k}`$ for the $`k`$ lifts $`\overline{L}_{j,1},\overline{L}_{j,2},\mathrm{},\overline{L}_{j,k}`$ of $`L_j`$ in $`\overline{L}_j`$.
Recall that $`\overline{S}_K^3`$ may be explicitly constructed according to the following outline: let $`Y^0=\mathrm{\Sigma }\times (1,1)`$ be an open bicollar of $`\mathrm{\Sigma }`$, and let $`Y=Y^0/(K\times (1,1)K)`$. Let $`Y^{}`$ be the manifold obtained by removing $`K`$ from $`Y`$. Glue $`k`$ copies of $`S^3\mathrm{\Sigma }`$ together along $`k`$ copies of $`Y^{}`$, alternating copies of $`S^3\mathrm{\Sigma }`$ with copies of $`Y^{}`$. Finally, glue $`K`$ back in to compactify. For a precise description of this construction, see \[R, pp 128 - 131 and pp 297-298\].
Now since $`(K,L_j)`$ is a boundary link in $`S^3`$, we see that the $`k`$ lifts of $`\mathrm{\Sigma }_j`$ to $`\overline{S}_K^3`$ are contained in the $`k`$ disjoint copies of $`S^3Y`$. It follows that these lifts of $`\mathrm{\Sigma }_j`$ are disjoint. Then clearly the Seifert surfaces $`\overline{\mathrm{\Sigma }}_{j,1},\overline{\mathrm{\Sigma }}_{j,2},\mathrm{},\overline{\mathrm{\Sigma }}_{j,k}`$ in $`\stackrel{~}{N}`$ are also disjoint.
Now the Dehn surgery on $`L_j`$ may be carried out in a neighborhood $`Z_j`$ of $`\mathrm{\Sigma }_j`$. Moreover, if we take $`Z_j`$ to be sufficiently small, then the inverse image $`\varphi ^1(Z_j)`$ consists of $`k`$ disjoint copies $`\overline{Z}_{j,i}`$ of $`Z_j`$ which are neighborhoods of the $`\overline{\mathrm{\Sigma }}_{j,i}`$. Clearly the Dehn surgery on $`L_j`$ in $`Z_j`$ induces an identical Dehn surgery on $`\overline{L}_{j,i}`$ in $`\overline{Z}_{j,i}.`$ Therefore every point of $`M`$ has $`k`$ distinct inverse images in $`\stackrel{~}{M}`$, each with a neighborhood which is carried homeomorphically to a neighborhood of the point in $`M`$. Thus, $`\stackrel{~}{M}`$ is a $`k`$-fold covering space of $`M`$.
Finally, note that for each $`j`$, the automorphism group $`/k`$ of $`\varphi :\stackrel{~}{N}N`$ cyclically permutes the $`k`$ disjoint copies of the Seifert surface for $`L_j`$ in $`\stackrel{~}{N}`$. Clearly, then, the automorphism group of the covering space $`\stackrel{~}{M}M`$ is also $`/k`$, and the images of the $`k`$ lifts of $`\overline{L}_{j,i}`$ after Dehn surgery are permuted by the automorphism group. Since the covering space $`\stackrel{~}{N}N`$ was regular, it follows that $`\stackrel{~}{M}M`$ is regular. $`\mathrm{}`$
We remark that the covering $`\stackrel{~}{M}M`$ is clearly torsion-split by construction.
## 3 Completeness of the construction
In this section, we show that the construction described in Section 1 is complete in that it generates all torsion-split regular $`k`$-fold cyclic covering space pairs $`(X,\stackrel{~}{X})`$ over rational homology spheres. We prove:
Theorem 1.3 Let $`(X,\stackrel{~}{X})`$ be a torsion-split regular $`k`$-fold cyclic covering space pair with base space $`X`$ a rational homology sphere. Then $`(X,\stackrel{~}{X})`$ has a pairwise Dehn surgery description.
We first prove the following lemma, which is a straightforward generalization of a lemma of S. Boyer and D. Lines \[BL\].
###### Lemma 3.1
Let $`W`$ be a rational homology sphere with $`\mathrm{H}_1(W;)=/nH`$ for some finite abelian group $`H`$. Assume further that the homology decomposition arises as a decomposition of the torsion linking pairing on $`\mathrm{H}_1(W)`$. Then there is a 3-manifold $`V`$ with $`\mathrm{H}_1(V;)=H`$, a knot $`𝒦`$ homologous to 0 in $`V`$, and an integer $`m`$ such that $`W`$ is the result of $`n/m`$-Dehn surgery on $`𝒦`$ in $`V`$.
Proof of Lemma 3.1: Let $`\alpha `$ be a generator of the $`/n`$ summand of $`\mathrm{H}_1(W;)`$. Represent $`\alpha `$ by a curve $`C`$ in $`W`$.
Note that the torsion subgroup of $`\mathrm{H}_1(WC;)`$ is just $`H`$. To see this, note that by assumption, $`link(\alpha ,\alpha )=t/n`$ for some integer $`t`$ relatively prime to $`n`$, where $`link(\text{\_},\text{\_})`$ is the torsion linking pairing on $`\mathrm{H}_1(W)`$. Therefore $`C`$ intersects the surface with boundary $`nC`$, and $`nC`$ does not bound in $`WC`$. On the other hand, there exist generators $`\beta _i`$ of $`H`$ such that $`link(\alpha ,\beta _i)=0`$. It follows that $`\beta _i`$ has the same finite order in $`\mathrm{H}(WC)`$ as in $`\mathrm{H}(W)`$.
Let $`T(C)`$ be a tubular neighborhood of $`C`$. Then $`\mathrm{H}_2(W,WT(C);)`$ is infinite cyclic with generator a meridional disk of $`T(C)`$. From the exact sequence
$$0\mathrm{H}_2(W,WT(C);)\mathrm{H}_1(WT(C);)\mathrm{H}_1(W;)0$$
it follows that $`\mathrm{H}_1(WT(C);)H`$.
Let $`C^{}`$ be a simple closed curve on $`T(C)`$ generating the infinite cyclic summand of $`\mathrm{H}_1(WT(C);)`$. Attach a solid torus to $`WT(C)`$ sending the meridian to $`C^{}`$; call the resulting manifold $`V`$. Let $`𝒦`$ denote the core of the surgery torus. Then $`\mathrm{H}_1(V;)H`$, and $`W`$ is the result of $`n/m`$-Dehn surgery on $`𝒦`$ for some integer $`m`$. Moreover, $`𝒦`$ is homologous to 0 in $`V`$, since all curves on $`T(C)`$ represent 0 in $`H`$. $`\mathrm{}`$
We now prove the theorem.
Proof of Theorem 1.3: Let $`\stackrel{~}{X}X`$ be a torsion-split regular $`k`$-fold cyclic covering space with $`X`$ a rational homology sphere. Let $`\mathrm{H}_1(X;)=/kpH`$ be a homology decomposition for $`X`$ obeying properties i) - iii) of Definition 1.1. We may apply Lemma 3.1 to find a 3-manifold $`V`$ with $`\mathrm{H}_1(V;)=H`$, a knot $`𝒦`$ homologous to 0 in $`V`$, and an integer $`q`$ such that $`X`$ is the result of $`kp/q`$-Dehn surgery on $`𝒦`$ in $`V`$. Let $`\mathrm{\Sigma }`$ be a Seifert surface for $`𝒦`$ in $`V`$.
It is well-known that there exists a link $`=(_1,_2,\mathrm{},_n)`$ in $`V`$ such that $`S^3`$ is the result of surgery on $``$. Let $`L=(L_1,L_2,\mathrm{},L_n)`$ be the image of $``$ in $`S^3`$. Then $`V`$ may be obtained from $`S^3`$ by surgeries on the components of $`L`$. Moreover, we may choose $``$ in such a way that the surgery coefficients for the components of $`L`$ are all $`\pm 1`$.
Choose $`\alpha _1,\alpha _2,\mathrm{},\alpha _{2g}`$ a collection of simple closed curves on $`\mathrm{\Sigma }`$ representing a homology basis for $`\mathrm{\Sigma }`$. Isotope $``$ without changing any crossings of $``$ so that the linking number of $`_i`$ and $`\alpha _j`$ is 0 for any $`i=1,2,\mathrm{},n`$ and any $`j=1,2,\mathrm{},2g`$ and so that $`_i\mathrm{\Sigma }=\mathrm{}`$. We continue to denote by $``$ and $`L`$ the images of $``$ and $`L`$ under the isotopy. Let $`K`$ be the image of $`𝒦`$ in $`S^3`$ after surgery on $``$. Abusing notation yet again, we denote by $`\mathrm{\Sigma }`$ and $`\alpha _j`$ the image of $`\mathrm{\Sigma }`$ and $`\alpha _j`$ in $`S^3`$. Note that the linking number of $`L_i`$ and $`\alpha _j`$ is 0 for any $`i=1,2,\mathrm{},n`$ and any $`j=1,2,\mathrm{},2g`$ by Lemma A.5. Furthermore $`L_i`$ does not meet $`\mathrm{\Sigma }`$.
We require the following
###### Lemma 3.2
Let $`C`$ be a knot in $`S^3`$, and let $`S`$ be a Seifert surface for $`C`$. Let $`x_1,x_2,\mathrm{},x_{2g}`$ be a collection of curves on $`S`$ representing a basis for $`\mathrm{H}_1(S;)`$. Let $`D`$ be a knot in $`S^3`$ such that $`DS=\mathrm{}`$ and $`lk(x_i,D)=0`$ for $`i=1,2,\mathrm{},2g`$. Then $`D`$ bounds a Seifert surface disjoint from $`S`$.
Proof of lemma: Let $`S^{}`$ be a Seifert surface for $`D`$ meeting $`S`$ transversely. If $`S^{}S=\mathrm{}`$, we are done. Otherwise, $`S^{}S`$ is a collection $`y_1,y_2,\mathrm{},y_m`$ of oriented simple closed curves. Now the homology class represented by $`y_1+y_2+\mathrm{}+y_m`$ in $`\mathrm{H}_1(S;)`$ must be 0. For if this class were non-zero, then there would be a curve $`x_i`$ on $`S`$ whose oriented intersection number with $`y_1+y_2+\mathrm{}+y_m`$ was non-zero. Then $`x_i`$ would have non-zero intersection number with $`S^{}`$ and hence would have non-zero linking number with $`D`$.
Now remove from $`S^{}`$ the components of $`S^{}S`$ which do not meet $`D`$. The resulting surface $`S^{\prime \prime }`$ has boundary $`y_1y_2\mathrm{}y_mD`$. Since the sum of the $`y_i`$’s represents 0 in $`\mathrm{H}_1(S;)`$, the curves $`y_i`$ cobound a collection of subsurfaces of $`S`$. Then we may glue a collection of parallel copies of these surfaces to the appropriate boundary components of $`S^{\prime \prime }`$ to form a new two-sided surface $`S^{\prime \prime \prime }`$ with boundary $`D`$. Pushing the parallel copies of subsurfaces of $`S`$ apart and away from $`S`$, we will find that $`S^{\prime \prime \prime }`$ is embedded and disjoint from $`S`$. $`\mathrm{}`$
Now returning to the proof of the theorem:
Applying the lemma to each of the link components $`L_i`$, we may find Seifert surfaces $`\mathrm{\Sigma }_1,\mathrm{\Sigma }_2,\mathrm{},\mathrm{\Sigma }_n`$ for $`L_1,L_2,\mathrm{},L_n`$, respectively, such that $`\mathrm{\Sigma }\mathrm{\Sigma }_i=\mathrm{}`$ for $`i=1,2,\mathrm{},n`$. We show
Claim: $`(K,L,k,p,q,I)`$ is a pairwise Dehn surgery description for $`(X,\stackrel{~}{X})`$.
Proof of Claim: It is clear from the construction that $`(K,L,k,p,q,I)`$ is a pairwise Dehn surgery description for some pair $`(X,X^{})`$ with base space $`X`$. Let $`\psi :X^{}X`$ be the projection map. Then $`\psi _{}(\pi _1(X^{}))`$ is the kernel of the homomorphism
$$\pi _1(X)\mathrm{H}_1(X;)/k,$$
where the latter homomorphism is the projection
$$\mathrm{H}_1(X;)\mathrm{H}_1(X;)/H=/kp/k.$$
But this kernel is precisely $`\varphi _{}(\pi _1(\stackrel{~}{X}))`$. Hence $`X^{}=\stackrel{~}{X}`$. $`\mathrm{}`$
## 4 Casson-Walker invariants for pairs
In 1985, Andrew Casson defined an invariant $`\lambda `$ for integral homology 3-spheres. Roughly, this invariant counts the signed equivalence classes of SU(2)-representations of the fundamental group of the 3-manifold. This invariant was extended to an invariant for oriented rational homology 3-spheres by Kevin Walker in \[W\], and Christine Lescop derived a combinatorial formula extending the invariant to arbitrary closed, oriented 3-manifolds in \[L\].
A number of mathematicians have explored the Casson-Walker invariant for branched covers of links in $`S^3`$. David Mullins computed the invariant for 2-fold branched covers in the case when the 2-fold branched cover is a rational homology sphere in \[M\]. More general results for $`k>2`$ may be found in \[GR\]. The invariants for $`n`$-fold branched covers $`\overline{S}_K^3`$ of particular families of knots have been computed by J. Hoste, A. Davidow, and K. Ishibe in \[H\], \[D\], and \[I\]. Garoufalidis generalizes several of these formulas in \[G\].
We study the Casson-Walker invariants of pairs $`(X,\stackrel{~}{X})`$. As in the early sections of the paper, we assume $`\stackrel{~}{X}X`$ is a torsion-split regular $`k`$-fold cyclic covering over a rational homology sphere. In this section, we assume further that $`\stackrel{~}{X}`$ is a rational homology sphere and that $`(X,\stackrel{~}{X})`$ has a Dehn surgery description $`(K,L,k,p,q,I)`$ with $`\overline{S}_K^3`$ a rational homology sphere.
In what follows, for any pair of non-zero integers $`x`$ and $`y`$ which are relatively prime, let $`s(x,y)`$ denote the Dedekind sum defined by
$$s(x,y)=sign(y)\underset{j=1}{\overset{|x|}{}}((j/y))((jx/y))$$
where
$$((z))=\{\begin{array}{cc}0\hfill & z\hfill \\ z[z]1/2\hfill & \text{else}\hfill \end{array}$$
For a knot $`C`$ in a rational homology sphere, let $`\mathrm{\Delta }_C`$ denote the Alexander polynomial of C, normalized so that it is symmetric in $`t^{1/2}`$ and $`t^{1/2}`$ and so that $`\mathrm{\Delta }_C(1)=1`$. We show
###### Theorem 4.1
Let $`(K,L,k,p,q,I)`$ be a pairwise Dehn surgery description for $`(X,\stackrel{~}{X})`$. Then
$$\lambda (\stackrel{~}{X})=k\lambda (X)+q/p(\mathrm{\Delta }_{\overline{K}}^{^{\prime \prime }}(1)\mathrm{\Delta }_K^{^{\prime \prime }}(1))ks(q,kp)+s(q,p)+\lambda (\overline{S}_K^3).$$
Here, $`\overline{S}_K^3`$ denotes the $`k`$-fold branched cyclic cover of $`S^3`$ branched along $`K`$ and $`\overline{K}`$ denotes the lift of $`K`$ to $`\overline{S}_K^3`$, as above.
Note that $`\mathrm{\Delta }_{\overline{K}}`$ can be computed from $`\mathrm{\Delta }_K`$. This relationship is described in the appendix. We remark further that in the case $`k=2`$, the invariant $`\lambda (\overline{S}_K^3)`$ can be computed whenever $`\stackrel{~}{X}`$ is a rational homology sphere using the work of Mullins. For $`k>2`$, the results of Garoufalidis and Rozansky apply. For certain families of knots, the invariant $`\lambda (\overline{S}_K^3)`$ can be computed for any value of $`k`$ using the results of Hoste, Davidow, Ishibe, and Garoufalidis.
Proof of Theorem 4.1: Retain all notation from previous sections. We begin by noting that
$$\lambda (N)=(q/kp)\mathrm{\Delta }_K^{^{\prime \prime }}(1)+s(q,kp)$$
(1)
and
$$\lambda (\stackrel{~}{N})=\lambda (\overline{S}_K^3)+(q/p)\mathrm{\Delta }_{\overline{K}}^{^{\prime \prime }}(1)+s(q,p)$$
(2)
by Proposition 6.2 of \[W\], since $`K`$ is a knot in $`S^3`$ and since $`\overline{K}`$ is homologous to 0 in $`\overline{S}_K^3`$.
Now $`X`$ is obtained from $`N`$ by $`I`$-surgery on $`L`$. We know that $`\lambda (X)\lambda (N)`$ may be obtained using the surgery formulae developed by Walker in \[W\]. These formulae depend on the coefficients $`I`$, as well as the link $`L`$ and the Alexander polynomials of the components of L. Similarly $`\stackrel{~}{X}`$ is obtained from $`\stackrel{~}{N}`$ by $`I`$-surgery on each lift $`\overline{L}_j`$ of $`L`$ in $`\stackrel{~}{N}`$.
Now by Proposition A.7 in the appendix, the Alexander polyomial of each component of $`\overline{L}`$ is equal to that of the corresponding component of $`L`$. Therefore, since a neighborhood of each component $`\overline{L}_{j,i}`$ of $`\overline{L}_j`$ in $`\stackrel{~}{N}`$ is carried homeomorphically onto a neighborhood of $`L_j`$ in $`N`$ and the coefficients $`I`$ correspond, it is clear that for any $`j`$, the Casson-Walker invariant of the manifold $`\stackrel{~}{N}_j`$ obtained by doing $`I`$-surgery on $`\overline{L}_j`$ in $`\stackrel{~}{N}`$ obeys the formula $`\lambda (\stackrel{~}{N}_j)\lambda (\stackrel{~}{N})=\lambda (X)\lambda (N)`$. But the lifts $`\overline{\mathrm{\Sigma }}_{j,i}`$ of $`\mathrm{\Sigma }_j`$ are contained in disjoint copies of $`S^3Y`$ in $`\stackrel{~}{N}`$, as shown in the proof of Theorem 1.2, so $`\overline{\mathrm{\Sigma }}_{j,i}\overline{\mathrm{\Sigma }}_{l,m}=\mathrm{}`$ if $`im`$. Then also $`\lambda (\stackrel{~}{N_{jk}})\lambda (\stackrel{~}{N}_j)=\lambda (\stackrel{~}{N}_j)\lambda (N)`$, where $`\stackrel{~}{N}_{jk}`$ is the manifold obtained by doing $`I`$-surgery on the image of $`\overline{L}_k`$ in $`\stackrel{~}{N}_j`$. It follows that
$$\lambda (\stackrel{~}{X})\lambda (\stackrel{~}{N})=k(\lambda (X)\lambda (N)).$$
(3)
The theorem follows after suitably combining equations (1) - (3). $`\mathrm{}`$
We remark that analogous results for the generalizations of $`\lambda `$ counting representations in SO(3), U(2), Spin(4), and SO(4) may be immediately obtained using the results of \[C\].
## 5 Cyclic and finite Dehn surgeries
In recent years, many new and exciting results have come to light concerning Dehn fillings of 3-manifolds with toral boundary. Among these are the Cyclic Surgery Theorem of Culler, Gordon, Luecke, and Shalen \[CGLS\] and the work of Boyer and Zhang concerning finite fillings \[BZ1\] and \[BZ2\].
We review some basic definitions.
###### Definition 5.1
Let $`K`$ be a knot homologous to 0 in an oriented 3-manifold $`Y`$. A slope of $`K`$ is the unoriented isotopy class of a non-trivial simple closed curve in $`(YK)`$. The distance between two slopes is their geometric intersection number.
Recall that for any knot homologous to 0 in an oriented 3-manifold, the set of slopes of the knot is canonically isomorphic to $`\overline{}`$. Thus, we denote by $`p/q`$ the slope which is the isotopy class of a curve which is $`p`$ times a meridian plus $`q`$ times a longitude.
###### Definition 5.2
With $`K`$ and $`Y`$ as above, suppose that $`p/q`$ is a slope of $`K`$ such that the manifold $`X`$ which is the result of $`p/q`$ Dehn surgery on $`K`$ in $`Y`$ has finite fundamental group. Then we call $`p/q`$ a finite surgery slope of $`K`$. Similarly, if $`X`$ has cyclic fundamental group, we call $`p/q`$ a cyclic surgery slope of $`K`$.
The results of Culler, Gordon, Luecke, and Shalen and of Boyer and Zhang state that for most irreducible knot complements $`YK`$ as above, there are at most 3 cyclic surgery slopes of $`K`$ and at most 6 finite surgery slopes. Moreover these slopes may be distance at most 1 apart in the cyclic case and at most 5 apart in the finite case. A detailed survey of work in this area may be found in \[B\].
Here, we note that given a knot $`K`$ homologous to 0 in a closed, oriented 3-manifold $`Y`$ and a finite or cyclic surgery slope of $`K`$, our work leads to explicit examples of more such fillings. Specifically, let $`\overline{Y}_K`$ denote the $`k`$-fold branched cyclic cover $`Y`$ branched along $`K`$, and let $`\overline{K}`$ denote the lift of $`K`$ to $`\overline{Y}_K`$. We show
###### Theorem 5.3
Let $`K`$ be a knot homologous to 0 in a closed oriented 3-manifold $`Y`$. Let $`k`$, $`p`$, and $`q`$ be integers with $`k>1`$, $`|p|1`$, and $`q1`$. Then $`kp/q`$ is a finite surgery slope of $`K`$ if and only if $`p/q`$ is a finite surgery slope of $`\overline{K}`$ in $`\overline{Y}_K`$. Moreover if $`kp/q`$ is a cyclic surgery slope of $`K`$, then $`p/q`$ is a cyclic surgery slope of $`\overline{K}`$. Finally, if $`p/q`$ is a cyclic surgery slope of $`\overline{K}`$ and $`p1`$, then $`kp/q`$ is a cyclic surgery slope of $`K`$.
Proof: Let $`X`$ denote the manifold resulting from $`kp/q`$-Dehn surgery on $`K`$, and let $`\stackrel{~}{X}`$ denote the result of $`p/q`$-Dehn surgery on $`\overline{K}`$. As in the proof of Theorem 1.3, we may find a knot $`C`$ and a link $`L=(L_1,L_2,\mathrm{},L_n)`$ in $`S^3`$ satisfying
* there exist Seifert surfaces $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }_1,\mathrm{\Sigma }_2,\mathrm{},\mathrm{\Sigma }_n`$ for $`C`$ and $`L_1,L_2,\mathrm{},L_n`$, respectively, with $`\mathrm{\Sigma }\mathrm{\Sigma }_j=\mathrm{}`$ for $`j=1,2,\mathrm{},n`$
* $`Y`$ is the result of $`I`$-surgery on $`L`$ for some $`I=(i_1,i_2,\mathrm{},i_n)`$ with $`i_j=\pm 1`$
* $`K`$ is the image of $`C`$ in $`Y`$.
(Note that none of these steps requires $`Y`$ to be a rational homology sphere.)
Now $`(C,L,k,p,q,I)`$ is a Dehn surgery description for some pair $`(Z,\stackrel{~}{Z})`$. But since $`(C,L_i)`$ is a boundary link for $`i=1,2,\mathrm{},n`$, it is clear that $`kp/q`$-surgery on $`C`$ followed by $`I`$-surgery on the image of $`L`$ yields the same manifold as $`I`$-surgery on $`L`$ followed by $`kp/q`$-surgery on the image of $`C`$. Hence $`XZ`$. Also, since $`(C,L_i)`$ is a boundary link, we see that $`\overline{Y}_K`$ may be obtained by $`I`$-surgery on the $`k`$ lifts of $`L`$ to $`\overline{S}_C^3`$. Moreover $`p/q`$-surgery on $`\overline{C}`$ in $`\overline{S}_C^3`$ followed by $`I`$-surgery on the images of each of the $`k`$ lifts of $`L`$ to $`\overline{S}_C^3`$ yields the same manifold as $`I`$-surgery on each of the lifts of $`L`$ to $`\overline{S}_C^3`$ followed by $`p/q`$-surgery on the image $`\overline{K}`$ of $`\overline{C}`$ in $`\overline{Y}_K`$. Hence $`\stackrel{~}{X}\stackrel{~}{Z}`$.
Now it follows from Theorem 1.2 that $`\stackrel{~}{X}`$ is a regular $`k`$-fold cyclic covering space of $`X`$. Therefore the fundamental group of $`\stackrel{~}{X}`$ is an index $`k`$ subgroup of that of $`X`$. Then clearly $`\pi _1(X)`$ is finite if and only if $`\pi _1(\stackrel{~}{X})`$ is finite, and $`\pi _1(\stackrel{~}{X})`$ is cyclic if $`\pi _1(X)`$ is cyclic.
Finally, if $`\pi _1(\stackrel{~}{X})`$ is cyclic and $`p1`$, so $`\pi _1(\stackrel{~}{X})=/p`$, then $`\pi _1(X)=\mathrm{H}_1(X)=/kp`$, since $`\mathrm{H}_1(X;)`$ has a $`/kp`$ summand and $`\pi _1(\stackrel{~}{X})`$ is index $`k`$ in $`\pi _1(X)`$. This proves the theorem. $`\mathrm{}`$
In light of \[CGLS\], \[BZ1\], and \[BZ2\], then, finding a finite or cyclic surgery slope $`p/q`$ for some knot $`K`$ homologous to 0 in an oriented 3-manifold $`Y`$ severely restricts the set of possible finite and cyclic surgery slopes not only of $`K`$ itself, but also of lifts of $`K`$ to branched covers of $`Y`$ branched along $`K`$ of all orders dividing $`p`$. For example, it is a result of Fintushel and Stern \[FS\] that 18- and 19-surgeries on the $`(2,3,7)`$ pretzel knot yield lens spaces. It follows that 1-surgery on the lift of the $`(2,3,7)`$ pretzel knot to either the 18- or 19-fold branched cover of $`S^3`$ branched along the pretzel knot yields $`S^3`$. Further, we find that $`p`$-surgery on the lift of the knot to the $`(18/p)`$-fold branched cover of $`S^3`$ branched along the knot also yields a lens space for $`p=2`$, 3, 6, or 9.
Finally, we note that the results of \[CGLS\], \[BZ1\], and \[BZ2\] may be strengthened for manifolds of the form $`\overline{Y}_K\overline{K}`$. We have the following corollary, where $`Y`$, $`K`$, $`\overline{Y}_K`$, and $`\overline{K}`$ are defined as above:
###### Corollary 5.4
Suppose $`YK`$ is irreducible and is not a Seifert fibered space. Then $`\overline{K}`$ has at most one cyclic surgery slope $`p/q`$ with $`p1`$.
If also $`YK`$ is not a cable on the twisted $`I`$-bundle over the Klein bottle, then the distance between any two finite surgery slopes of $`\overline{K}`$ is at most $`5/k`$.
If $`YK`$ is also hyperbolic, then the distance between a cyclic surgery slope $`p/q`$ with $`p1`$ and any finite surgery slope is at most $`2/k`$. Moreover the distance between any two finite surgery slopes is at most $`3/k`$.
Proof: By Theorem 5.3, for any $`p`$ and $`q`$ with $`|p|1`$ and $`q1`$, we know that $`p/q`$ is a finite surgery slope of $`\overline{K}`$ if and only if $`kp/q`$ is a finite surgery slope of $`K`$ and that $`kp/q`$ is a cyclic surgery slope of $`K`$ if $`p/q`$ is a cyclic surgery slope of $`\overline{K}`$ with $`p1`$. The distance between $`kp^{}/q^{}`$ and $`kp/q`$ is $`k`$ times the distance between $`p^{}/q^{}`$ and $`p/q`$. The assertions follow from the restrictions on the distance between cyclic (resp. finite) surgery slopes $`kp/q`$ for $`K`$ imposed by \[CGLS\], \[BZ1\], and \[BZ2\]. $`\mathrm{}`$
## Appendix A Alexander polynomials in $`\overline{S}_K^3`$
In this section, we relate the Alexander polynomials $`\mathrm{\Delta }_{\overline{K}}`$ and $`\mathrm{\Delta }_{\overline{L}_{j,i}}`$ in $`\overline{S}_K^3`$ to $`\mathrm{\Delta }_K`$ and $`\mathrm{\Delta }_{L_j}`$, respectively. These results are used in Section 4.
We begin by recalling the definition of the Alexander polynomial of a knot in a rational homology sphere. Details can be found in \[M\].
Let $`C`$ be a knot in a rational homology sphere, and denote by $`Z`$ the complement of a tubular neighborhood of $`C`$. Further, denote by $`\stackrel{~}{Z}`$ the infinite cyclic cover of $`Z`$ determined by $`\pi _1(Z)\mathrm{H}_1(Z;)/torsion`$. Then $`\mathrm{H}_1(\stackrel{~}{Z};)`$ is a module over $`[]`$.
Now $`[]`$ is a principal ideal domain, so
$$\mathrm{H}_1(\stackrel{~}{Z};)[]/(p_1)[]/(p_2)\mathrm{}[]/(p_r).$$
We define the order of $`\mathrm{H}_1(\stackrel{~}{Z};)`$ to be the product ideal $`(p_1p_2\mathrm{}p_r)`$.
###### Definition A.1
An Alexander polynomial $`\mathrm{\Delta }_C(t)`$ of $`C`$ is any polynomial generating the order of $`\mathrm{H}_1(\stackrel{~}{Z};)`$. We also call $`\mathrm{\Delta }_C(t)`$ an Alexander polynomial of the $`[]`$-module $`\mathrm{H}_1(\stackrel{~}{Z};)`$.
This is well-defined up to multiplication by polynomials of the form $`ct^k`$, where $`c`$, where $`t`$ generates the deck transformations of the covering $`\stackrel{~}{Z}Z`$, and where $`k`$. Henceforth we write $`p(t)q(t)`$ if $`p(t)`$ and $`q(t)`$ are polynomials with $`p(t)=ct^kq(t)`$ for some $`c`$ and some $`k`$.
We first relate $`\mathrm{\Delta }_{\overline{K}}`$ to $`\mathrm{\Delta }_K`$. We prove the following theorem, which was pointed out to me by Steve Boyer. The proof offered is that of Boyer.
###### Proposition A.2
(Boyer)
$$\mathrm{\Delta }_{\overline{K}}(t^k)\underset{j=0}{\overset{k1}{}}\mathrm{\Delta }_K(\zeta ^jt)$$
where $`\zeta `$ is a primitive $`k`$th root of unity.
Proof (Boyer): Let $`A`$ be the matrix of multiplication by $`t`$ in $`\mathrm{H}_1(\stackrel{~}{Z};)`$ with respect to the $``$-vector space structure. We first show
###### Lemma A.3
(Boyer) $`\mathrm{\Delta }_K(t)|AtI|`$
Proof of lemma (Boyer): Since $`\mathrm{H}_1(\stackrel{~}{Z};)`$ decomposes as
$$[t,t^1]/p_1(t)[t,t^1]/p_2(t)\mathrm{}[t,t^1]/p_r(t),$$
we see that $`A=_{j=1}^rA_j`$, where $`A_j`$ is the matrix of multiplication by $`t`$ in $`[t,t^1]/p_j(t)`$. Thus it suffices to show that
$$|A_jtI|p_j(t).$$
But writing $`p_j(t)=b_0+b_1t+b_2t^2+\mathrm{}+b_{s1}t^{s1}+t^s`$, we see that $`[t,t^1]/p_j(t)`$ has basis $`1,t,t^2,\mathrm{},t^{s1}`$. Then
$$A_j=\left[\begin{array}{ccccc}0\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & b_0\hfill \\ 1\hfill & 0\hfill & \mathrm{}\hfill & 0\hfill & b_1\hfill \\ 0\hfill & 1\hfill & \mathrm{}\hfill & 0\hfill & b_2\hfill \\ \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill & \mathrm{}\hfill \\ 0\hfill & 0\hfill & \mathrm{}\hfill & 1\hfill & b_{s1}\hfill \end{array}\right]$$
It follows that $`|A_jtI|=(1)^sp_j(t)`$. $`\mathrm{}`$
Now define a new $`[t,t^1]`$-module structure on $`\mathrm{H}_1(\stackrel{~}{Z};)`$ with mutliplication $`tm=t^km`$. Clearly $`\mathrm{H}_1(\stackrel{~}{Z};)`$ with the new multiplication is again a finitely generated, torsion $`[t,t^1]`$-module and hence a finite dimensional $``$-vector space. We show
###### Lemma A.4
(Boyer) Let $`\mathrm{\Delta }_k(t)`$ denote an Alexander polynomial for $`\mathrm{H}_1(\stackrel{~}{Z};)`$ with the new $`[t,t^1]`$-module structure. Then
$$\mathrm{\Delta }_k(t^k)\underset{j=0}{\overset{k1}{}}\mathrm{\Delta }(\zeta ^jt)$$
Proof of lemma (Boyer): If $`A`$ is the matrix of multiplication by $`t`$ in $`\mathrm{H}_1(\stackrel{~}{Z};)`$ with the original $`[t,t^1]`$-module structure, then $`A^k`$ is the matrix of multiplication by $`t`$ in $`\mathrm{H}_1(\stackrel{~}{Z};)`$ with the new $`[t,t^1]`$-module structure. Therefore by the previous lemma,
$`\mathrm{\Delta }_k(t^k)`$ $``$ $`|A^kt^kI|`$ (4)
$`=`$ $`{\displaystyle \underset{j=0}{\overset{k1}{}}}|A\zeta ^jtI|`$ (5)
$`=`$ $`{\displaystyle \underset{j=0}{\overset{k1}{}}}\mathrm{\Delta }(\zeta ^jt)`$ (6)
This proves the lemma. $`\mathrm{}`$
But now note that an Alexander polynomial $`\mathrm{\Delta }_{\overline{K}}(t)`$ is an Alexander polynomial of $`\mathrm{H}_1(\stackrel{~}{Z};)`$ with the second $`[t,t^1]`$-module structure, while $`\mathrm{\Delta }_K(t)`$ is an Alexander polynomial of $`\mathrm{H}_1(\stackrel{~}{Z};)`$ with the original $`[t,t^1]`$-module structure. The theorem follows. $`\mathrm{}`$
Thus, renaming $`u=t^n`$, say, we obtain $`\mathrm{\Delta }_{\overline{K}}(u)`$ in terms of $`\mathrm{\Delta }_K`$.
This can be further simplified as follows: write
$$\mathrm{\Delta }_K(t)=c_0+c_1(t+t^1)+c_2(t^2+t^2)+\mathrm{}+c_n(t^n+t^n).$$
(7)
Applying Boyer’s theorem to $`(7)`$ and noting that all cross-terms cancel, we see that
$$\mathrm{\Delta }_{\overline{K}}(t^k)c_0^k+c_1^k\underset{j=0}{\overset{k1}{}}(\zeta ^jt+\zeta ^jt^1)+c_2^k\underset{j=0}{\overset{k1}{}}(\zeta ^{2j}t^2+\zeta ^{2j}t^2)+\mathrm{}+c_n^k\underset{j=0}{\overset{k1}{}}(\zeta ^{nj}t^n+\zeta ^{nj}t^n).$$
(8)
In particular, if $`k`$ is odd, note that $`_{j=0}^{k1}(\zeta ^{rj}t^r+\zeta ^{rj}t^r)=t^{rk}+t^{rk}`$ for any $`r`$. Then equation (8) becomes
$$\mathrm{\Delta }_{\overline{K}}(t^k)c_0^k+c_1^k(t^k+t^k)+c_2^k(t^{2k}+t^{2k})+\mathrm{}+c_n^k(t^{nk}+t^{nk})$$
(9)
and finally
$$\mathrm{\Delta }_{\overline{K}}(u)c_0^k+c_1^k(u+u^1)+c_2^k(u^2+u^2)+\mathrm{}+c_n^k(u^n+u^n).$$
(10)
We now turn to relating $`\mathrm{\Delta }_{\overline{L}_{j,i}}`$ to $`\mathrm{\Delta }_{L_j}`$. We return to the convention that $`\mathrm{\Delta }_C`$ is symmetric in $`t^{1/2}`$ and $`t^{1/2}`$ and that $`\mathrm{\Delta }_C(1)=1`$, so that $`\mathrm{\Delta }_C`$ is uniquely defined.
For the remainder of the paper, for any rational homology sphere $`M`$, let $`lk_M(\mathrm{\_},\mathrm{\_})`$ denote the linking number in $`M`$.
###### Lemma A.5
Given knots $`C`$ and $`C^{}`$ in a rational homology sphere $`M`$ such that $`lk_M(C,C^{})=0`$. Let $`\widehat{M}`$ denote the manifold resulting from $`p/q`$-Dehn surgery on $`C^{}`$ in $`M`$. Then $`lk_{\widehat{M}}(\widehat{C},\widehat{D})=lk_M(C,D)`$ for any knot $`D`$ in $`M`$, where $`\widehat{C}`$ and $`\widehat{D}`$ are the images of $`C`$ and $`D`$ in $`\widehat{M}`$.
Proof: Fix a knot $`D`$ in $`M`$. Suppose $`C`$ represents an element of order $`m`$ in $`\mathrm{H}_1(M;)`$. Since $`lk_M(C,C^{})=0`$, we may choose a two-sided surface $`\mathrm{\Sigma }`$ with boundary $`m`$ times $`C`$ which is disjoint from a neighborhood of $`C^{}`$ and which meets $`D`$ transversely. Then Dehn surgery on $`C^{}`$ does not affect a neighborhood of $`\mathrm{\Sigma }`$ and hence does not affect the oriented intersection of $`\mathrm{\Sigma }`$ and $`D`$. The assertion follows. $`\mathrm{}`$
We now show
###### Lemma A.6
Let $`C`$ and $`D`$ be knots in a rational homology sphere $`M`$. Suppose $`C`$ is homologous to 0 and bounds a Seifert surface $`\mathrm{\Sigma }`$ disjoint from $`D`$ satisfying $`lk_M(D,\alpha _j)=0`$ for $`j=1,2,\mathrm{},2g`$, where $`\alpha _1,\alpha _2,\mathrm{},\alpha _{2g}`$ is a collection of simple closed curves on $`\mathrm{\Sigma }`$ representing a basis of $`\mathrm{H}_1(\mathrm{\Sigma };)`$. Let $`M^{}`$ denote the manifold resulting from $`p/q`$-Dehn surgery on $`D`$ in $`M`$, and let $`C^{}`$ denote the image of $`C`$ in $`M^{}`$. Then the Alexander polynomial of $`C^{}`$ in $`M^{}`$ is equal to the Alexander polynomial of $`C`$ in $`M`$.
Proof: Let $`\mathrm{\Sigma }^{}`$ be the image of $`\mathrm{\Sigma }`$ in $`M^{}`$. Let $`\alpha _1^{},\alpha _2^{},\mathrm{},\alpha _{2g}^{}`$ denote the images of $`\alpha _1,\alpha _2,\mathrm{},\alpha _{2g}`$, respectively, in $`M^{}`$. The Alexander polynomial $`\mathrm{\Delta }_C`$ of $`C`$ in $`M`$ is the determinant of the matrix $`A`$ with entries $`a_{i,j}=lk_M(\alpha _i^+,\alpha _j)tlk_M(\alpha _i^{},\alpha _j)`$, while the Alexander polynomial $`\mathrm{\Delta }_C^{}`$ of $`C^{}`$ in $`M^{}`$ is the determinant of the matrix $`A^{}`$ with entries $`a_{i,j}^{}=lk_M^{}(\alpha _i^{}_{}{}^{}+,\alpha _j^{})tlk_M^{}(\alpha _i^{{}_{}{}^{}},\alpha _j^{})`$. (Here $`x^+`$ and $`x^{}`$ denote the plus- and minus-pushoffs of a simple closed curve $`x`$ on $`\mathrm{\Sigma }`$ or $`\mathrm{\Sigma }^{}`$.) (See \[W,Appendix B\], for example.)
But since $`lk_M(D,\alpha _j)=0`$ for $`j=1,2,\mathrm{},2g`$, we may apply Lemma A.5 to show that $`lk_M^{}(\alpha _i^{}_{}{}^{}+,\alpha _j^{})=lk_M(\alpha _i^+,\alpha _j)`$ and $`lk_M^{}(\alpha _i^{{}_{}{}^{}},\alpha _j^{})=lk_M(\alpha _i^{},\alpha _j)`$ for all $`i`$ and $`j`$. It follows that $`A^{}=A`$, and hence $`\mathrm{\Delta }_C^{}=\mathrm{\Delta }_C`$. $`\mathrm{}`$
###### Proposition A.7
Let $`(K,L)`$ be a boundary link in $`S^3`$, and let $`\overline{L}`$ be a lift of $`L`$ in the branched $`k`$-fold cover $`\overline{S}_K^3`$ for some $`k`$. The Alexander polynomial of $`\overline{L}`$ in $`\overline{S}_K^3`$ is equal to the Alexander polynomial of $`L`$ in $`S^3`$.
Proof: Choose a link $`C=(C_1,C_2,\mathrm{},C_n)`$ in $`S^3`$ and surgery coefficients $`r=(r_1,r_2,\mathrm{},r_n)`$ so that $`r`$-Dehn surgery on $`C`$ yields $`S^3`$ and so that the image $`K^{}`$ of $`K`$ is an unknot. (See \[R, Section 6D\], for example.) Isotoping $`C`$ as necessary, changing crossings of $`C`$ and $`L`$ but changing no crossings of $`CK`$, we may assume that $`L`$, $`K`$, and $`C_i`$ bound Seifert surfaces $`\mathrm{\Sigma }_L`$, $`\mathrm{\Sigma }_K`$, and $`\mathrm{\Sigma }_i`$, respectively, with $`\mathrm{\Sigma }_L\mathrm{\Sigma }_K=\mathrm{\Sigma }_L\mathrm{\Sigma }_i=\mathrm{}`$ for $`i=1,2,\mathrm{},n`$.
Let $`L^{}`$ denote the image of $`L`$ under $`r`$-surgery on $`C`$. Let $`\alpha _1,\alpha _2,\mathrm{},\alpha _{2g}`$ be a collection of simple closed curves on $`\mathrm{\Sigma }`$ representing a basis of $`\mathrm{H}_1(\mathrm{\Sigma };)`$. By Lemma A.5, the linking number of the image of any component $`C_i`$ with the image of any curve $`\alpha _j`$ is 0 at any stage of the sequence of surgeries on $`C_1,C_2,\mathrm{},C_n`$. Then we may apply Lemma A.6 repeatedly to show that the Alexander polynomial $`\mathrm{\Delta }_L^{}=\mathrm{\Delta }_L`$. Furthermore, the linking number of $`K^{}`$ and the image of any $`\alpha _j`$ in $`S^3`$ at the end of the surgery sequence is 0. Then applying Lemma 3.2, we find that $`(L^{},K^{})`$ and $`(L^{},C_i^{})`$ are boundary links, where the $`C_i^{}`$ are the images of the $`C_i`$ in $`S^3`$ after the surgery. In fact, since $`K^{}`$ is an unknot, we see that $`(L^{},K^{})`$ is splittable.
Now consider the $`k`$-fold branched cover of $`S^3`$ branched along the unknot $`K^{}`$, which is again $`S^3`$. Since $`(L^{},K^{})`$ is splittable, we see that $`L^{}`$ has $`k`$ lifts $`\overline{L}_1^{},\overline{L}_2^{},\mathrm{},\overline{L}_k^{}`$ lying in disjoint 3-balls, and the Alexander polynomials satisfy $`\mathrm{\Delta }_{\overline{L}_i^{}}=\mathrm{\Delta }_L^{}`$ for $`i=1,2,\mathrm{},k`$. Let $`\overline{C}^{}`$ denote the inverse image of $`C^{}`$ in $`\overline{S}_K^{}^3`$.
Finally, it is clear that surgery on $`\overline{C}^{}`$ with appropriate coefficients yields $`\overline{S}_K^3`$ and takes $`\overline{L}_i^{}`$ onto $`\overline{L}_i`$ for $`i=1,2,\mathrm{},k`$. The assertion follows. $`\mathrm{}`$
Department of Mathematics and Statistics
The College of New Jersey
Ewing, NJ 08628
USA
ccurtis@tcnj.edu
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# On the extension of the concept of thin shells to the Einstein-Cartan theory
## 1 Introduction
The purpose of this paper is to extend the notion of singular hypersurface existing in general relativity to the context of the Einstein-Cartan theory in a coherent way in order to include the effects of torsion originating in the spinning properties of the matter distribution. The role and the importance of such objects has been extensively studied within the context of the Einstein theory of gravitation and we refer the reader to existing reviews like and the references contained therein for their various aspects. A singular hypersurface is a geometric representation in spacetime of the interface existing between two coexisting phases of matter such that the curvature tensor suffers a discontinuity of the Dirac type at the inter-phase. This results in a two-dimensional concentrated distribution of matter whose history in spacetime corresponds to this hypersurface, a distribution of matter which is usually referred to as a thin shell or surface layer. Singular hypersurfaces stand for the idealization of a vast range of physical situations such as for example the false-vacuum bubbles used in inflationary cosmology (see the introduction of and the references listed) as well as in extended inflation in Brans-Dicke models (see for a review). They also have been successfully used to describe impulsive gravitational waves coexisting with null shells of matter by recently providing a model of supernovae in which the burst of gravitational radiation and the burst of neutrinos are respectively modeled by a gravitational impulsive wave and a null shell of matter such that the wave-fronts of the wave and the shell share a singular null hypersurface as their common history. More recently, the formalism of singular hypersurfaces derived in has been successfully applied to the description of the peeling properties of a light-like signal propagating through a general Bondi-Sachs model .
The Einstein-Cartan theory of gravitation was originally suggested by Elie Cartan in 1922 as a natural generalisation of Einstein’s theory of general relativity by extending the model of spacetime into a four-dimensional manifold endowed with a linear connection still compatible with the metric tensor but non symmetric. As a result, the torsion tensor, defined by Cartan as the antisymmetric part of the connection, plays a non trivial role in the geometry and Cartan’s idea was to relate this torsion to the density of intrinsic angular momentum . Cartan’s ideas reemerged in the 60’s with the works of Sciama and Kibble who independently developed a gauge theory approach to gravitation from which the modern $`U_4`$ theory emerged. For this reason, the Einstein-Cartan theory is also known as the Einstein-Cartan-Sciama-Kibble theory. Later, it was proved that this theory of gravitation could lead to non singular cosmological models . The mathematical structure of the Einstein-Cartan equations was also analyzed into a new perspective by A. Trautman , using tensor-valued differential forms, who proved that the Cartan equations determine the linear connection only up to projective transformations and that the condition for the connection to be metric-compatible is exactly the condition which removed this arbitrarness. For a more detailed historical development, we refer to the section I.6 of .
It is therefore of great interest to investigate the problem of discontinuities within this theory by analysing the matching conditions the various fields have to satisfy across a given singular hypersurface whose definition will be exposed within this paper. Gravitational matching conditions have first been considered in the Einstein-Cartan theory by Arkuszewski, Kopczyński and Pomomariev who derived the matching and the junction conditions for surfaces characterized by the continuity of the second fundamental form, also called extrinsic curvature. This corresponds to boundary surfaces in general relativity where the first derivatives of the metric tensor are continuous across the hypersurface, such as for example the surface separating a star from the surrounding vacuum. But this reference didn’t consider the case of more irregular discontinuities arising when the metric tensor is assumed to be only continuous across the surface with a jump in the extrinsic curvature which corresponds to thin shells. Moreover, they considered only timelike hypersurfaces. This paper is an attempt to describe these more general discontinuities by providing a formalism to derive the physical content of thin shells of any type, in particular null. This will be achieved by extending the Barrabès-Israel formalism of general relativity to the Einstein-Cartan theory in order to include the effects of spin.
The paper is organized as follows. In section 2, we present the Einstein-Cartan theory and summarize all the useful equations and identities we will need in this paper. We then discuss the type of spacetimes we will restrict to in order to derive our shell formalism. The definition of a singular hypersurface and the shell formalism are presented in section 3. The stress-energy tensor of the shell is derived both in its four dimensional and intrinsic versions. Finally, in section 4, we construct a general family of static Einstein-Cartan solutions with non-trivial torsion which contains a natural family of null hypersurfaces and we use it to construct a null shell of matter. The surface energy density of the shell is characterized by the jump in the local mass function of the two surrounding spacetimes and the surface pressure is given in terms of the jump in the components of the torsion across the hypersurface.
## 2 The Einstein-Cartan theory
The Einstein-Cartan theory considers a four-dimensional manifold $``$ as a model of spacetime, endowed with a metric tensor $`𝐠`$ of signature $`(,+,+,+)`$ satisfying the same regularity conditions as in General Relativity (three times continuously differentiable) but characterized by a non symmetric connection. The components of this connection are denoted $`\mathrm{\Gamma }_{\mu \nu }^\lambda `$ in a local coordinate system $`\{x^\mu \}`$, Greek indices taking the values 0,1, 2, 3. This implies the existence of a non-vanishing torsion tensor $`𝐐`$ such that
$$Q_{\mu \nu }^\lambda =\mathrm{\Gamma }_{\nu \mu }^\lambda \mathrm{\Gamma }_{\mu \nu }^\lambda .$$
(1)
The metric tensor $`g_{\mu \nu }`$ is still assumed to be metric-compatible with the covariant derivative associated with the above connection i.e. $`_\mu g_{\nu \lambda }=\mathrm{\hspace{0.17em}0}`$. It is useful to display the existing relation between the Einstein-Cartan metric compatible linear connection and the Riemannian connection constructed out of the metric $`g_{\mu \nu }`$ represented by its Christoffel symbols $`{}_{}{}^{R}\mathrm{\Gamma }_{\mu \nu }^{\lambda }`$, namely
$$\mathrm{\Gamma }_{\mu \nu }^\lambda ={}_{}{}^{R}\mathrm{\Gamma }_{\mu \nu }^{\lambda }+\chi _{\mu \nu }^\lambda ,$$
(2)
with
$${}_{}{}^{R}\mathrm{\Gamma }_{\mu \nu }^{\lambda }=\frac{1}{2}g^{\lambda \rho }(_\mu g_{\nu \rho }+_\nu g_{\mu \rho }_\rho g_{\mu \nu }),$$
(3)
where we introduced
$$\chi _{\lambda \mu \nu }=\frac{1}{2}\left(Q_{\lambda \mu \nu }+Q_{\mu \nu \lambda }+Q_{\nu \mu \lambda }\right)=\chi _{\mu \lambda \nu },$$
(4)
which is often referred as the “contortion”or “defect” of the connection (see references and ).
In the local coordinate system $`\{x^\mu \}`$, the components of the curvature tensor have the familiar expression
$$R_{\nu \lambda \rho }^\mu =_\lambda \mathrm{\Gamma }_{\nu \rho }^\mu _\rho \mathrm{\Gamma }_{\nu \lambda }^\mu +\mathrm{\Gamma }_{\kappa \lambda }^\mu \mathrm{\Gamma }_{\nu \rho }^\kappa \mathrm{\Gamma }_{\kappa \rho }^\mu \mathrm{\Gamma }_{\nu \lambda }^\kappa .$$
(5)
The symmetries of the Riemann tensor which remain within this context are
$$R_{\mu \nu \lambda \rho }=R_{\mu \nu \rho \lambda },$$
(6)
and
$$R_{\mu \nu \lambda \rho }=R_{\nu \mu \lambda \rho }.$$
(7)
The symmetry (6) comes from the general definition of a curvature operator and the symmetry (7) follows from the fact that the connection is metric compatible.
The cyclic identity is replaced by
$$R_{[\nu \lambda \rho ]}^\mu =_{[\nu }Q_{\lambda \rho ]}^\mu +Q_{\kappa [\nu }^\mu Q_{\lambda \rho ]}^\kappa ,$$
(8)
which is also called Bianchi identity for the torsion.
The property of symmetry by interchanging pairs of indices in the Riemann tensor is no longer true in this non-Riemannian spacetime but is replaced by the following identity, obtained using four times the identity (8):
$$R_{\mu \nu \lambda \rho }=R_{\lambda \rho \mu \nu }+S_{\mu \nu \lambda \rho }^{(1)}+S_{\mu \nu \lambda \rho }^{(2)},$$
(9)
with
$$S_{\mu \nu \lambda \rho }^{(1)}=_\nu \chi _{\lambda \rho \mu }+_\lambda \chi _{\mu \nu \rho }+_\rho \chi _{\nu \mu \lambda }+_\mu \chi _{\rho \lambda \nu },$$
(10)
and
$$S_{\mu \nu \lambda \rho }^{(2)}=\psi _{\mu \nu \lambda \rho }+\psi _{\lambda \mu \rho \nu }+\psi _{\rho \mu \nu \lambda }+\psi _{\nu \lambda \mu \rho },$$
(11)
where we introduced
$$\psi _{\mu \nu \lambda \rho }={\scriptscriptstyle \frac{1}{2}}\left(Q_{\mu \kappa \nu }Q_{\lambda \rho }^\kappa +Q_{\mu \kappa \lambda }Q_{\rho \nu }^\kappa +Q_{\mu \kappa \rho }Q_{\nu \lambda }^\kappa \right)=\psi _{\mu \lambda \nu \rho }=\psi _{\mu \rho \lambda \nu }.$$
(12)
The Bianchi identity is now
$$_{[\lambda }R_{\nu \rho \sigma ]}^\mu +R_{\nu \kappa [\lambda }^\mu Q_{\rho \sigma ]}^\kappa =\mathrm{\hspace{0.17em}0}.$$
(13)
The basic field equations of the Einstein-Cartan theory are
$$G_{\mu \nu }=\mathrm{\hspace{0.17em}8}\pi T_{\mu \nu },$$
(14)
and
$$Q_{\nu \lambda }^\mu +\delta _\nu ^\mu Q_{\lambda \rho }^\rho \delta _\lambda ^\mu Q_{\nu \rho }^\rho =\mathrm{\hspace{0.17em}8}\pi S_{\nu \lambda }^\mu ,$$
(15)
where $`S_{\nu \lambda }^\mu `$ is the spin tensor representing the density of intrinsic angular momentum, related to the torsion tensor in a purely algebraic way. $`T_{\mu \nu }`$ represents the stress-energy tensor of the spacetime matter content. In the field equation (14), $`G_{\mu \nu }`$ stands for the Einstein tensor constructed out of the Ricci tensor $`R_{\mu \nu }=R_{\mu \lambda \nu }^\lambda `$ by means of the non symmetric connection $`\mathrm{\Gamma }_{\mu \nu }^\lambda `$. We should insist on the fact that this makes the Einstein tensor non symmetric and therefore, in general, the stress-energy tensor of the spacetime is a non symmetric quantity within this framework. When the torsion vanishes so does the spin and these field equations reduce to the Einstein field equations of general relativity.
Contraction of the identity (8) yields
$$R_{\nu \rho }=R_{\rho \nu }+_\mu Q_{\nu \rho }^\mu +_\nu Q_{\rho \mu }^\mu _\rho Q_{\nu \mu }^\mu +Q_{\kappa \mu }^\mu Q_{\rho \nu }^\kappa .$$
(16)
Regrouping the covariant derivatives and taking into account the relation between spin and torsion given in (15), this equation leads directly to the generalized evolution equation of the spin tensor
$$_\mu S_{\nu \rho }^\mu =\mathrm{\hspace{0.17em}2}T_{[\nu \rho ]}+2S_{\kappa \mu }^\mu Q_{\nu \rho }^\kappa ,$$
(17)
which is therefore not conserved in general. In particular, if the stress-energy tensor of the spacetime is chosen to be symmetric, this evolution equation is
$$_\mu S_{\nu \rho }^\mu =\mathrm{\hspace{0.17em}2}S_{\kappa \mu }^\mu Q_{\nu \rho }^\kappa .$$
(18)
However, one must keep in mind that the torsion does not propagate (cf. ). Indeed, from the algebraic field equation (15), it is clear there cannot exist torsion outside the spinning matter distribution itself (the vanishing of the one is equivalent to the vanishing of the other) and therefore the torsion cannot propagate through vacuum. Outside the spinning matter distribution, the influence of the spin on the evolution of spacetime can only be seen on the metric tensor but this effect is of higher order than the effect of mass on the spacetime metric tensor. In some exact cosmological models known in the Einstein-Cartan literature (see for instance ), the stress-energy tensor of the spacetime is chosen to be symmetric and the torsion tensor trace-free. The spin tensor is in this sense conserved.
The twice contracted Bianchi identity (13) (which in general relativity directly leads to the conservation of the Einstein tensor and consequently to the conservation of the stress-energy tensor) here results in
$$_\nu G_\rho ^\nu =R_\lambda ^\mu Q_{\rho \mu }^\lambda +\frac{1}{2}R_{\lambda \rho }^{\mu \nu }Q_{\mu \nu }^\lambda .$$
(19)
This proves, in view of the field equation (14), that the stress-energy tensor is not conserved in general. This was remarked by Cartan himself who arrived at the conclusion that, in order for the stress-energy momentum to be conserved, a geometrical constraint relating the torsion and the curvature has to be imposed, namely
$$𝒞_\rho :=R_\lambda ^\mu Q_{\rho \mu }^\lambda +\frac{1}{2}R_{\lambda \rho }^{\mu \nu }Q_{\mu \nu }^\lambda =\mathrm{\hspace{0.17em}0}.$$
(20)
In the following sections, we will restrict ourselves to the class of spacetimes with torsion and curvature satifying this constraint. The existence of this constraint will also reveal itself to be necessary in order to properly describe thin shells within this framework as we will now investigate.
## 3 The Shell formalism
In this section, we develop a technique to describe thin shells of matter in the Einstein-Cartan theory by extending the known formalism of general relativity. A thin shell corresponds to a distribution of matter ideally assumed to be concentrated on a two-dimensional surface. Its history in spacetime is therefore geometrically represented by a three-dimensional hypersurface which can be either timelike, spacelike or lightlike. Such a shell is generally constructed using the “cut and paste” approach of two spacetimes along a hypersurface on which the induced metrics of both spacetimes are required to agree. Within the framework of general relativity, the Barrabès-Israel formalism provides a recipe to obtain a unified description of shells of any type by extending the known formalism for timelike and spacelike shells. Its main feature is a description of the surface properties in terms of the jump of the transverse extrinsic curvature (generalizing the usual extrinsic curvature - see below for the definition) and thus directly obtained as functions of the shell intrinsic coordinates. This makes possible the four-dimensional coordinates of each side of the layer to be chosen freely and independently. In particular, no continuous coordinate system is required and we can choose coordinates adapted to the peculiar symmetries of each spacetime which greatly simplifies the practical calculations. In the following paragraphs, we will extend this formalism in the presence of a non symmetric connection.
Let us consider a hypersurface $`\mathrm{\Sigma }`$ which can be either timelike, spacelike or lightlike. We assume that $`\mathrm{\Sigma }`$ results from the isometric soldering of two hypersurfaces $`\mathrm{\Sigma }^+`$ and $`\mathrm{\Sigma }^{}`$ bounding respectively two spacetimes $`M^+`$ and $`M^{}`$ -see figure 1. We also assume that $`M^{}`$ is endowed with a metric tensor, a non symmetric connection and a torsion tensor whose components are given respectively by $`g_{\alpha \beta }^{}`$, $`{}_{}{}^{}\mathrm{\Gamma }_{\beta \gamma }^{\alpha }`$ and $`{}_{}{}^{}Q_{\beta \gamma }^{\alpha }`$ with respect to a local coordinate system $`\{x_{}^\alpha \}`$ and that $`M^+`$ is endowed with a metric tensor, a non symmetric connection and a torsion tensor whose components are given respectively by $`g_{\alpha \beta }^+`$, $`{}_{}{}^{+}\mathrm{\Gamma }_{\beta \gamma }^{\alpha }`$ and $`{}_{}{}^{+}Q_{\beta \gamma }^{\alpha }`$ with respect to a local coordinate system $`\{x_+^\alpha \}`$. $`\mathrm{\Sigma }`$ can thereby be considered as a hypersurface imbedded in the hybrid manifold $`=M^{}M^+`$. The metrics $`g_{\alpha \beta }^{}`$ and $`g_{\alpha \beta }^+`$ are at least $`C^3`$-continuously differentiable respectively in $`M^{}`$ and $`M^+`$. The connections are at least $`C^2`$-continuously differentiable respectively in $`M^{}`$ and $`M^+`$ (from relation (2), this is equivalent to requiring the torsion tensors to be at least $`C^2`$-continuously differentiable respectively in each domain). If $`\{\xi ^a\}`$ is a set of intrinsic coordinates on $`\mathrm{\Sigma }`$ (where the Latin indices take the values 1, 2, 3), the vectors $`e_{(a)}=\frac{}{\xi ^a}`$ form a basis of tangent vectors to $`\mathrm{\Sigma }`$ at each point of $`\mathrm{\Sigma }`$.
Figure 1 : Isometric soldering of two spacetimes
For any quantity $`F`$ regular in $`/\mathrm{\Sigma }`$ (which can be a function aswell as a tensor field) such that $`F`$ tends uniformly to finite left and right limits on $`\mathrm{\Sigma }`$, we will denote as usual by $`[F]`$ the difference of these limits at a point on $`\mathrm{\Sigma }`$. $`[F]`$ is by definition a function defined on $`\mathrm{\Sigma }`$ and continuous on $`\mathrm{\Sigma }`$ and is called the jump of $`F`$ across $`\mathrm{\Sigma }`$.
We will call a singular hypersurface in spacetime $``$ a hypersurface such that the following regularity conditions occur across $`\mathrm{\Sigma }`$ :
$$[g_{\mu \nu }]=\mathrm{\hspace{0.17em}0}$$
(21)
$$[e_{(a)}^\mu e_{(b)}^\nu e_{(c)}^\lambda Q_{\mu \nu \lambda }]=[Q_{abc}]=\mathrm{\hspace{0.17em}0}$$
(22)
$$[\mathrm{\Gamma }_{\nu \lambda }^\mu ]\mathrm{\hspace{0.17em}0}$$
(23)
Condition (21) expresses the continuity of the metric across $`\mathrm{\Sigma }`$ which we assumed by construction of $`\mathrm{\Sigma }`$ and this means that the induced metrics on both sides of the hypersurface are the same i.e.
$$g_{\alpha \beta }e_{(a)}^\alpha e_{(b)}^\beta |_\pm =g_{ab}=e_{(a)}.e_{(b)},$$
(24)
where $`g_{ab}`$ stands for the the three-dimensional metric intrinsic to $`\mathrm{\Sigma }`$. Condition (22) is the requirement that the purely tangential part (i.e. the projection on $`\mathrm{\Sigma }`$) of the torsion tensor has to be continuous across $`\mathrm{\Sigma }`$. As three-dimensional manifold embedded in an Einstein-Cartan four-dimensional spacetime, $`\mathrm{\Sigma }`$ is also endowed with a three-dimensional torsion it inherits from the non symmetric connections of one or other side. Condition (22) is therefore the requirement that the induced torsions from both sides coincide with the intrinsic torsion of $`\mathrm{\Sigma }`$. If we define $`x_{\mu \nu \lambda }=[\chi _{\mu \nu \lambda }]`$ the jump of the defect across the hypersurface and denote by $`x_{abc}=x_{\mu \nu \lambda }e_{(a)}^\mu e_{(b)}^\nu e_{(c)}^\lambda `$ its projection onto $`\mathrm{\Sigma }`$, we immediately see from the definition (4) that $`x_{abc}={\scriptscriptstyle \frac{1}{2}}[Q_{abc}+Q_{bca}+Q_{cba}]`$ and $`x_{abc}+x_{bca}=[Q_{bca}]`$ proving that the vanishing of $`[Q_{abc}]`$ is equivalent to the vanishing of $`x_{abc}`$. This gives another set of conditions equivalent to (21)-(23) by simply replacing (22) by $`[\chi _{abc}]=\mathrm{\hspace{0.17em}0}`$. We point out that this is not equivalent to saying that the jump across $`\mathrm{\Sigma }`$ of the tangential part $`S_{abc}`$ of the spin tensor vanishes. Indeed, the field equation (15) only tells us that $`[Q_{abc}+Q_{bca}+Q_{cab}]=[S_{abc}+S_{bca}+S_{cab}]`$ showing that a necessary (but not sufficient) condition for (22) is $`[S_{abc}+S_{bca}+S_{cab}]=\mathrm{\hspace{0.17em}0}`$ and not $`[S_{abc}]=\mathrm{\hspace{0.17em}0}`$.
The condition (23) means that we assume a jump in the connection across the hypersurface and it is clear from the relation (2) that this discontinuity results from the discontinuities of the first derivatives of the metric tensor and from the discontinuities of the torsion.
Such a hypersurface, as we will see below, is characterized by a jump in the (transverse) extrinsic curvature of the hypersurface and should be distinguish from the boundary surfaces studied in reference where the extrinsic curvature is assumed to be continuous (this reference, incidently, is confined to the timelike case).<sup>1</sup><sup>1</sup>1This continuity condition is proved to be equivalent to the regularity of the full Riemann tensor, a condition which is not required in the context of thin shells.
Following the procedure of reference , we now describe how to obtain the expression for the stress-energy tensor carried by the singular hypersurface $`\mathrm{\Sigma }`$ in its four-dimensional form and in its intrinsic form. We first consider a single coordinate system $`\{y^\mu \}`$ of $``$ maximally smooth across $`\mathrm{\Sigma }`$. Let $`n`$ be a normal vector field to $`\mathrm{\Sigma }`$ (uniquely defined up to a sign in the timelike or spacelike case but not in the null case) normalized by
$$n.n=ϵ,$$
(25)
$`ϵ`$ being positive, negative or zero respectively when the hypersurface is timelike, spacelike or null. It satisfies $`n.e_{(a)}=\mathrm{\hspace{0.17em}0}`$. Let $`\mathrm{\Phi }(y^\mu )=\mathrm{\hspace{0.17em}0}`$ be the local equation of the hypersurface in this coordinate system and assume that the normal is related to the gradient via
$$n_\mu =\alpha ^1_\mu \mathrm{\Phi },$$
(26)
where $`\alpha `$ is a non-vanishing function of the local coordinates. We introduce a transverse vector field -see figure 1- non uniquely defined by
$$N.n=\eta ^1,$$
(27)
where $`\eta `$ is a non-vanishing function defined on $`\mathrm{\Sigma }`$. For a fixed $`\eta `$, $`N`$ is free up to a tangential displacement
$$NN+\lambda ^a(\xi ^b)e_{(a)},$$
(28)
where $`\lambda ^a`$ is an arbitrary vector field on $`\mathrm{\Sigma }`$. According to the regularity condition (23), we will employ the following notation
$$w_{\nu \rho }^\mu =[\mathrm{\Gamma }_{\nu \rho }^\mu ].$$
(29)
Because the metric is supposed to be only continuous, the Riemannian part of the connection $`{}_{}{}^{R}\mathrm{\Gamma }_{\mu \nu }^{\lambda }`$ is also jumping across $`\mathrm{\Sigma }`$ and we get from (2)
$$w_{\nu \rho }^\mu ={}_{}{}^{R}w_{\nu \rho }^{\mu }+x_{\nu \rho }^\mu ,$$
(30)
where $`{}_{}{}^{R}w_{\nu \rho }^{\mu }`$ is the jump across $`\mathrm{\Sigma }`$ of the Riemannian connection $`{}_{}{}^{R}\mathrm{\Gamma }_{\nu \rho }^{\mu }`$. We now assume that the Einstein-Cartan equations (14)-(15) apply in the sense of distributions on the manifold $``$, following the now classical techniques first introduced by Lichnerowicz and Taub . From the regularity conditions (21)-(23) and the general expression of the curvature tensor (5), the curvature tensor contains a distributional Dirac $`\delta (\mathrm{\Phi })`$ part $`\widehat{R}_{\mu \nu \rho \lambda }`$ which is given by
$$\alpha ^1\widehat{R}_{\mu \nu \lambda \rho }=n_\lambda w_{\mu \nu \rho }n_\rho w_{\mu \nu \lambda }.$$
(31)
Requiring that the distributional curvature tensor must satisfy the identity (9) and extracting the impulsive part of this identity, we arrive with the help of (31) at
$$n_\lambda w_{\mu \nu \rho }n_\rho w_{\mu \nu \lambda }=n_\mu w_{\lambda \rho \nu }n_\nu w_{\lambda \rho \mu }+x_{\mu \nu \rho }n_\lambda +x_{\nu \mu \lambda }n_\rho +x_{\rho \lambda \nu }n_\mu +x_{\lambda \rho \mu }n_\nu .$$
(32)
Taking into account the relation (30), equation (32) turns out to be an identity involving only the Riemannian connection, namely
$$n_\lambda {}_{}{}^{R}w_{\mu \nu \rho }^{}n_\rho {}_{}{}^{R}w_{\mu \nu \lambda }^{}=n_\mu {}_{}{}^{R}w_{\lambda \rho \nu }^{}n_\nu {}_{}{}^{R}w_{\lambda \rho \mu }^{}.$$
(33)
Contracting this last identity by $`N^\rho `$ leads to the following algebraic form of $`{}_{}{}^{R}w_{\mu \nu \lambda }^{}`$
$$\eta ^1{}_{}{}^{R}w_{\mu \nu \lambda }^{}=\frac{1}{2}\left(p_{\mu \nu }^Rn_\lambda +p_{\lambda \mu }^Rn_\nu p_{\lambda \nu }^Rn_\mu \right),$$
(34)
where we have introduced the quantity
$$p_{\mu \nu }^R=\mathrm{\hspace{0.17em}2}{}_{}{}^{R}w_{\mu \nu \rho }^{}N^\rho .$$
(35)
The symmetry of the Riemannian connection applied to the identity (34) now leads directly to the expression for the jump of the first derivatives of the metric
$$[_\nu g_{\lambda \mu }]=\eta \gamma _{\lambda \mu }^Rn_\nu ,$$
(36)
with
$$\gamma _{\lambda \mu }^R=N^\nu [_\nu g_{\lambda \mu }],$$
(37)
which is exactly the same result obtained in in general relativity, with the same notation. From relation (30), we now finally arrive at the jump of the Einstein-Cartan connection
$$w_{\mu \nu \lambda }=\frac{\eta }{2}\left(\gamma _{\mu \nu }^Rn_\lambda +\gamma _{\lambda \mu }^Rn_\nu \gamma _{\lambda \nu }^Rn_\mu \right)+x_{\mu \nu \lambda }.$$
(38)
The impulsive part of the Riemann tensor (31) is therefore
$$\alpha ^1\widehat{R}_{\mu \nu \lambda \rho }=n_\lambda x_{\mu \nu \rho }n_\rho x_{\mu \nu \lambda }+\frac{\eta }{2}\left(\gamma _{\mu \rho }^Rn_\nu n_\lambda \gamma _{\nu \rho }^Rn_\mu n_\lambda +\gamma _{\nu \lambda }^Rn_\mu n_\rho \gamma _{\mu \lambda }^Rn_\nu n_\rho \right).$$
(39)
By successive contractions, we deduce from (39) the expression for the impulsive part of the Einstein-tensor
$$\alpha ^1\widehat{G}_{\nu \rho }=x_{\mu \nu \rho }n^\mu x_{\nu \mu }^\mu n_\rho +x_{\mu \lambda }^\lambda n^\mu g_{\nu \rho }+\alpha ^1\widehat{G}_{\nu \rho }^R,$$
(40)
which clearly splits into two parts, a pure Einstein-Cartan part involving the jump in the defect tensor $`x_{\mu \nu \rho }`$, i.e. the jump in the torsion, and an Einstein-Riemann part $`\widehat{G}_{\nu \rho }^R`$ corresponding to the impulsive part of the Einstein tensor obtained in general relativity and given by
$$\alpha ^1\widehat{G}_{\nu \rho }^R=\frac{\eta }{2}\left(\gamma _\nu ^Rn_\rho +\gamma _\rho ^Rn_\nu \gamma ^Rn_\nu n_\rho \gamma _R^{}g_{\nu \rho }ϵ(\gamma _{\nu \rho }^R\gamma ^Rg_{\nu \rho })\right),$$
(41)
where we have set, following ,
$$\gamma _\nu ^R=\gamma _{\nu \mu }^Rn^\mu ,\gamma ^R=g^{\mu \nu }\gamma _{\mu \nu }^R,\gamma _R^{}=\gamma _{\mu \nu }^Rn^\mu n^\nu .$$
(42)
From the field equation (14), the distributional stress-energy tensor $`T^{\nu \rho }`$ of the hybrid space-time $``$ exhibits a Dirac part $`\alpha \mathrm{\Sigma }_{\nu \rho }`$ given by
$$8\pi \alpha \mathrm{\Sigma }^{\nu \rho }=\widehat{G}^{\nu \rho },$$
(43)
which we can express as
$$\mathrm{\Sigma }^{\nu \rho }=\mathrm{\Sigma }_R^{\nu \rho }+{\scriptscriptstyle \frac{1}{8\pi }}\left(x^{\mu \nu \rho }n_\mu x_\mu ^{\mu \nu }n^\rho +x_{\mu \lambda }^\lambda n^\mu g^{\nu \rho }\right),$$
(44)
where $`\mathrm{\Sigma }_R^{\nu \rho }`$ is the stress-energy tensor of the shell obtained in general relativity in given by
$$16\pi \eta ^1\mathrm{\Sigma }_R^{\nu \rho }=\gamma _R^\nu n^\rho +\gamma _R^\rho n^\nu \gamma _Rn^\nu n^\rho \gamma _R^{}g^{\nu \rho }ϵ(\gamma _R^{\nu \rho }\gamma _Rg^{\nu \rho }).$$
(45)
Using the relations (4), (15) and (29), we can obtain the following useful form of $`\mathrm{\Sigma }^{\nu \rho }`$ in terms of the spin tensor
$$\mathrm{\Sigma }^{\nu \rho }=\mathrm{\Sigma }_R^{\nu \rho }{\scriptscriptstyle \frac{1}{2}}[\left(S^{\mu \nu \rho }+S^{\nu \rho \mu }+S^{\rho \nu \mu }\right)n_\mu ].$$
(46)
In order to interpret the quantity (44) as the stress-energy tensor of the matter shell having the hypersurface $`\mathrm{\Sigma }`$ as history in spacetime, our regularity conditions (21)-(23) must ensure that this tensor is indeed purely tangential to $`\mathrm{\Sigma }`$. As $`\mathrm{\Sigma }^{\nu \rho }`$ is not symmetric in general, this leads to the two conditions
$$(\mathrm{a})\mathrm{\Sigma }^{\nu \rho }n_\nu =0;(\mathrm{b})\mathrm{\Sigma }^{\nu \rho }n_\rho =0.$$
(47)
It is easy to verify that condition (a) is automatically satisfied by the algebraic form (44) and this can been considered with our assumptions ($`𝒞_\rho =\mathrm{\hspace{0.17em}0}`$) as an algebraic consequence of the twice contracted Bianchi identity $`_\nu G_\mu ^\nu =\mathrm{\hspace{0.17em}0}`$ by extracting its $`\delta ^{}`$-distributional part. Condition (b) is equivalent to the following condition
$$x_{\mu \nu \rho }n^\mu n^\rho +(x_{\rho \mu }^\mu n^\rho )n_\nu ϵx_{\nu \mu }^\mu =\mathrm{\hspace{0.17em}0},$$
(48)
and from (46), we get that this condition is equivalent to the following requirement on the spin tensor
$$[S_{\nu \rho }^\mu n_\mu n^\nu ]=\mathrm{\hspace{0.17em}0}.$$
(49)
We shall analyze in detail the intrinsic content of this constraint and prove that it is satisfied.
We first consider the case of a timelike or spacelike shell ($`ϵ\mathrm{\hspace{0.17em}0}`$). In this situation, the transversal $`N`$ can be chosen to be the normal $`n`$ and the projector operator on $`\mathrm{\Sigma }`$ is defined by
$$g^{ab}e_{(a)}^\mu e_{(b)}^\nu =g^{\mu \nu }ϵ^1n^\mu n^\nu ,$$
(50)
where $`g^{ab}`$ represents the inverse induced metric which exists in this case. The constraint (48) is now reduced to the intrinsic trace-free condition
$$g^{bc}x_{bac}=x_{ab}^b=\mathrm{\hspace{0.17em}0}.$$
(51)
From the definition (4) and field equation (15), this three-dimensional constraint turns to be
$$[S_{ab}^b]=\mathrm{\hspace{0.17em}0},$$
(52)
which is only a part of the constraint obtained in but is weaker because we do not assume continuity of the extrinsic curvature.
For a lightlike hypersurface ($`ϵ=\mathrm{\hspace{0.17em}0}`$), the induced metric is degenerate and an inverse does not exist. Following the technique introduced in , we introduce a pseudo-inverse symmetric tensor $`g_{}^{ab}`$ not uniquely defined by the following relation
$$g_{}^{ac}g_{cb}=\delta _b^a\eta l^aN_b,$$
(53)
where $`N_b=N.e_{(b)}`$ and where the components $`l^a`$ are such that in the oblique basis $`(N,e_{(a)})`$, the normal vector $`n`$ expands as
$$n=ϵ\eta N+l^ae_{(a)}.$$
(54)
The relation (50) is now generalized to
$$g^{\mu \nu }=g_{}^{ab}e_{(a)}^\mu e_{(b)}^\nu +2\eta l^ae_{(a)}^{(\mu }N^{\nu )}+\eta ^2ϵN^\mu N^\nu .$$
(55)
It can easily be proved that the constraint (48) reduces now to the two conditions
$`x_{abc}l^bl^c=\mathrm{\hspace{0.17em}0},`$ (56)
$`l^ag_{}^{bc}x_{abc}=\mathrm{\hspace{0.17em}0}.`$ (57)
As we pointed out, the regularity condition (22) is equivalent to $`x_{abc}=\mathrm{\hspace{0.17em}0}`$. The constraints (51) and (56) respectively in the timelike/spacelike and lightlike cases are thereby automatically satisfied proving that the assumption (22) is sufficient. We shall prove now that it is necessary.
In order to describe the intrinsic matter content of the shell represented by its three-dimensional stress-energy tensor, we introduce, following reference , the “transverse” extrinsic curvature $`𝒦_{ab}`$ of the hypersurface $`\mathrm{\Sigma }`$ generalizing the usual extrinsic curvature (whose utility was confined to the timelike and spacelike shells) via
$$𝒦_{ab}=N._{e_{(b)}}e_{(a)}.$$
(58)
In the Einstein-Cartan theory, this quantity is non symmetric. As in general relativity, $`𝒦_{ab}`$ is not a three-tensor under changes of intrinsic coordinates $`\xi ^a`$. Moreover, it is not invariant under the transformation on $`N`$ (28) but transforms according to
$$𝒦_{ab}𝒦_{ab}\lambda ^c\mathrm{\Gamma }_{c,ab}=𝒦_{ab}\lambda ^c{}_{}{}^{R}\mathrm{\Gamma }_{c,ab}^{}\lambda ^c\chi _{cab},$$
(59)
where $`\mathrm{\Gamma }_{c,ab}`$ represents the three-dimensional connection induced on the submanifold $`\mathrm{\Sigma }`$ by the connection $`\mathrm{\Gamma }_{\nu \lambda }^\mu `$ on $``$. Let us now consider the jump of $`𝒦_{ab}`$ across $`\mathrm{\Sigma }`$ by considering as in reference the following quantity
$$\gamma _{ab}=\mathrm{\hspace{0.17em}2}[𝒦_{ab}].$$
(60)
In order to make the formalism independent of the choice of the transversal $`N`$, we must require $`\gamma _{ab}`$ to be gauge invariant under (28). In general relativity, the Riemann-Christoffel symbols $`{}_{}{}^{R}\mathrm{\Gamma }_{c,ab}^{}`$ of the three-dimensional manifold $`\mathrm{\Sigma }`$ are constructed only out of tangential derivatives of the metric which must be continuous across $`\mathrm{\Sigma }`$. The jump of $`\gamma _{ab}`$ across $`\mathrm{\Sigma }`$ is therefore zero. This gauge invariance is thereby a direct consequence of the regularity condition (21). In the Einstein-Cartan theory, the jump of the Riemann-Christoffel symbols associated to the metric is also zero but the transformation (28) gives an additional term in (59) which involves the defect. It is now clear that the gauge invariance of $`\gamma _{ab}`$ requires the jump $`\lambda ^cx_{cab}`$ to be zero for an arbitrary vector field $`\lambda ^c(\xi ^a)`$ on $`\mathrm{\Sigma }`$, and thus $`x_{abc}=\mathrm{\hspace{0.17em}0}`$. This leads directly to the condition (22). This condition also ensures that, under change of intrinsic coordinates $`\xi ^b\xi ^{}_{}{}^{}a(\xi ^b)`$, the quantity $`\gamma _{ab}`$ behaves as a three tensor intrinsic to $`\mathrm{\Sigma }`$.
The four-dimensional tensor $`\mathrm{\Sigma }^{\mu \nu }`$ (44) of the shell being purely tangential to $`\mathrm{\Sigma }`$, there exists a three-dimensional tensor $`\mathrm{\Sigma }^{ab}`$ such that
$$\mathrm{\Sigma }^{\mu \nu }=\mathrm{\Sigma }^{ab}e_{(a)}^\mu e_{(b)}^\nu ,$$
(61)
which is not simply obtained from $`\mathrm{\Sigma }_{ab}=e_{(a)}^\mu e_{(b)}^\nu \mathrm{\Sigma }_{\mu \nu }`$ by raising the three-indices since in general, the induced metric could be degenerate. After some algebra, we obtain the following intrinsic stress-energy tensor of the shell
$$16\pi \eta ^1\mathrm{\Sigma }^{ab}=\left(g_{}^{ac}l^bl^d+l^al^cg_{}^{bd}g_{}^{ab}l^cl^dl^al^bg_{}^{cd}\right)\gamma _{cd}ϵ\left(g_{}^{ac}g_{}^{bd}g_{}^{ab}g_{}^{cd}\right)\gamma _{cd},$$
(62)
which is formally the same expression as in general relativity but $`\gamma _{ab}`$ is now a non symmetric quantity which splits into a Riemann part and a Cartan part, namely
$$\gamma _{ab}={}_{}{}^{R}\gamma _{ab}^{}+\beta _{ab},$$
(63)
where $`{}_{}{}^{R}\gamma _{ab}^{}`$ is the Riemannian quantity used in , i.e. the jump in the Riemannian transverse extrinsic curvature defined by the formula (58) but calculated with the Riemannian connection. In (63), we introduced the projection $`\beta _{ab}`$ of the tensor $`\beta _{\mu \nu }`$ defined by
$$\beta _{\mu \nu }=2N^\rho x_{\rho \mu \nu }.$$
(64)
The quantity (62) is, as in general relativity, both independent of the choice of the transversal and of the choice of the pseudo-tensor $`g_{}^{ab}`$ and this independence is mainly ensured by the continuity requirement $`x_{abc}=\mathrm{\hspace{0.17em}0}`$. Let us insist on the fact that, as in general relativity, despite the fact that the four dimensional expression (44) has been obtained using distributional theory applied in a smooth coordinate system, the construction of the intrinsic expression (62) does not require the construction of such spacetime coordinates that match continuously at $`\mathrm{\Sigma }`$ and in which the four dimensional metric is continuous. The expression (62) can be calculated on either side of the shell using its adapted coordinates.
Let us now summarize the procedure to construct the surface stress-energy tensor of the shell :
* In the hybrid spacetime $`=M^+M^{}`$ where $`\mathrm{\Sigma }`$ is embedded, choose a transversal $`N`$ and a normal $`n`$ such that (25) and (27) are satisfied and such that
$$[N_a]=[N.e_{(a)}]=\mathrm{\hspace{0.17em}0},$$
(65)
in order for $`N`$ to be geometrically well defined in $``$ as it can be given in $`M^+`$ and $`M^{}`$ by two different sets of coordinates $`\{N_\alpha ^+\}`$ and $`\{N_\alpha ^{}\}`$.
* Choose one of the two coordinates systems $`\{x_{}^\alpha \}`$ or $`\{x_+^\alpha \}`$ (we then drop the positive or negative suffix) and extend the quantity $`\gamma _{ab}`$ obtained from (63) to any four tensor $`\gamma _{\mu \nu }`$ such that
$$\gamma _{\mu \nu }e_{(a)}^\mu e_{(b)}^\nu =\gamma _{ab}.$$
(66)
Here $`\gamma _{\mu \nu }`$ has the form $`\gamma _{\mu \nu }={}_{}{}^{R}\gamma _{\mu \nu }^{}+\beta _{\mu \nu }`$ where $`{}_{}{}^{R}\gamma _{\mu \nu }^{}`$ is any four tensor such that $`{}_{}{}^{R}\gamma _{\mu \nu }^{}e_{(a)}^\mu e_{(b)}^\nu ={}_{}{}^{R}\gamma _{ab}^{}`$. Thus we obtain the four dimensional form (44). We can equivalently obtain directly the intrinsic form (62) by means of the pseudo-inverse $`g_{}^{ab}`$.
In the case of a lightlike shell, the intrinsic stress-energy tensor (62) is reduced to
$$16\pi \eta ^1\mathrm{\Sigma }^{ab}=g_{}^{ac}l^b\left(\gamma _{cd}l^d\right)+l^ag_{}^{bd}\left(l^c\gamma _{cd}\right)g_{}^{ab}\left(\gamma _{cd}l^cl^d\right)l^al^b\left(g_{}^{cd}\gamma _{cd}\right),$$
(67)
showing, as in general relativity -see reference , that there is a part $`\widehat{\gamma }_{ab}`$ of $`\gamma _{ab}`$ which does not contribute to $`\mathrm{\Sigma }^{ab}`$. $`\widehat{\gamma }_{ab}`$ satisfies the following system of 7 equations
$$\{\begin{array}{ccc}\widehat{\gamma }_{cd}l^d& =& 0\\ \widehat{\gamma }_{cd}l^c& =& 0\\ g_{}^{cd}\widehat{\gamma }_{cd}& =& 0\end{array},$$
(68)
and is the generalisation of the quantity obtained in reference (where $`\gamma _{ab}`$ was symmetric and therefore $`\widehat{\gamma }_{ab}`$ satisfied only 4 equations). From the general splitting (63), it is clear that $`\widehat{\gamma }_{ab}`$ is obtained as
$$\widehat{\gamma }_{ab}={}_{}{}^{R}\widehat{\gamma }_{ab}^{}+\widehat{\beta }_{ab},$$
(69)
where $`{}_{}{}^{R}\widehat{\gamma }_{ab}^{}`$ is the symmetric quantity obtained in and where $`\widehat{\beta }_{ab}`$ is the part of $`\beta _{ab}`$ which satisfies the system (68). If we choose the transversal to be a null vector (which is always possible), one obtains
$$\widehat{\gamma }_{ab}=\gamma _{ab}{\scriptscriptstyle \frac{1}{2}}\left(g_{}^{cd}\gamma _{cd}\right)g_{ab}\eta \left(l^c\gamma _{cb}\right)N_a\eta \left(\gamma _{ac}l^c\right)N_b+\eta ^2\left(\gamma _{cd}l^cl^d\right)N_aN_b,$$
(70)
with
$$\widehat{\beta }_{ab}=\beta _{ab}{\scriptscriptstyle \frac{1}{2}}\left(g_{}^{cd}\beta _{cd}\right)g_{ab}\eta \left(l^c\beta _{cb}\right)N_a\eta \left(\beta _{ac}l^c\right)N_b+\eta ^2\left(\beta _{cd}l^cl^d\right)N_aN_b,$$
(71)
and with
$${}_{}{}^{R}\widehat{\gamma }_{ab}^{}={}_{}{}^{R}\gamma _{ab}^{}{\scriptscriptstyle \frac{1}{2}}\left(g_{}^{cd}{}_{}{}^{R}\gamma _{cd}^{}\right)g_{ab}2\eta l^c{}_{}{}^{R}\gamma _{c(a}^{}N_{b)}+\eta ^2\left({}_{}{}^{R}\gamma _{cd}^{}l^cl^d\right)N_aN_b.$$
(72)
By analogy with general relativity, we interpret $`\widehat{\gamma }_{ab}`$ (the part of $`\gamma _{ab}`$ which does not contribute to the expression of the intrinsic stress-energy tensor $`\mathrm{\Sigma }^{ab}`$ of the null shell) as being an impulsive gravitational wave propagating in $``$ whose history in $``$ is the null hypersurface $`\mathrm{\Sigma }`$ and which propagates independently of the null shell. Since $`\gamma _{ab}`$ has 9 independent components and as $`\widehat{\gamma }_{ab}`$ satisfies the system (68) of 7 independent equations, it follows that $`\widehat{\gamma }_{ab}`$ has two independent components which can be interpreted as representing the two degrees of polarization of the gravitational impulsive wave.
In the following section, we will apply this formalism to the construction of a null shell.
## 4 Application: an example of a null shell
We first construct a solution of the Einstein-Cartan equations which has an easily identifiable family of null hypersurfaces. To achieve this we look for a metric tensor of the general static form
$$ds^2=du\left(f(r)du+2dr\right)+r^2\left(d\theta ^2+\mathrm{sin}^2\theta d\varphi ^2\right),$$
(73)
and a torsion tensor with all components vanishing except
$$Q_{ur}^u=\alpha (r).$$
(74)
In order to apply the shell formalism of the preceeding section, we first impose the space-time constraint (20) which corresponds to the vanishing of the components of the vector field $`𝒞_\rho `$. These components vanish identically except for
$$𝒞_u=\frac{\alpha \left(f^{}+2\alpha f\right)}{r}.$$
(75)
This constraint leads therefore immediately to
$$\alpha (r)=\frac{1}{2}\frac{f^{}}{f},$$
(76)
where the stands for $`d/dr`$.
The evolution equation for the torsion tensor is
$$_\mu Q_{\nu \rho }^\mu =\mathrm{\hspace{0.17em}2}G_{[\nu \rho ]}Q_{\kappa \mu }^\mu Q_{\nu \rho }^\kappa .$$
(77)
But with the metric (73) and the torsion (74), the Einstein tensor turns out to be symmetric and with (76) is given by
$$G_{\mu \nu }=\left[\begin{array}{cccc}\frac{f\left(f^{}r1+f\right)}{r^2}& \frac{f^{}r1+f}{r^2}& 0& 0\\ \multicolumn{4}{c}{}\\ \frac{f^{}r1+f}{r^2}& \frac{f^{}}{rf}& 0& 0\\ \multicolumn{4}{c}{}\\ 0& 0& \frac{1}{2}rf^{}& 0\\ \multicolumn{4}{c}{}\\ 0& 0& 0& \frac{1}{2}rf^{}\mathrm{sin}^2\theta \end{array}\right].$$
(78)
It can be checked that our assumptions lead to
$$Q_{\kappa \mu }^\mu Q_{\nu \rho }^\kappa =\mathrm{\hspace{0.17em}0}.$$
(79)
As a result, the evolution equation for the torsion tensor reduces to
$$_\mu Q_{\nu \rho }^\mu =\mathrm{\hspace{0.17em}0},$$
(80)
an equation which is now identically satisfied. It can also be checked after tedious algebra that the two Bianchi identities (8) and (13) are both satisfied by virtue of the torsion choice (76). The stress-energy tensor of the spacetime is therefore given by $`T_{\mu \nu }=\mathrm{\hspace{0.17em}1}/8\pi G_{\mu \nu }`$. This $`T_{\mu \nu }`$ is a conserved symmetric tensor whose proper values are all real <sup>2</sup><sup>2</sup>2These proper values $`\lambda `$, obtained as the roots of the equation $`\mathrm{det}(T_{\alpha \beta }\lambda g_{\alpha \beta })=\mathrm{\hspace{0.17em}0}`$, are $`\lambda _0=\frac{rf^{}+f1}{8\pi r^2}`$, $`\lambda _1=\frac{f1}{8\pi r^2}`$ and $`\lambda _2=\frac{f^{}}{16\pi r}`$ (double root)..
From the second Cartan field equation (15), we obtain the only non zero components of the spin tensor:
$$S_{r\theta }^\theta =S_{r\varphi }^\varphi =\frac{\alpha }{8\pi }=\frac{1}{16\pi }\frac{f^{}}{f}.$$
(81)
In the spacetime with metric (73), the $`u=const`$ form a family of null hypersurfaces. Using the technique described in the preceeding section, we would like to glue two spacetimes $`M_+`$ and $`M_{}`$ both endowed with a metric of the general form (73) in two different coordinate systems $`(u,r_+,\theta _+,\varphi _+)`$ and $`(u,r,\theta ,\varphi )`$ with torsion given by (74) and (76) but with different functions $`f_+(r_+)`$ and $`f_{}(r)`$ and different $`\alpha _+(r_+)`$ and $`\alpha _{}(r)`$. We join these spacetimes along the null hypersurface $`\mathrm{\Sigma }`$ with equation $`u=\mathrm{\hspace{0.17em}0}`$ by requiring the positive and negative sides of $`\mathrm{\Sigma }`$ to be isometrically soldered via the identity matching
$$r_+=r,\theta _+=\theta ,\varphi _+=\varphi .$$
(82)
This enables the metric continuity condition (21) to be automatically satisfied, the common induced metric being $`ds_{|_\mathrm{\Sigma }}^2=r^2d\mathrm{\Omega }^2`$ and $`r`$ is chosen as a common parameter along the generators of each side of the hypersurface. Let us choose $`\xi ^a=(r,\theta .\varphi )`$ as the intrinsic coordinates on the three-dimensional manifold $`\mathrm{\Sigma }`$. The continuity condition (22) for the torsion tensor is then automatically satisfied since on both sides of $`\mathrm{\Sigma }`$, the torsion tensor satisfies $`Q_{abc}=\mathrm{\hspace{0.17em}0}`$ as well as the necessary condition $`[S_{abc}+S_{bca}+S_{cab}]=\mathrm{\hspace{0.17em}0}`$.
The future-directed normal generator of $`\mathrm{\Sigma }`$ is the null vector
$$n=\frac{}{r},$$
(83)
and as transversal we choose the null vector
$$N=\frac{}{u}{\scriptscriptstyle \frac{1}{2}}f\frac{}{r}.$$
(84)
On both sides of $`\mathrm{\Sigma }`$, we have
$`N.N`$ $`=`$ $`0,`$ (85)
$`N.n`$ $`=`$ $`\eta =1,`$ (86)
$`N.e_{(A)}`$ $`=`$ $`0,(A=2,3),`$ (87)
which makes $`N`$ a geometrically well-defined object.
The non zero components of the Riemannian transverse extrinsic curvature are
$`K_{\theta \theta }^R|_\pm `$ $`=`$ $`{\scriptscriptstyle \frac{1}{2}}rf^\pm ,`$ (88)
$`K_{\varphi \varphi }^R|_\pm `$ $`=`$ $`{\scriptscriptstyle \frac{1}{2}}r\mathrm{sin}^2\theta f^\pm ,`$ (89)
from which we get that the Riemannian part $`\gamma _{ab}^R`$ of $`\gamma _{ab}`$ is
$$\gamma _{ab}^R=\left[\begin{array}{ccc}0& 0& 0\\ 0& r\left[f\right]& 0\\ 0& 0& r\mathrm{sin}^2\theta \left[f\right]\end{array}\right].$$
(90)
Since $`l^a=\delta _r^a`$, $`N_b=N.e_{(b)}=\delta _b^r`$ and $`\eta =1`$, we can choose $`g_{}^{ab}`$ as being $`g^{AB}`$ ($`A,B=2,3`$) bordered by zeros and take
$$g_{}^{ab}=\left[\begin{array}{ccc}0& 0& 0\\ 0& r^2& 0\\ 0& 0& r^2\mathrm{sin}^2\theta \end{array}\right],$$
(91)
and according to (62) and (63), the Riemann part of the stress-energy tensor of the shell is given by
$$16\pi \mathrm{\Sigma }_R^{ab}=2\frac{\left[f\right]}{r}l^al^b=\mathrm{\hspace{0.17em}4}\pi \frac{\left[m\right]}{r^2}l^al^b,$$
(92)
where we introduced as in a local mass function $`m(r)`$ given by $`f(r)=12m/r`$.
In order to determine the Cartan part of the surface stress-energy tensor, we now calculate the tensor $`\beta _{ab}`$ defined in (64). With $`N`$ given by (84), it turns to be $`\beta _{ab}=2x_{uab}`$. It is not difficult to see that $`x_{urr}`$ is the only non zero component of $`x_{uab}`$ and we get
$$\beta _{ab}=\mathrm{\hspace{0.17em}2}\left[\alpha \right]\delta _a^r\delta _b^r.$$
(93)
The Cartan part of the surface stress-energy tensor is finally obtained as
$$8\pi \mathrm{\Sigma }_C^{ab}=\left[\alpha \right]g_{}^{ab}.$$
(94)
In conclusion, we have found that the hypersurface $`\mathrm{\Sigma }`$ is the history of a null shell of matter whose surface stress-energy tensor is given by
$$8\pi \mathrm{\Sigma }^{ab}=\mathrm{\hspace{0.17em}2}\frac{\left[m\right]}{r^2}l^al^b+\left[\alpha \right]g_{}^{ab}.$$
(95)
This means that the matter on the shell is characterized by a surface energy density $`\sigma `$ expressed as
$$4\pi r^2\sigma =\left[m\right],$$
(96)
which owes its existence to the jump in the mass function $`m`$, and by a pressure $`P`$ given by
$$8\pi P=\left[\alpha \right],$$
(97)
whose existence is due to the jump in the torsion function. Finally, it can be easily seen that the wave part $`\widehat{\gamma }_{ab}`$ defined in (70) vanishes identically showing that there is no gravitational impulsive wave associated with this null shell of matter.
## Acknowledgements
I would like to express my gratitude to Professor P. A. Hogan for many valuable and stimulating discussions, comments and for constant advice and encouragement. The author wishes also to thank the Departement of Education and Science (HEA) for a post-doctoral fellowship.
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# 1 Introduction
## 1 Introduction
Recently Born–Infeld gauge theory has attracted considerable interest as the bosonic light–brane approximation or limit of superstring theory, and has turned out to be a simple and transparent model in this context . Branes, defined as extended objects in spacetime, can be fundamental or solitonic. The connection of these branes with a $`U(1)`$ gauge field was motivated by the presence of this field in the massless part of the spectrum of open strings, and by realising that branes with open strings attached to them which satisfy Dirichlet boundary conditions, or more generally one brane attached to another, can become classically stable, solitonic objects. It is for this reason that the dynamics of $`D`$-branes in Born–Infeld theory is being studied in detail and generalised ,,,,. Since, in general, a brane may or may not be a solitonic configuration or BPS state, the exploration of this question deserves particular attention. It is often stated that a brane is BPS in view of the vanishing of a fraction of the supersymmetry variation of the associated gaugino field. However, since BPS states (as classically and topologically stable states) and Bogomol’nyi bounds have been studied in great detail in a host of other theories, and the approach in these is practically standard, one would like to understand aspects of Born–Infeld particles in a similar way, also because it is not absolutely clear that $`D`$–branes are solitons of string theory in precisely the same way as more familiar topological solitons in field theory. Therefore our first intention in the following is to study Born–Infeld particles with standard methods of nonlinear dynamics in the simplest case of a flat spacetime. We begin with the free Born–Infeld particles, i.e. BIon and catenoid . Using a scale transformation argument we show that these static configurations – which differ from ordinary solitons of nonlinear theories in requiring a special consideration of source terms or boundary conditions (cf. also ,) – require the number of space dimensions $`p`$ to be larger than 2. We assume spherical symmetry and study the local stability of these configurations by considering the second variational derivatives of their respective actions. Our conditions for stability are a) that the eigenfunctions of the corresponding operator be square integrable, and b) that the charge $`e`$ be fixed, with angular fluctuations ignored. We then consider the case of the scalar field corresponding to a single transverse coordinate coupled to the gauge field (here only the electric component), i.e. the catenoid or brane with associated open fundamental string. We distinguish between two types of arguments in deriving the linearised fluctuation equation, and infer the stability of this stringy $`D`$-brane. In ref. an explicit and detailed consideration of the Bogomol’nyi bound in a special model of Born–Infeld theory has been given where the central charge of the supersymmetry algebra plays the role of the topological or winding number of ordinary solitons.
Our second intention in the following is the explicit study of the small fluctuation equation about the $`D3`$–brane in the high energy domain. This equation with singular potential has the remarkable property of being convertible into a modified Mathieu equation which depends only on one coupling parameter which is a product of energy and electric charge. The $`S`$–matrix for scattering of the fluctuation off the brane can be obtained in explicit form. The $`D3`$–brane is therefore one of the very rare examples allowing a detailed study of its properties with explicit expressions for all relevant physical quantities in both low and high energy domains. We therefore expect that also $`S`$–duality can be uncovered and studied in this case (although we do not attempt this here). Various other $`Dp`$–brane models have been discovered recently whose small fluctuation equations can be reduced to modified Mathieu equations which have then been investigated mainly by computational methods. For the AdS/CFT correspondence the logarithmic corrections to the low energy absorption probability are of particular interest, since these permit a direct relation to the discontinuity of the cut in the correlation function of the dual two–dimensional quantum field theory. The first such logarithmic correction to the absorption probability was originally obtained in refs. without resorting to the use of Mathieu functions. Subsequently the authors of ref. considered the modified Mathieu equation and used computational methods to generate explicit series expansions up to several orders for the low energy absorption probability. In a different choice of expansions was considered to obtain leading expressions more easily. It is natural to supplement such investigations by exploring also the high energy case, the first such consideration being that of ref.. The analytical high energy results obtained in the following and the complementary low energy results of ref. (we also demonstrate how the $`S`$–matrices are related) are therefore directly applicable to these. Singular potentials have been studied from time to time, and have mostly been discarded as pathological. It seems, however, that their real significance lies in the context of curved spaces with black–hole type of absorption .
Sections 2 and 3 deal with the BIon and the catenoid, sections 4 and 5 with the Bogomol’nyi limit of the $`D3`$–brane and the derivation of the linearised fluctuation equation about it. In section 6 we consider this equation in detail in the high energy domain and calculate the rate of absorption of partial waves of the fluctuation field by the brane. That this absorption occurs is attributed to the singularity of the potential. The absorptivity part of the paper may be looked at as the high energy complement to the low energy case of ref. with the same expression of the $`S`$–matrix. All these calculations require a matching of wave functions. In the low energy $`S`$–wave case simple considerations of Bessel and Hankel functions suffice as was shown in refs. . The low energy limit is, in fact, independent of the choice of matching point, as was shown recently . Our considerations here, however, are general.
## 2 The BIon
We consider first purely static cases and write the Lagrangian of the static BIon in $`p+1`$ spacetime dimensions (cf.)
$$L=d^px,=1\sqrt{1(_i\varphi )^2}\mathrm{\Sigma }_pe\varphi \delta (𝐫),\mathrm{\Sigma }_p=\frac{p\pi ^{\frac{p}{2}}}{(\frac{p}{2})!}$$
(1)
($`i=1,\mathrm{},p`$) with the charge $`e`$ held fixed by the constraint
$$e+\frac{1}{\mathrm{\Sigma }_p}𝑑\sigma _i\frac{_i\varphi }{\sqrt{1(_j\varphi )^2}}=0$$
(2)
Eq. (1) is the Lagrangian one obtains from the world brane action of the pure Born–Infeld $`U(1)`$ electromagnetic action reduced to the purely electric case with field $`E_i=_0A_i_iA_0`$ and no transverse coordinate. The field $`A_\mu `$ is assumed to depend on the world brane coordinates $`x_\mu ,\mu =0,\mathrm{},p`$. The static BIon equation of motion is
$$_i\left(\frac{_i\varphi }{\sqrt{1(_i\varphi )^2}}\right)=\mathrm{\Sigma }_pe\delta (𝐫)$$
(3)
In the special case $`p=3`$ the classical $`SO(3)`$ symmetric solution, called a BIon, is given by
$$\varphi _c(r)=_r^{\mathrm{}}\frac{dx}{\sqrt{1+\frac{x^4}{e^2}}}=\varphi _c(0)_0^r\frac{dx}{\sqrt{1+\frac{x^4}{e^2}}}\stackrel{r0}{}\left[\varphi _c(0)r+\frac{r^5}{10e^2}\right]$$
(4)
and $`\varphi _c(0)=\frac{1}{4}B(\frac{1}{4},\frac{1}{4}).e^{\frac{1}{2}}=1.854074677.e^{\frac{1}{2}}`$, $`B`$ being the Bernoulli function. It is easily verified that this solution satisfies the constraint (2) for any value of $`r`$. Defining $`𝐄=\varphi _c`$ (so that $`\varphi _c=A_0`$ with $`A_0(x_i,t)/t=0`$ in the static case), and defining $`𝐃=\frac{}{𝐄}=\frac{𝐄}{\sqrt{1𝐄^2}}`$ we have (with $`F_{0i}=E_i`$)
$$T_{00}=F_{0i}\frac{}{F_{0i}}=𝐄𝐃=\frac{1}{\sqrt{1𝐄^2}}1+4\pi e\varphi \delta (𝐫)$$
(5)
The energy $`H_c`$ of the BIon (obtained by integration over $`𝐑^3`$) is then found to be finite, i.e.
$$H_c=𝑑𝐱T_{00}=4\pi (3.09112).e^{\frac{3}{2}}$$
(6)
and in $`p`$ dimensions the total energy of the BIon scales correspondingly as $`e^{\frac{p}{p1}}`$. The finiteness of the energy depends on the minus sign in (1) and so with (3) on the relation
$$\sqrt{1\left(\varphi _c^{}\right)^2}=\frac{r^2}{e}\varphi _c^{}=\frac{r^2}{e\sqrt{1+\frac{r^4}{e^2}}}$$
(7)
for $`0r\mathrm{}`$. It may be noted that by defining $`𝐃`$ such that the left hand side of eq.(3) is $`_iD_i`$, the singularity of the right hand side is associated with $`𝐃`$ rather than with $`𝐄`$ which is the decisive difference between Maxwell and Born–Infeld electrodynamics. A similar observation applies to the catenoid equation below. The energy of the BIon is seen to be independent of its position which hints at the existence of some kind of collective coordinate. However, exploring this point further is expected to be difficult since a moving charge generates a magnetic field, and hence the electric field alone would not suffice.
We can use a scaling argument to show that here finite energy configurations require $`p`$ to be larger than or equal to $`3`$. Under a scale transformation $`xx^{}=\lambda x,\varphi (x)\varphi _\lambda (x)=\varphi (\lambda x),_i\varphi (x)[_i\varphi (x)]_\lambda =\lambda _i\varphi (\lambda x)`$. The charge $`e`$ defined by the constraint (2) also changes under the scale transformation, i.e.
$$ee_\lambda =\frac{1}{\lambda ^{p2}\mathrm{\Sigma }_p}\frac{d\sigma _i_i\varphi }{\sqrt{1\lambda ^2(_j\varphi )^2}}$$
(8)
In particular for $`p=3`$ and radial symmetry
$$e_\lambda ^{(p=3)}=\frac{r^2}{\lambda }\frac{1}{\sqrt{1\lambda ^2+\frac{r^4}{e^2}}}$$
(9)
and for arbitrary values of $`\lambda `$ the $`r`$–dependence drops out only if the limit $`r\mathrm{}`$ is taken in the evaluation of the integral. Then
$$\frac{e_\lambda ^{(p=3)}}{e}\stackrel{r\mathrm{}}{}\frac{1}{\lambda }$$
(10)
But also $`e_{\lambda =1}=e`$ for any $`r`$. If $`\varphi _c`$ is stable and $`0`$, the energy must be stationary for $`\lambda =1`$, i.e. $`(H_c/\lambda )_{\lambda =1}=0`$. From this one finds that $`p3`$. Also $`(^2H_c/\lambda ^2)_{\lambda =1}>0`$ for $`p3`$. Eqs.(6) and (10) show that changing the scale changes both the charge and the energy, i.e. if the charge were variable, one could lower the energy and hence the configuration could be unstable. But fixing the charge (e.g. by a quantisation condition) no instability is implied by the scaling condition.
We investigate the stability of the BIon further in the special and exemplary case of $`p=3`$ by considering the second functional variation of the static Lagrangian evaluated at $`\varphi _c(r)`$. This can be written and simplified in the following form (ignoring total divergences on the way)
$$\delta ^2L=\frac{1}{2}d^3x\delta \varphi \widehat{A}\delta \varphi ,$$
(11)
where
$`\widehat{A}`$ $`=`$ $`_i{\displaystyle \frac{1}{[1(_j\varphi _c)^2]^{1/2}}}_i_i{\displaystyle \frac{_i\varphi _c_j\varphi _c}{[1(_k\varphi _c)^2]^{3/2}}}_j`$ (12)
$`=`$ $`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{d}{dr}}{\displaystyle \frac{r^2}{(1\varphi _{c}^{}{}_{}{}^{2})^{3/2}}}{\displaystyle \frac{d}{dr}}`$
The operator $`\widehat{A}`$ can also be written
$$\widehat{A}=\frac{1}{(1\varphi _{c}^{}{}_{}{}^{2})^{3/2}}\left\{\frac{1}{r^2}\frac{d}{dr}r^2\frac{d}{dr}\frac{6}{r}\varphi _{c}^{}{}_{}{}^{2}\frac{d}{dr}\right\}$$
(13)
The classical stability of $`\varphi _c`$ is therefore decided by the spectrum $`\{\omega _n\}`$ of the small fluctuation equation
$$\frac{1}{r^2}\frac{d}{dr}\frac{r^2}{(1\varphi _{c}^{}{}_{}{}^{2})^{3/2}}\frac{d}{dr}\psi _n=\omega _n\psi _n$$
(14)
We explore first the existence of a zero mode $`\psi _0`$, i.e. the case $`\omega =0`$. In this case
$$\frac{r^2}{(1\varphi _{c}^{}{}_{}{}^{2})^{3/2}}\frac{d}{dr}\psi _0=C$$
(15)
and so with $`\psi _0(\mathrm{})=0`$
$$\psi _0=\frac{C}{e^3}_r^{\mathrm{}}\frac{x^4}{(1+\frac{x^4}{e^2})^{3/2}}𝑑x=C\frac{}{e}_r^{\mathrm{}}\frac{dx}{(1+\frac{x^4}{e^2})^{1/2}}=C\frac{\varphi _c}{e}$$
(16)
The derivative of the classical configuration $`\varphi _c`$ with respect to the charge $`e`$ indicates that a perturbation along $`\varphi _c/e`$ around $`\varphi _c`$leaves the static action invariant, i.e. $`\varphi _c(e,r)`$ and $`\varphi _c(e+\delta e,r)`$ have the same action since
$$\frac{\varphi _c(e+\delta e,r)}{(\delta e)}|_{\delta e=0}=\frac{\varphi _c(e,r)}{e}.$$
We now show that the operator $`\widehat{A}`$ does not possess negative eigenvalues, and that therefore the BIon is a classically stable configuration. We let $`\psi _n`$ be an eigenfunction of the operator $`\widehat{A}`$. Then
$`{\displaystyle d^3x\psi _n\widehat{A}\psi _n}`$ $`=`$ $`4\pi {\displaystyle _0^{\mathrm{}}}𝑑r\psi _n{\displaystyle \frac{d}{dr}}{\displaystyle \frac{r^2}{(1\varphi _{c}^{}{}_{}{}^{2})^{3/2}}}{\displaystyle \frac{d\psi _n}{dr}}`$ (17)
$`=`$ $`4\pi {\displaystyle _0^{\mathrm{}}}𝑑r\left\{{\displaystyle \frac{d}{dr}}\psi _n{\displaystyle \frac{r^2}{(1\varphi _{c}^{}{}_{}{}^{2})^{3/2}}}{\displaystyle \frac{d\psi _n}{dr}}{\displaystyle \frac{r^2}{(1\varphi _{c}^{}{}_{}{}^{2})^{3/2}}}\left({\displaystyle \frac{d\psi _n}{dr}}\right)^2\right\}`$
$`=`$ $`F+4\pi {\displaystyle _0^{\mathrm{}}}𝑑r{\displaystyle \frac{r^2}{(1\varphi _{c}^{}{}_{}{}^{2})^{3/2}}}\left({\displaystyle \frac{d\psi _n}{dr}}\right)^2`$
where $`F:=F(r)|_0^{\mathrm{}}`$ and
$$F(r)=4\pi \psi _n\frac{r^2}{(1\varphi _{c}^{}{}_{}{}^{2})^{3/2}}\frac{d\psi _n}{dr}=4\pi e^3\frac{\psi _n}{r^4}\left(1+\frac{r^4}{e^2}\right)^{3/2}\frac{d\psi _n}{dr}.$$
(18)
The second term on the right hand side of eq.(17) is strictly positive. Hence nonpositive eigenvalues imply a nonvanishing negative value of $`F`$. ¿From the condition $`_0^{\mathrm{}}\psi _n^2r^2𝑑r<\mathrm{}`$, (i.e. $`\psi _n\stackrel{r\mathrm{}}{}1/(r^{1+ϵ}),ϵ>0`$) it follows that $`r^2\psi _nd\psi _n/dr0`$ with $`r\mathrm{}`$, so that
$$F(r)\stackrel{r\mathrm{}}{}4\pi r^2\psi _n\frac{d\psi _n}{dr}0$$
and $`F(\mathrm{})=0`$. Hence
$$F=F(0)4\pi e^3\frac{\psi _n}{r^4}\frac{d\psi _n}{dr}|_{r0}$$
(19)
As $`r0`$ eq.(14) becomes
$$\frac{1}{r^2}\frac{d}{dr}\frac{1}{r^4}\frac{d}{dr}\psi _n=\frac{\omega _n}{e^3}\psi _n$$
(20)
In the case of the zero mode
$$\psi _0C_1+C_2r^5,r0$$
(21)
In this case $`F=20\pi e^3C_1C_2`$. For $`C_1C_2<0`$ this is in full compliance with (16) and (4) from which we obtain
$$\psi _0C\left(\frac{B(\frac{1}{4},\frac{1}{4})}{8e^{1/2}}\frac{r^5}{5e^3}\right).$$
For $`\omega _n0`$ the small–r behaviour of $`\psi _n`$ is
$$\psi _nC_n\left(1\frac{1}{24}\frac{\omega _n}{e^3}r^8+O(r^{16})\right)$$
(22)
so that
$$F=\frac{4}{3}\pi C_{n}^{}{}_{}{}^{2}\frac{1}{r^4}r^7|_{r=0}=0$$
Thus the conclusion is that for all eigenfunctions $`\psi _n`$
$$<\psi _n|\widehat{A}|\psi _n>0$$
(23)
This inequality excludes the possibility of the existence of negative eigenvalues. Hence the BIon is in this sense classically stable.
## 3 The catenoid
The Lagrangian of the static catenoid in $`p+1`$ spacetime dimensions and with a source term is given by (cf.)
$$L=d^px,=1\sqrt{1+(_iy)^2}\mathrm{\Sigma }_pr_0^{p1}y\delta (𝐫)$$
(24)
where the signs have been chosen such that the energy is positive. Here the scalar field $`y(x_i,t)`$ originates from gauge field components $`A_a`$ for $`a=p+1,\mathrm{},(d1),d=`$dimension, which represent transverse displacements of the brane; here we consider the case of only one such transverse coordinate, i.e. $`y`$, all $`dp1`$ of which are essentially Kaluza–Klein remnants of the $`d=10`$ dimensional $`N=1`$ electrodynamics after dimensional reduction to $`p+1`$ dimensions. The Euler–Lagrange equation of the static catenoid $`y_c`$ (static meaning $`y(x_i,t)/t=0`$) is given by
$$_i\left(\frac{_iy_c}{\sqrt{1+\left(y_c^{}\right)^2}}\right)=\mathrm{\Sigma }_pr_0^{p1}\delta (𝐫)$$
(25)
so that after integration
$$\frac{y_c}{\sqrt{1+\left(y_c^{}\right)^2}}=r_0^{p1}\frac{𝐫}{r^p}$$
(26)
or for $`rr_0`$
$$y_c^{}=\stackrel{}{(+)}\frac{r_{0}^{}{}_{}{}^{p1}}{\sqrt{r^{2p2}r_{0}^{}{}_{}{}^{2p2}}},\sqrt{1+y_{c}^{}{}_{}{}^{2}}=\stackrel{}{(+)}\frac{r^{p1}}{\sqrt{r^{2p2}r_0^{2p2}}}$$
(27)
In the case of the catenoid without source term the right hand side of eq.(26) can be taken to originate from a boundary condition such as $`\left(\frac{𝐫}{r^p}\right)=0`$. The domain $`rr_0`$ is the nonsingular throat region (i.e. $`y_c(r_0))`$ is finite). One may observe that the singularity on the right hand side of eq.(25) is associated with the entire expression on the left whereas, like $`_i\varphi _c`$ in the BIon case, so now here $`y_c`$ is finite, i.e. the $`p`$-brane or single throat solution is given by
$$y_c(r)=\stackrel{}{(+)}_r^{\mathrm{}}dr\frac{r_0^{p1}}{\sqrt{r^{2p2}r_0^{2p2}}}$$
(28)
Thus $`y`$ is double valued. The two possible signs can be taken to define a brane and its antibrane. We show at the end of this section that the solution with the minus sign is the minimum of the action and the solution with the plus sign the maximum of the action. This function is finite at $`r=r_0`$ and can be expressed in terms of elliptic integrals. For $`r_0=1`$ it is even simpler and has the value $`y_c(1)=\stackrel{}{(+)}\frac{1}{\sqrt{2}}𝒦(\frac{1}{\sqrt{2}})`$ where $`𝒦`$ is the complete elliptic integral of the first kind. Plotted as a function of $`r`$, $`y_c(r)`$ is a monotonically decreasing function starting from $`r_0`$; pictured on a 2–dimensional space it looks like an inverted funnel (i.e. the surface swept out by a catenary with boundaries at the openings), thus suggesting the name catenoid. As pointed out in ref., the two possible signs of the square root allow a smooth joining of one such funnel–shaped branch to an inverted one connected by a throat of finite thickness, the resulting structure then representing a brane–antibrane pair. This brane–antibrane pair is joined by the throat of finite thickness $`r_0`$ and finite length. In fact, we can rewrite eq.(28) in terms of $`\stackrel{~}{y}_c(x)=y_c(r_0x),x=\frac{r}{r_0}`$, and for the special case of $`p=3`$ as
$`\stackrel{~}{y}_c(x)`$ $`=`$ $`\stackrel{}{(+)}{\displaystyle _x^{\mathrm{}}}{\displaystyle \frac{dx}{\sqrt{x^41}}}=\stackrel{}{(+)}{\displaystyle _1^{\mathrm{}}}{\displaystyle \frac{dx}{\sqrt{x^41}}}\stackrel{+}{()}{\displaystyle _1^x}{\displaystyle \frac{dx}{\sqrt{x^41}}}`$ (29)
$`=`$ $`\stackrel{}{(+)}{\displaystyle \frac{1}{\sqrt{2}}}\left[𝒦\left({\displaystyle \frac{\sqrt{2}}{2}}\right)cn^1({\displaystyle \frac{1}{x}},{\displaystyle \frac{\sqrt{2}}{2}})\right]`$
where $`x>1`$ and we used formulae of ref.. Inverting this expression we obtain the periodic function
$$x(y)=\left[cn\left(𝒦\left(\frac{\sqrt{2}}{2}\right)\stackrel{+}{()}\sqrt{2}y,\frac{\sqrt{2}}{2}\right)\right]^1$$
(30)
Plotting this expression with $`x`$ as ordinate, one obtains the picture of a cross section through a chain of periodically recurring funnel–shaped structures to the one side of the throat, i.e. the series $`\mathrm{}`$ representing a series of brane–antibrane pairs along the abscissa.
Proceeding as in the above case of the static BIon and calculating the second variational derivative we obtain
$$\delta ^2L=\frac{1}{2}d^px\delta y\widehat{B}\delta y$$
(31)
where for $`rr_0`$
$`\widehat{B}`$ $`=`$ $`_i{\displaystyle \frac{i}{\left[1+(_iy_c)^2\right]^{1/2}}}_i_i{\displaystyle \frac{_iy_c_jy_c}{\left[1+(_iy_c)^2\right]^{3/2}}}_j`$ (32)
$`=`$ $`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{d}{dr}}{\displaystyle \frac{r^2}{\left(1+y_{c}^{}{}_{}{}^{2}\right)^{3/2}}}{\displaystyle \frac{d}{dr}}=\stackrel{}{(+)}{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{d}{dr}}{\displaystyle \frac{(r^4r_0^4)^{3/2}}{r^4}}{\displaystyle \frac{d}{dr}}`$
The operator $`\widehat{B}`$ can also be written
$$\widehat{B}=\frac{1}{(1+y_{c}^{}{}_{}{}^{2})^{3/2}}\left\{\frac{1}{r^2}\frac{d}{dr}r^2\frac{d}{dr}+\frac{6}{r}y_{c}^{}{}_{}{}^{2}\frac{d}{dr}\right\}$$
(33)
Since the gauge field components $`A_a,a=p+1,\mathrm{},d1`$ (of which we retain only one), are dynamical, the Lagrangian in the nonstatic case is
$$=1\sqrt{1(_\mu y)(^\mu y)}\mathrm{\Sigma }_pr_0^{p1}\delta (𝐫)$$
(34)
and we can obtain the same condition of stability by considering the dynamical fluctuation $`\eta `$, i.e.
$$y(t,𝐱)=y_c(r)+\eta (t,𝐱),\eta =\xi (r)e^{i\sqrt{\omega }t}$$
and linearising the time–dependent Euler–Lagrange equation. The square integrable perturbations $`\xi (r)`$ are the socalled “$`L^2`$ deformations” of ref.. The classical stability of $`y_c`$ is therefore decided by the spectrum $`\{\omega \}`$ of the small fluctuation equation
$$\frac{1}{r^2}\frac{d}{dr}\frac{r^2}{(1+y_{c}^{}{}_{}{}^{2})^{3/2}}\frac{d}{dr}\psi =\frac{1}{r^2}\frac{d}{dr}\left\{\stackrel{}{(+)}\frac{(r^4r_0^4)^{3/2}}{r^4}\right\}\frac{d}{dr}\psi =\omega \psi $$
(35)
We explore first the existence of a zero mode $`\psi _0`$, i.e. the case $`\omega =0`$. In this case
$$\frac{r^2}{(1+y_{c}^{}{}_{}{}^{2})^{3/2}}\frac{d}{dr}\psi _0=C$$
(36)
and so in the case $`p=3`$ and $`rr_0`$
$$\psi _0=C_r^{\mathrm{}}𝑑x\frac{x^4}{(x^4r_0^4)^{3/2}}=\frac{C}{2r_0}\frac{y_c}{r_0}$$
(37)
so that
$$\psi _0=C\frac{y_c}{r_{0}^{}{}_{}{}^{2}}$$
Here again the derivative of the classical configuration $`y_c`$ with respect to the parameter $`r_{0}^{}{}_{}{}^{2}`$ is indicative of stationarity of the action in a shift of $`r_0^2`$.
We now demonstrate that the operator $`\widehat{B}`$ with the minus sign has no negative eigenvalues, and that therefore the free catenoid is a classically stable configuration like the BIon for fixed throat radius $`r_0`$. Then
$`{\displaystyle d^3x\psi \widehat{B}\psi }`$ $`=`$ $`4\pi {\displaystyle _{r_0}^{\mathrm{}}}𝑑r\psi {\displaystyle \frac{d}{dr}}{\displaystyle \frac{r^2}{(1+y_{c}^{}{}_{}{}^{2})^{3/2}}}{\displaystyle \frac{d\psi }{dr}}`$
$`=`$ $`4\pi {\displaystyle \frac{r^2}{(1+y_{c}^{}{}_{}{}^{2})^{3/2}}}\psi {\displaystyle \frac{d\psi }{dr}}|_{r_0}^{\mathrm{}}4\pi {\displaystyle _{r_0}^{\mathrm{}}}𝑑r{\displaystyle \frac{r^2}{(1+y_{c}^{}{}_{}{}^{2})^{3/2}}}\left({\displaystyle \frac{d\psi }{dr}}\right)^2`$
$`=`$ $`\stackrel{}{(+)}4\pi {\displaystyle \frac{(r^4r_0^4)^{3/2}}{r^4}}\psi {\displaystyle \frac{d\psi }{dr}}|_{r_0}^{\mathrm{}}\stackrel{+}{()}4\pi {\displaystyle _{r_0}^{\mathrm{}}}𝑑r{\displaystyle \frac{(r^4r_{0}^{}{}_{}{}^{2})^{3/2}}{r^4}}\left({\displaystyle \frac{d\psi }{dr}}\right)^2`$
where we used eq.(27). The second term is always positive if the upper sign is chosen. The first term vanishes at infinity with $`𝑑rr^2\psi ^2<\mathrm{}`$, since
$$4\pi \frac{(r^4r_0^4)^{3/2}}{r^4}\psi \frac{d\psi }{dr}\stackrel{r\mathrm{}}{}4\pi r^2\psi \frac{d\psi }{dr}0.$$
On the other hand, in the case $`rr_0`$, we have
$$4\pi \frac{(r^4r_0^4)^{3/2}}{r^4}\psi \frac{d\psi }{dr}32\pi \sqrt{r_0}(rr_0)^{3/2}\psi \frac{d\psi }{dr}$$
As $`rr_0`$ eq.(35) becomes
$$\frac{8}{r_0^{3/2}}\frac{d}{dr}(rr_0)^{3/2}\frac{d\psi }{dr}=\omega \psi $$
(39)
In the case of the zero mode $`\psi _0`$ with $`\omega =0`$ the considerations are analogous to those of the BIon case and the sum of the two terms in eq.(LABEL:38) vanishes. In the case of $`\omega 0`$ we therefore have
$$\psi C\left(1\frac{\omega }{4}r_{0}^{}{}_{}{}^{3/2}\sqrt{rr_0}\right)$$
and
$$\underset{rr_0}{lim}(rr_0)^{3/2}\psi \frac{d\psi }{dr}=\frac{\omega }{8}C^2.\underset{rr_0}{lim}r_{0}^{}{}_{}{}^{3/2}(rr_0)=0$$
This proves that for all eigenfunctions $`\psi `$
$$<\psi |\widehat{B}|\psi >0.$$
Thus $`\widehat{B}`$ has no negative eigenvalues, and the free throat is classically stable with fixed $`r_0`$ for the sign chosen as in eq.(LABEL:38). Obviously the operator $`\widehat{B}`$ with the plus sign has no positive eigenvalues, which means that we have the maximum of the action. Of course, if we change $`r_0`$ (and so consider a different theory), the expectation value of $`\widehat{B}`$ also changes. One should note that the free throat we discuss here is that with vanishing gauge field. The double valuedness of the solution of eq.(25) implies that if one solution is classically stable, the other one is not. Thus a multi–throat solution constructed from these by matching both solutions, if it exists, like the brane–antibrane solution of ref., is expected to be unstable in view of negative as well as positive eigenvalues, and is therefore neither a maximum nor a minimum of the action. In fact, as argued in ref. (after eq.(132)) equilibrium between these should not be possible. The reason for this is that a symmetrical configuration, symmetrical about the plane $`x_3=0`$ for instance, implies $`_3y=0`$ there. Evaluating the stress tensor element $`T_{33}`$ (even for vanishing gauge field), one obtains a negative quantity which is interpreted as implying an attractive force between the brane and its antibrane in this symmetrically constructed configuration. This is, in fact, the general instability of this configuration discussed in ref..
## 4 Coupled fields: The $`D`$–brane in the Bogomol’nyi limit
In the case of coupled fields $`\varphi `$ and $`y`$ (the former with source, the latter without), the Lagrangian of the static case is (cf.)
$`L`$ $`=`$ $`{\displaystyle d^px},=1Q\mathrm{\Sigma }_pe\varphi \delta (𝐫),`$
$`Q`$ $`=`$ $`[1(_i\varphi )^2+(_iy)^2+(_i\varphi ._iy)^2(_i\varphi )^2(_jy)^2]^{\frac{1}{2}}`$ (40)
¿From the first variation of $`L`$ we obtain the coupled equations of the fields $`\varphi `$ and $`y`$, i.e. from
$`\delta L=`$ $``$ $`{\displaystyle \delta \varphi _i\frac{1}{Q}[_i\varphi (\varphi y)_iy+(y)^2_i\varphi ]d^px}`$ (41)
$`+`$ $`{\displaystyle \delta y_i\frac{1}{Q}[_iy+(\varphi y)_i\varphi (\varphi )^2_iy]d^px}`$
$``$ $`\mathrm{\Sigma }_pe{\displaystyle \delta \varphi \delta (𝐫)d^px}`$
(ignoring total divergences).
The source term of the electric field again suggests spherical symmetry. In deriving the two coupled Euler–Lagrange equations one new constant $`c`$ (apart from $`e`$) arises in the integration of the catenoid equation, i.e.
$$_r\left(r^{p1}\frac{}{(_ry)}\right)=0,r^{p1}\frac{}{(_ry)}=c$$
We have no source term of the $`y`$ field because, as before, the appropriate effect is provided by the boundary condition defining the width of the throat. The two equations with spherical symmetry are found to be
$$\frac{\varphi ^{}}{\left[1(\varphi ^{})^2+(y^{})^2\right]^{\frac{1}{2}}}=\frac{e}{r^{p1}},\frac{y^{}}{\left[1(\varphi ^{})^2+(y^{})^2\right]^{\frac{1}{2}}}=\frac{c}{r^{p1}}$$
(42)
so that
$$\frac{\varphi ^{}}{y^{}}=\frac{e}{c}\frac{1}{a}$$
(43)
Then
$$(\varphi ^{})^2=\frac{1}{\frac{r^{2(p1)}}{e^2}+1a^2},(y^{})^2=\frac{a^2}{\frac{r^{2(p1)}}{e^2}+1a^2}$$
(44)
Thus the family of solutions can be parametrised in terms of the single parameter $`a`$ as already pointed out in ref. . This parameter is seen to interpolate between the two types of static solutions. The solution $`y`$ of (35) for various values of $`a^2`$ is now the $`p`$-brane, i.e.
$$y(r)=\stackrel{+}{()}ae_r^{\mathrm{}}dr\frac{1}{\sqrt{r^{2(p1)}r_0^{2(p1)}}}$$
(45)
where $`r_0^{2(p1)}=e^2(a^21)`$ and for the solution to make sense we must have $`a^21`$. If $`ae`$ in eq.(36) is replaced by $`ae`$, the expression represents the corresponding antibrane. Taking $`e^20,a^2e^2`$ const. the electric field is eliminated and we regain the free catenoid solution. In approaching the limit $`a^21`$ the width of the throat becomes infinitesimal with nonvanishing electric field and the configuration can then be considered to be a fundamental string, as argued in ref.. We distinguish between three cases:
$`|a|<1`$ $`:`$
$`\varphi ={\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{dx}{\sqrt{1a^2+x^4/e^2}}},y=a{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{dx}{\sqrt{1a^2+x^4/e^2}}},`$
$`|a|>1`$ $`:`$
$`\varphi =e{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{dx}{\sqrt{x^4r_{0}^{}{}_{}{}^{4}}}},y=ae{\displaystyle _r^{\mathrm{}}}{\displaystyle \frac{dx}{\sqrt{x^4r_{0}^{}{}_{}{}^{4}}}}`$
$`a=\pm 1`$ $`:`$ $`\varphi ={\displaystyle \frac{e}{r}},y=\pm {\displaystyle \frac{e}{r}}`$ (46)
We see that for $`a^2=1`$ eq. (34) becomes the first order Bogomol’nyi equation or linearised field equation for $`y`$ (as in ref.)
$$F_{0r}\pm \frac{y}{r}=0$$
(47)
where $`F_{0r}=E_c`$ is the static electric field. This is the same equation as that obtained from the vanishing of the supersymmetry variation of the gaugino field $`\mathrm{\Sigma }`$ for half the number of 16 supersymmetries (for $`d=10`$ and $`p=3`$) $`ϵ_+,ϵ_{}`$ of $`ϵ`$ for which $`\delta \mathrm{\Sigma }=0`$, i.e.
$$\delta _+\mathrm{\Sigma }=0,\delta _{}\mathrm{\Sigma }0$$
where – as discussed in the literature $`ϵ`$ is the constant spinor of the supersymmetry variation and $`ϵ_\pm `$ are its chiral components. Thus $`a^2=1`$ implies BPS configurations, wheras those with $`a^21`$ are non–BPS. Taking $`a^2=0`$ in eq.(36) we regain the BIon configuration as a local minimum of the energy whereas for vanishing electric field one expects a local maximum, i.e. a sphaleron configuration (as pointed out in ).
Next we investigate the second variation of the static $`L`$ with spherical symmetry. We set
$$\delta ^2L=\frac{1}{2}\left\{\delta \varphi \widehat{M}\delta \varphi +\delta y\widehat{N}\delta y+\delta \varphi \widehat{L}\delta y+\delta y\widehat{L}^{}\delta \varphi \right\}d^px$$
(48)
Again ignoring total divergences one finds
$`\widehat{M}`$ $`=`$ $`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{d}{dr}}r^2{\displaystyle \frac{1+y_{}^{}{}_{}{}^{2}}{Q^3}}{\displaystyle \frac{d}{dr}},`$
$`\widehat{N}`$ $`=`$ $`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{d}{dr}}r^2{\displaystyle \frac{1\varphi _{}^{}{}_{}{}^{2}}{Q^3}}{\displaystyle \frac{d}{dr}},`$
$`\widehat{L}`$ $`=`$ $`{\displaystyle \frac{1}{r^2}}{\displaystyle \frac{d}{dr}}r^2{\displaystyle \frac{\varphi ^{}y^{}}{Q^3}}{\displaystyle \frac{d}{dr}}`$
with $`\widehat{L}=\widehat{L}^{}`$. We can now rewrite $`\delta ^2L`$ as
$$\delta ^2L=\frac{1}{2}d^3x(\delta \varphi ,\delta y)\widehat{H}\left(\begin{array}{c}\delta \varphi \\ \delta y\end{array}\right)$$
(50)
where
$$\widehat{H}=\left(\begin{array}{c}\begin{array}{cc}M& L\\ L^{}& N\end{array}\end{array}\right)=\frac{1}{r^2}\frac{d}{dr}r^2h\frac{d}{dr}$$
(51)
and
$$h=\frac{1}{Q^3}\left(\begin{array}{c}\begin{array}{cc}1y_{}^{}{}_{}{}^{2}& y^{}\varphi ^{}\\ y^{}\varphi ^{}& 1\varphi _{}^{}{}_{}{}^{2}\end{array}\end{array}\right),\widehat{H}^{}=\widehat{H},$$
(52)
with
$$h^1=Q\left(\begin{array}{c}\begin{array}{cc}1+\varphi _{}^{}{}_{}{}^{2}& y^{}\varphi ^{}\\ y^{}\varphi ^{}& 1+y_{}^{}{}_{}{}^{2}\end{array}\end{array}\right),deth=\frac{1}{Q^4},$$
(53)
The small fluctuation equation therefore becomes
$$\widehat{H}\psi =\frac{1}{r^2}\frac{d}{dr}r^2h\frac{d}{dr}\psi =\omega \psi $$
(54)
Again we first explore the existence of a zero mode $`\psi _0`$ with
$$r^2h\frac{d}{dr}\psi _0=\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)$$
(55)
where $`\alpha `$ and $`\beta `$ are constants. Setting
$$\psi _0=\left(\begin{array}{c}\varphi _0\\ y_0\end{array}\right)$$
and evaluating $`\psi _0`$ for the solutions of eq.(46) we obtain with
$$\phi =_r^{\mathrm{}}\varphi _{}^{}{}_{}{}^{3}(x)𝑑x$$
the relation
$$\psi _0=\frac{\varphi }{e}\left\{\left(\begin{array}{c}\alpha \\ \beta \end{array}\right)\frac{\phi }{e}(\alpha +a\beta )\left(\begin{array}{c}1\\ a\end{array}\right)\right\}$$
(56)
In the BPS limit with $`y^{}=\varphi ^{}=E_c,Q=1`$, the operator $`\widehat{H}`$ of eq.(42) becomes
$$\widehat{H}=\frac{1}{r^2}\frac{d}{dr}r^2\left(\begin{array}{c}\begin{array}{cc}1E_c^2& E_c^2\\ E_c^2& 1E_c^2\end{array}\end{array}\right)\frac{d}{dr}$$
(57)
Setting
$$\psi _s=\left(\begin{array}{c}\delta \varphi \\ \delta y\end{array}\right)=\rho (x)\left(\begin{array}{c}1\\ 1\end{array}\right)$$
for an arbitrary function $`\rho (x)`$ we have
$$\widehat{H}\psi _s=\frac{1}{r^2}\frac{d}{dr}r^2\rho ^{}\left(\begin{array}{c}\begin{array}{cc}1E_c^2& E_c^2\\ E_c^2& 1E_c^2\end{array}\end{array}\right)\left(\begin{array}{c}1\\ 1\end{array}\right)=\frac{1}{r^2}\frac{d}{dr}r^2\rho ^{}\left(\begin{array}{c}1\\ 1\end{array}\right)$$
(58)
Thus for arbitrary $`\rho (x)`$, we have $`\psi _s\widehat{H}\psi _s=0`$ implying $`\delta ^2L=0`$ or $`L`$ constant in a specific direction about the BPS configuration. This behaviour may be interpreted as indicative of a local symmetry, in this case of supersymmetry, and so of the cancellation of fermionic and bosonic contributions in the one loop approximation. Here, of course, we have no fermionic contributions and consequently those of the two bosonic fields have opposite signs.
## 5 Fluctuations about the $`D`$–brane
In the following we distinguish clearly between two different types of fluctuations. We consider the above BPS solution for the string attached to the 3–brane as background and consider first a scalar field propagating in a direction along the string and perpendicular to the brane and its anti–brane. The linearised equation of small fluctuations about this background is obtained from the second variational derivative of the action which is the standard procedure and we therefore consider this first (cf. also ). Our treatment here is somewhat different (see below) from that in refs.. The resulting fluctuation equation has also been given in ref.. It is necessary to return to the fully time-dependent version, i.e.
$$S=\frac{1}{(2\pi )^pg_s}d^{p+1}x\left[1\sqrt{det(\eta _{\mu \nu }+F_{\mu \nu })}\mathrm{\Sigma }_pe\varphi \delta (𝐫)\right]$$
(59)
where in $`3+1`$ dimensions $`F_{\mu \nu }=F_{\mu \nu }(x_0,x_1,x_2,x_3)`$. In the electrostatic case with only one scalar field $`y`$ we have $`A_\mu =(A_0,A_1,A_2,A_3,y,0,0,0,0,0),F_{0i}=E_i`$ and $`F_{\mu 4}=_\mu y`$ for $`i=1,2,3`$ and $`\mu =0,1,2,3`$. Then
$$det(\eta _{\mu \nu }+F_{\mu \nu })=\left|\begin{array}{c}\begin{array}{cc}\begin{array}{ccc}\begin{array}{cccc}\begin{array}{ccccc}1& E_1& E_2& E_3& _0y\\ E_1& 1& 0& 0& _1y\\ E_2& 0& 1& 0& _2y\\ E_3& 0& 0& 1& _3y\\ _0y& _1y& _2y& _3y& 1\end{array}& & & \end{array}& & \end{array}& \end{array}\end{array}\right|$$
(60)
and so
$$det(\eta _{\mu \nu }+F_{\mu \nu })=(1𝐄^2)(1+y^2)(𝐄y)^2+(_0y)^2$$
(61)
We consider first the Lagrangian density (remembering that the relevant fields are $`A_0,A_i`$ and $`y`$)
$$=1Q,Q=\left[(1𝐄^2)(1+y^2)+(𝐄y)^2\dot{y}^2\right]^{\frac{1}{2}}$$
(62)
The equations of the static BIon and the static catenoid discussed above follow again from the first variations
$`{\displaystyle \frac{}{E_i}}`$ $`=`$ $`{\displaystyle \frac{1}{Q}}\left[E_i(1+y^2)_iy(𝐄y)\right],`$
$`{\displaystyle \frac{}{_iy}}`$ $`=`$ $`{\displaystyle \frac{1}{Q}}\left[_iy(1𝐄^2)+E_i(𝐄y)\right],`$
$`{\displaystyle \frac{}{_0y}}`$ $`=`$ $`{\displaystyle \frac{1}{Q}}`$ (63)
In the BPS background given by
$$_iy=E_i,𝐄^2=(y)^2=𝐄y=\frac{e^2}{r^4}E_c^2y_c^2,Q=1$$
(64)
one finds
$$\frac{^2}{E_iE_j}=(1+E_c^2)\delta _{ij},\frac{^2}{_iy_jy}=(1E_c^2)\delta _{ij},\frac{^2}{E_i_jy}=E_c^2\delta _{ij},\frac{^2}{\dot{y}^2}=1$$
(65)
This enables us to write (ignoring again total divergences in shifting derivatives)
$`\delta ^2=(\delta A_0,\delta A_i,\delta y)`$
$`\left(\begin{array}{c}\begin{array}{cc}\begin{array}{ccc}_i(1+E_c^2)_i& _i(1+E_c^2)_0& _iE_c^2_i\\ _0(1+E_c^2)_i& _0(1+E_c^2)_0& _0E_c^2_i\\ _iE_c^2_i& +_iE_c^2_0& _i(1E_c^2)_i_0_0\end{array}& \end{array}\end{array}\right)\left(\begin{array}{c}\delta A_0\\ \delta A_i\\ \delta y\end{array}\right)`$ (74)
In the linear approximation the Euler–Lagrange equations of the fluctuations $`\delta y\eta ,\delta E_i=_0\delta A_i_i\delta A_0`$ are therefore given by the following set of three equations
$`{\displaystyle \frac{d^2}{dt^2}}\eta +_i(1E_c^2)_i\eta +_iE_c^2(_0\delta A_i_i\delta A_0)`$ $`=`$ $`0,`$ (75)
$`{\displaystyle \frac{d}{dt}}(1+E_c^2)(_0\delta A_i_i\delta A_0){\displaystyle \frac{d}{dt}}E_c^2_i\eta `$ $`=`$ $`0,`$ (76)
$`_i(1+E_c^2)(_0\delta A_i_i\delta A_0)_iE_c^2_i\eta `$ $`=`$ $`0`$ (77)
The last of these three equations can be seen to be a constraint by appying $`/t`$ and using the second equation. Substituting from the last
$$_iE_c^2(_0\delta A_i_i\delta A_0)=_iE_c^2_i\eta _i(_0\delta A_i_i\delta A_0)$$
into the first equation we obtain
$$\frac{d^2}{dt^2}\eta +\mathrm{}\eta _i(_0\delta A_i_i\delta A_0)=0$$
(78)
The second of the three equations can be written in the form
$$(1+E_c^2)(_0\delta A_i_i\delta A_0)E_c^2_i\eta =(1+E_c^2)C_i(r)$$
(79)
where $`𝐂(r)`$ is an arbitrary function. Dividing eq.(60) by $`(1+E_c^2)`$ and taking the derivative $`_i`$, we obtain
$`_i(_0\delta A_i_i\delta A_0)`$ $`=`$ $`_i{\displaystyle \frac{E_c^2}{1+E_c^2}}_i\eta +_iC_i`$ (80)
$`=`$ $`{\displaystyle \frac{E^2}{1+E_c^2}}\mathrm{}\eta +{\displaystyle \frac{2E_cE_c^{}}{(1+E_c^2)^2}}{\displaystyle \frac{x_i}{r}}_i\eta +_iC_i`$
Replacing on the right hand side $`E_c^2_i\eta `$ by the expression in eq.(60)this becomes
$$_i(_0\delta A_i_i\delta A_0)=\frac{E_c^2}{1+E_c^2}\mathrm{}\eta +\frac{2E_c^{}}{E_c(1+E_c^2)}\frac{x_i}{r}\left[(_0\delta A_i_i\delta A_0)C_i\right]+_iC_i$$
(81)
Choosing as gauge fixing condition the relation
$$\frac{2E_c^{}}{E_c(1+E_c^2)}\frac{x_i}{r}\left[(_0\delta A_i_i\delta A_0)C_i\right]+_iC_i=0$$
one obtains the following fluctuation equation for $`\eta `$
$$(1+E_c^2)\frac{d^2\eta }{dt^2}+\mathrm{}\eta =0$$
(82)
All the relations from (60) to (82) describe perturbations along the string and perpendicular to the brane. Eq. (82) cannot be considered independently of the others as is apparent from the linkage of the fields in the above equations. Thus if one wants to determine the radiation of the string between the brane and the antibrane, one must connect the asymptotic behaviour of the field $`\eta `$ with that of the vector field $`\delta A_\mu `$.
However, an equation like (82)is also obtained if one evaluates the determinant in the Born–Infeld Lagrangian at the BPS background and with an additional time–dependent scalar $`\eta `$, representing the fluctuation field along a new spatial direction (cf. also ref.). In this case this new scalar field in the $`D`$–brane background has no relevance to the string radiation, and we have
$$det(\eta _{\mu \nu }+F_{\mu \nu })|_{BPS,\eta }=\left|\begin{array}{c}\begin{array}{cc}\begin{array}{ccc}\begin{array}{cccc}\begin{array}{ccccc}\begin{array}{cccccc}1& E_1& E_2& E_3& 0& _0\eta \\ E_1& 1& 0& 0& E_1& _1\eta \\ E_2& 0& 1& 0& E_2& _2\eta \\ E_3& 0& 0& 1& E_3& _3y\\ 0& E_1& E_2& E_3& 1& 0\\ _0\eta & _1\eta & _2\eta & _3\eta & 0& 1\end{array}& & & & \end{array}& & & \end{array}& & \end{array}& \end{array}\end{array}\right|$$
(83)
and so
$$det(\eta _{\mu \nu }+F_{\mu \nu })|_{BPS,\eta }=(1+𝐄_{c}^{}{}_{}{}^{2})(_0\eta )^2(_i\eta )^21$$
(84)
Thus the Lagrangian density becomes
$$=1\sqrt{1+(\eta )^2(1+𝐄_{c}^{}{}_{}{}^{2})(_0\eta )^2}$$
(85)
By expanding the square root and retaining only the lowest order terms, we again obtain a fluctuation equation like (65), but this time for $`\eta `$ with no relevance to radiation of the string. This is equivalent to studying the scattering of the scalar $`\eta `$ off a corresponding supergravity background.
## 6 Absorption of scalar in background of $`D3`$ brane
We now consider the equation of small fluctuations, i.e. eq.(82), in more detail. The fluctuation $`\eta (t,𝐱)`$ represents a scalar field that impinges on the brane which reflects part of it and absorbs part of it depending on the energy $`\omega `$ of the field. The absorption results from and takes place into the singularity of the real potential which corresponds to the black hole with zero event horizon in the analogous case of the dilaton–axion system of e.g. ref.. This absorption is a classical phenomenon. We therefore consider the equation
$$\mathrm{}_r\xi +\omega ^2\left[1+\frac{e^2}{r^4}\right]\xi =0$$
(86)
One can argue that the absorption is a consequence of the nonhermiticity of the potential.
The radial part of this equation is with $`\xi =r^1\mathrm{\Psi }Y_{lm}`$ and angular momentum $`l`$
$$\frac{d^2\mathrm{\Psi }}{dr^2}+\left[\frac{l(l+1)}{r^2}+\omega ^2\left(1+\frac{e^2}{r^4}\right)\right]\mathrm{\Psi }=0$$
(87)
This equation is a radial Schrödinger equation for an attractive singular potential $`r^4`$ but depends only on the single coupling parameter $`\kappa =e\omega ^2`$ for constant positive Schrödinger energy, i.e. for $`S`$-waves the equation is with $`x=\omega r`$ simply
$$\left(\frac{d^2}{dx^2}+1+\frac{\kappa ^2}{x^4}\right)\mathrm{\Psi }=0$$
(88)
In the following we consider the general case, i.e. $`l0`$. The simplified case of the singular potential replaced by an effective delta–function potential has been considered in refs. and . The solutions and properties of such equations have been studied in detail in the literature, in both the small– and large–$`\kappa `$ domains and with inclusion of the centrifugal term $`l(l+1)/r^2`$ in eq.(87) for the calculation of Regge trajectories $`l\alpha _n(\omega ^2)`$ , ,,. A recent investigation which attempts to treat arbitrary power singular potentials is ref.. Eq.(87) describes waves above the singular potential well. With the substitutions
$$\mathrm{\Psi }(r)=r^{\frac{1}{2}}\psi (r),r=\sqrt{e}e^z,h^2=e\omega ^2,a=l+\frac{1}{2},$$
(89)
the equation becomes the modified Mathieu equation
$$\frac{d^2\psi }{dz^2}+\left[2h^2\mathrm{cosh}2za^2\right]\psi =0$$
(90)
which has been studied in detail in the literature (though some properties, such as large–h asymptotic expansions of Fourier coefficients, have even now not yet been published). Here we study the S–matrix in the domain of finite values of angular momentum $`l`$ and $`h^20`$, i.e. in the domain of $`h^2`$ large. Relevant solutions and matching conditions for this case have been developed in and . We follow the latter of these references here since this makes full use of the symmetries of the solutions. Moreover we can determine also the Floquet exponent $`\nu `$ which ref. leaves undetermined and only remarks that the notion that this is a known function of (our) $`a^2`$ and $`h^2`$ is “partly a convenient fiction”.
For convenience we set in eq. (90) as in ref.
$$a^2=2h^2+2hq+\frac{\mathrm{}(q,h)}{8}$$
(91)
where $`q`$ is a parameter to be determined as the solution of this equation and $`\mathrm{}/8`$ is the remainder of the large–h asymptotic expansion (91), the various terms of which are determined concurrently with corresponding iteration contributions of the solutions $`\psi `$ of the equation and are known explicitly to many orders . Then setting in eq. (90)
$$\psi (q,h;z)=A(q,h;z)exp[\pm 2hi\mathrm{sinh}z]$$
(92)
we obtain an equation for $`A`$ which can be written
$$\mathrm{cosh}z\frac{dA}{dz}+\frac{1}{2}(\mathrm{sinh}z\pm iq)A=\pm \frac{1}{4hi}\left[\frac{\mathrm{}}{8}A\frac{d^2A}{dz^2}\right]$$
(93)
We let $`A_q(z)`$ be the solution of this equation when the right hand side is replaced by zero (i.e. in the limit $`h\mathrm{}`$). Then one finds easily
$$A_q(z)=\frac{1}{\sqrt{\mathrm{cosh}z}}\left(\frac{1+i\mathrm{sinh}z}{1i\mathrm{sinh}z}\right)^{q/4}\stackrel{z\mathrm{}}{}\sqrt{2}e^{z/2}e^{i\pi q/4}$$
(94)
Correspondingly the various solutions $`\psi `$ are
$`\psi (q,h;z)`$ $`=`$ $`A_q(z)exp[\pm 2hi\mathrm{sinh}z]\stackrel{z\mathrm{}}{}{\displaystyle \frac{exp(\pm ihe^z)}{\sqrt{\mathrm{cosh}z}}}e^{i\pi q/4},`$
$`\psi (q,h;z)`$ $`=`$ $`A_q(z)exp[\pm 2hi\mathrm{sinh}z]\stackrel{z\mathrm{}}{}{\displaystyle \frac{exp(ihe^{|z|})}{\sqrt{\mathrm{cosh}z}}}e^{i\pi q/4}`$ (95)
We make the important observation that given one solution $`\psi (q,h;z)`$ we can obtain the linearly independent one either as $`\psi (q,h;z)`$ or as $`\psi (q,h;z)`$, the expression (91) remaining unchanged. With the solutions as they stand, of course $`\psi (q,h;z)=\psi (q,h;z)`$. Below we require solutions $`He^{(i)}(z),i=1,2,3,4`$, with some specific asymptotic behaviour. We define these in terms of the function
$$Ke(q,h;z):=\frac{exp[i\pi q/4]}{\sqrt{2ih}}A_q(z)exp[2hi\mathrm{sinh}z]k(q,h)\psi (q,h;z)$$
(96)
Since this function differs from a solution $`\psi `$ by a factor $`k(q,h)`$, it is still a solution but not with the symmetry property $`\psi (q,h;z)=\psi (q,h;z)`$. Instead, after performing this cycle of replacements the function picks up a factor, i.e.
$$Ke(q,h;z)=\frac{k(q,h)}{k(q,h)}Ke(q,h;z),\frac{k(q,h)}{k(q,h)}=e^{i\frac{\pi }{2}(q+1)}$$
(97)
in leading order. One can easily show that the quantity $`\mathrm{\Phi }_0`$ of ref. is related to $`q`$ by $`\mathrm{\Phi }_0=iq\pi /2+O(1/h)`$. In Fig. 1 we show the behaviour of $`q`$ as a function of $`h`$. In order to be able to obtain the $`S`$–matrix, we have to match a solution valid at $`z=\mathrm{}`$ to a combination of solutions valid at $`z=\mathrm{}`$. This is achieved with the help of Floquet solutions $`Me_{\pm \nu }(z,h^2)`$. As such, these satisfy the same circuit relation as a solution $`M_{\pm \nu }^{(1)}(z,h^2)`$ of eq.(90) expanded in a series of Bessel functions, i.e. we have the proportionality
$$Me_\nu (z,h^2)=\alpha _\nu M_\nu ^{(1)}(z,h^2),\alpha _\nu (h^2)=Me_\nu (0,h^2)/M_\nu ^{(1)}(0,h^2)$$
(98)
The functions $`Me_{\pm \nu }(z,h^2)`$ are expansions of the modified (hence ‘M’ instead of ‘m’) Mathieu equation in terms of exponentials (hence ‘e’) which are uniformly convergent in any finite domain of $`z`$. For large values of the argument $`2h\mathrm{cosh}z`$ of the Bessel functions of the modified Mathieu function $`M_\nu ^{(1)}(z,h^2)`$ can be reexpressed in terms of Hankel functions. With the dominant terms of these we can obtain the large $`2h\mathrm{cosh}z`$ asymptotic behaviour of the Floquet function $`Me_{\pm \nu }(z,h^2)`$, i.e. for $`|z|\mathrm{}`$
$$Me_{\pm \nu }(z,h^2)exp[\pm i\pi \gamma /2]\frac{\mathrm{cos}(2h\mathrm{cosh}z\nu \pi /2\pi /4)}{\sqrt{2h\mathrm{cosh}z}}$$
(99)
where (with $`Me_\nu (z,h^2)=Me_\nu (z,h^2)`$)
$$exp[i\pi \gamma ]=\frac{\alpha _\nu (h^2)}{\alpha _\nu (h^2)}=M_\nu ^{(1)}(0,h^2)/M_\nu ^{(1)}(0,h^2)$$
(100)
We now define the following set of solutions of eq.(90) by setting
$`He^{(2)}(z,q,h)`$ $`=`$ $`Ke(q,h,z),He^{(1)}(z,q,h)=He^{(2)}(z,q,h),`$
$`He^{(3)}(z,q,h)`$ $`=`$ $`He^{(1)}(z,q,h),He^{(4)}(z,q,h)=He^{(2)}(z,q,h)`$ (101)
The solutions so defined have the following asymptotic behaviour (where $`ϵ(z)=(2h\mathrm{cosh}z)^{1/2}`$):
$`He^{(1)}(z,q,h)`$ $`=`$ $`ϵ(z)exp[ihe^zi{\displaystyle \frac{\pi }{4}}],\mathrm{}z>>0,`$
$`\stackrel{r\mathrm{}}{}`$ $`{\displaystyle \frac{exp[i\omega ri\pi /4]}{\sqrt{\omega r}}},`$
$`He^{(2)}(z,q,h)`$ $`=`$ $`ϵ(z)exp[ihe^z+i{\displaystyle \frac{\pi }{4}}],\mathrm{}z>>0,`$
$`\stackrel{r\mathrm{}}{}`$ $`{\displaystyle \frac{exp[i\omega r+i\pi /4]}{\sqrt{\omega r}}},`$
$`He^{(3)}(z,q,h)`$ $`=`$ $`ϵ(z)exp[ihe^|z|i{\displaystyle \frac{\pi }{4}}],\mathrm{}z<<0,`$
$`He^{(4)}(z,q,h)`$ $`=`$ $`ϵ(z)exp[ihe^|z|+i{\displaystyle \frac{\pi }{4}}],\mathrm{}z<<0,`$ (102)
$`\stackrel{r0}{}`$ $`{\displaystyle \frac{r^{1/2}exp[ie\omega /r+i\frac{\pi }{4}]}{(e\omega )^{1/2}}},`$
For the following reasons we choose the latter, i.e. the solution $`He^{(4)}(z,q,h)`$, as our solution at $`r=0`$. The time–dependent wave function with this asymptotic behaviour is proportional to
$$e^{i\omega t+ie\omega /r+i\pi /4}$$
Fixing the wave front by setting $`\phi =\omega t+e\omega /r+\pi /4=const.`$ and considering the propagation of this wave front, we have
$$r=\frac{e\omega }{\phi +\omega t\pi /4}$$
so that when $`t\mathrm{}:r0`$. This means that the origin of coordinates acts as a sink.
With eq. (99) we therefore equate in the domain $`\mathrm{}z>>0`$:
$`Me_\nu (z,h^2)`$ $`=`$ $`{\displaystyle \frac{i}{2}}exp[i\pi \gamma /2]\left\{exp[i\nu {\displaystyle \frac{\pi }{2}}]He^{(1)}(z,q,h)exp[i\nu {\displaystyle \frac{\pi }{2}}]He^{(2)}(z,q,h)\right\},`$
$`Me_\nu (z,h^2)`$ $`=`$ $`{\displaystyle \frac{i}{2}}exp[i\pi \gamma /2]\{exp[i\nu {\displaystyle \frac{\pi }{2}}]He^{(1)}(z,q,h)`$ (103)
$``$ $`exp[i\nu {\displaystyle \frac{\pi }{2}}]He^{(2)}(z,q,h)\},`$
where the second relation was obtained by changing the sign on $`\nu `$ in the first. Changing the sign of $`z`$ we obtain in the domain $`\mathrm{}z<<0`$:
$`Me_\nu (z,h^2)`$ $`=`$ $`Me_\nu (z,h^2)`$
$`=`$ $`{\displaystyle \frac{i}{2}}exp[i\pi \gamma /2]\left\{exp[i\nu {\displaystyle \frac{\pi }{2}}]He^{(3)}(z,q,h)exp[i\nu {\displaystyle \frac{\pi }{2}}]He^{(4)}(z,q,h)\right\},`$
$`Me_\nu (z,h^2)`$ $`=`$ $`{\displaystyle \frac{i}{2}}exp[i\pi \gamma /2]\{exp[i\nu {\displaystyle \frac{\pi }{2}}]He^{(3)}(z,q,h)`$ (104)
$``$ $`exp[i\nu {\displaystyle \frac{\pi }{2}}]He^{(4)}(z,q,h)\},`$
These relations are now valid over the entire range of $`z`$. Substituting eqs.(104)into eqs.(103) and eliminating $`He^{(3)}`$ we obtain
$$\mathrm{sin}\pi \nu .He^{(4)}(z,q,h)=\mathrm{sin}\pi (\gamma +\nu ).He^{(1)}(z,q,h)\mathrm{sin}\pi \gamma .He^{(2)}(z,q,h)$$
(105)
In a similar way one obtains the relations
$`\mathrm{sin}\pi \nu .He^{(2)}(z,q,h)`$ $`=`$ $`\mathrm{sin}\pi (\gamma +\nu ).He^{(3)}(z,q,h)\mathrm{sin}\pi \gamma .He^{(4)}(z,q,h)`$
$`\mathrm{sin}\pi \nu .He^{(1)}(z,q,h)`$ $`=`$ $`\mathrm{sin}\pi \gamma .He^{(3)}(z,q,h)+\mathrm{sin}\pi (\gamma \nu ).He^{(4)}(z,q,h)`$
¿From eqs.(97) and (101) we see that $`He^{(2)}(z,q,h)`$ is proportional to $`He^{(3)}(z,q,h)`$. ¿From (97) and (LABEL:98) we see that the proportionality factor is given by
$$exp[i\frac{\pi }{2}(q+1)]=\frac{\mathrm{sin}\pi (\gamma +\nu )}{\mathrm{sin}\pi \nu }$$
(107)
¿From eq.(105) we can now deduce the S–matrix $`S_le^{2i\delta _l}`$, where $`\delta _l`$ is the phase shift. The latter is defined by the following large $`r`$ behaviour of the solution chosen at $`r=0`$, which in our case is the solution $`He^{(4)}`$. Thus here the S–matrix is defined by (using (105))
$``$ $`\mathrm{sin}\pi \nu {\displaystyle \frac{r^{1/2}e^{ie\omega /r+i\pi /4}}{(e\omega )^{1/2}}}`$ (108)
$`\stackrel{r\mathrm{}}{=}`$ $`(1)^l{\displaystyle \frac{\mathrm{sin}\pi (\gamma +\nu )e^{i\pi /4}}{\sqrt{\omega r}}}\left[{\displaystyle \frac{\mathrm{sin}\pi \gamma (1)^l}{\mathrm{sin}\pi (\gamma +\nu )}}e^{i\pi /2}e^{i\omega r}(1)^le^{i\omega r}\right]`$
$``$ $`{\displaystyle \frac{e^{i\delta }e^{il\pi /2}}{2i\sqrt{r}}}\left[S_le^{i\omega r}(1)^le^{i\omega r}\right]`$
¿From this we deduce that
$$S_l=\frac{\mathrm{sin}\pi \gamma }{\mathrm{sin}\pi (\gamma +\nu )}e^{i\pi (l+1/2)}=\frac{\mathrm{sin}\pi \gamma }{\mathrm{sin}\pi \nu }e^{i\pi (l\frac{1}{2}q)}$$
(109)
We can see the relation of this high–energy (i.e. large $`|h|`$) expression of the S–matrix to the low–energy expression of ref. by recalling that $`R`$ of the latter is here $`exp(i\pi \gamma )`$. With this identification we can write $`S_l`$
$$S_l=\frac{R\frac{1}{R}}{(Re^{i\pi \nu }\frac{e^{i\pi \nu }}{R})}e^{i\pi (l+1/2)},Re^{i\pi \gamma },$$
which agrees with the $`S`$–matrix of ref., i.e. we thus obtained the same exact expression of the $`S`$–matrix here with our large–$`h`$ considerations. In fact, comparison with the considerations given there allows one to write down the reflection and transmission amplitudes $`A_r`$ and $`A_t`$ as $`A_r=2i\mathrm{sin}\pi \gamma `$ and $`A_t=\mathrm{sin}\pi \nu `$ respectively. We thus have one and the same expression for the $`S`$–matrix for the two asymptotic regions, i.e. in the low energy and high energy domains. One should therefore be able to proceed directly to the large–$`h`$ case from the exact $`S`$–matrix derived in the small–$`h`$ domain. This is an interesting calculation which we do not attempt to go into here. We only indicate in Appendix A the first necessary step in that direction, i.e. the derivation of large–$`h`$ asymptotic expansions for the Fourier coefficients of Mathieu functions. In this connection we make the following two observations. 1) Eq.(88) is invariant under interchanges $`x\kappa /x,\mathrm{\Psi }\mathrm{\Psi }x`$ which means that the inner or string region is equivalent or dual to the outer or brane region. 2) Due to the $`SL(2,R)`$ invariance of the $`D3`$–brane its action is mapped into that of an equivalent $`D3`$–brane by $`S`$–duality transformations or, in other words, weak–strong duality takes the $`D3`$–brane into itself. It would be interesting to find some connection between these properties, or equivalently the symmetry which the $`SL(2,R)`$ invariance of the $`D3`$–brane action imposes on the $`S`$–matrix.
The quantity $`\gamma `$ is now to be determined from eq. (107). One finds
$$\mathrm{sin}\pi \gamma =\mathrm{sin}\pi \nu \left\{ie^{i\frac{\pi }{2}q}\mathrm{cos}\pi \nu \pm \sqrt{1+e^{i\pi q}\mathrm{sin}^2\pi \nu }\right\}$$
(110)
It remains to determine the Floquet exponent $`\nu `$ in terms of $`q`$ and $`h`$. In Appendix B we derive the appropriate large–$`h`$ behaviour of $`\nu `$ for the case of the periodic Mathieu equation. Replacing there the eigenvalue $`\lambda `$ by $`a=(l+\frac{1}{2})^2`$ and observing that $`h^2`$ remains $`h^2`$, the appropriate relation for our considerations is
$`\mathrm{cos}\pi \nu +1`$ $`=`$ $`{\displaystyle \frac{\pi e^{4h}}{(8h)^{q/2}}}\left[{\displaystyle \frac{1+\frac{3(q^2+1)}{64h}}{\mathrm{\Gamma }[\frac{3}{4}\frac{q}{4}]\mathrm{\Gamma }[\frac{1}{4}\frac{q}{4}]}}+O({\displaystyle \frac{1}{h^2}})\right]`$ (111)
$`=`$ $`{\displaystyle \frac{e^{4h}}{(8h)^{q/2}}}\left[{\displaystyle \frac{\left(1+\frac{3(q^2+1)}{64h}\right)\mathrm{\Gamma }(\frac{q+1}{2})\mathrm{cos}(\frac{q\pi }{2})}{\sqrt{2\pi }2^{q/2}}}+O({\displaystyle \frac{1}{h}})\right]`$
Since the right hand side grows exponentially with increasing $`h`$ the Floquet exponent $`\nu `$ must have a large imaginary part. Since the right hand side is real, the real part of $`\nu `$ must be an integer. Using Stirling’s formula we can approximate the equation for $`qh`$ (i.e. irrespective of what the value of $`l`$ is) as
$$\mathrm{cos}\pi \nu +1=\sqrt{\frac{h}{2}}\mathrm{cos}(\frac{h\pi }{2})(e^7/32)^{h/2}\sqrt{\frac{h}{2}}e^{1.8h}\mathrm{cos}(\frac{h\pi }{2})$$
(112)
¿From eq.(109) and eq.(110) we obtain
$$S_l=ie^{il\pi }\left(\mathrm{cos}\pi \nu \sqrt{\mathrm{cos}^2\pi \nu 1e^{iq\pi }}\right)$$
(113)
¿From this we obtain the absorptivity $`A(l,h)`$ of the l–th partial wave, i.e.
$$A(l,h):=1|S_l|^2$$
(114)
with near asymptotic behaviour
$$A(l,h)1\frac{2\pi (16h)^q}{e^{8h}\left\{\mathrm{\Gamma }(\frac{q+1}{2})\right\}^2}$$
(115)
In Figs. 2, 3 and 4 we plot $`A(l,h)`$ as a function of $`h`$. One can clearly see the expected asymptotic approach to unity and in Fig. 2 some sign of rapidly damped oscillations. This behaviour agrees with that obtained on general grounds in ref.. We also observe that in the high energy limit logarithmic contributions as in the low energy expansions, discovered originally in , and typical of the low energy expansions of and , do not arise. Of course, these plots do not extend down to $`h=0`$, since our asymptotic solutions become meaningless in that domain. The continuation to $`h=0`$ can be obtained, however, from small–h expansions such as those derived in refs. and . Thus the absorptivity $`A(l,h)`$ is known over the entire range of $`h`$. We observe that $`S_l=0`$ for $`q=1,3,5,\mathrm{}`$, with $`[(l+1/2)^2+2h^2]/2h1,3,5,\mathrm{}`$. Only in the plot for $`l=2`$ is $`h`$ sufficiently large to hint at these zeros.
## 7 Concluding remarks
Branes, whether fundamental or solitonic, play an important role in all aspects of string theory. In particular $`D`$–branes have been looked at as string–theory analogues of solitons of simple field theories, and some of their important properties such as charges are well understood. Our first objective in the above was to investigate properties of solitonic objects of Born–Infeld theory in ways familiar from field theory, in particular their classical stability. It was shown that the BIon and the catenoid as distinct, i.e. free objects, are stable configurations whereas the brane–antibrane system is unstable; we also recognised the zero modes associated with these and their significance. We then considered the $`D3`$–brane of Born–Infeld theory and recognised this as a BPS state that preserves half of the number of supersymmetries as discussed in detail already in . The equation of small fluctuations about this $`D3`$–brane was derived and shown to be convertible into a modified Mathieu equation. The low energy solutions of this equation, the $`S`$–matrix for scattering of a massless scalar off the brane and the corresponding absorption and reflection amplitudes are similar to those for the dilaton–axion system investigated first in refs., where the important logarithmic contributions were discovered, and then investigated in extensive detail in and . Here we performed the high energy calculations which complement in particular those of , thus completing the investigation of the modified Mathieu equation for the purpose of obtaining absorption cross sections for all such cases. In particular the behaviour of the important Floquet exponent involved in these calculations (in general a complex quantity) is now fully understood, the Floquet exponent being vital in the evaluation of the $`S`$–matrix which we derive and the calculation of the corresponding absorption amplitudes and cross sections. According to our findings the high energy limit of the absorption cross section does not involve logarithmic contributions, quite contrary to the low energy limit.
The high energy case considered here is not only of interest in the immediate context of the Born–Infeld model considered here, but together with the low–energy case also of considerable interest in connection with the concept of duality which links weak coupling with strong coupling. The $`D3`$–brane with Schrödinger potential coupling $`e\omega ^2`$, which links the gauge field charge $`e`$ with energy$`\omega `$ of the incoming scalar field is presumably the ideal example for the investigation of this property. Investigations elucidating this aspect are of considerable interest. We also envisage interest in the study of non–BPS configurations, including sphalerons and bounces, as a matter of principle, i.e. even if the effect of these is not of dominant importance. Finally we remark that it should be possible to proceed directly from the $`S`$–matrix derived in ref. to the high–energy case here by using appropriate asymptotic expansions for the cylindrical functions and expansion coefficients involved (for the latter such expansions do not seem to have been given in the published literature so far, but we comment on these in Appendix A).
Acknowledgements
D.K.P, S.T. and J.-z. Z. are indebted to the Deutsche Forschungsgemeinschaft (Germany) for financial support of visits to Kaiserslautern; the work of J.-z.Z. has also been supported in part by the National Natural Science Foundation of China under Grant No. 19674014 and the Shanghai Education Development Foundation.
Appendix A
In ref. on the absorptivity of the $`D3`$–brane of the dilaton–axion system it was shown that the $`S`$–matrix for scattering of a massless scalar field off the brane is given by
$$S=\frac{(R\frac{1}{R})e^{i\pi l}}{Re^{i\nu \pi }\frac{e^{i\nu \pi }}{R}}$$
(A.1)
where
$$R=\frac{M_\nu ^{(1)}(0,h)}{M_\nu ^{(1)}(0,h)},$$
$`M_\nu ^{(1)}(z,h)`$ being the modified Mathieu function expanded in terms of Bessel functions, i.e.
$$Me_\nu (0,h)M_\nu ^{(1)}(z,h)=\underset{r=\mathrm{}}{\overset{\mathrm{}}{}}c_{2r}^\nu (h^2)J_{\nu +2r}(2h\mathrm{cosh}z)$$
(an expansion with better convergence to use in practice is one in terms of products of Bessel functions as shown in ref.) where $`Me_\nu (z,h)`$ is the Fourier or Floquet solution of the Mathieu equation. In the published literature the coefficients $`c_{2r}^\nu (h^2)`$ have only been considered as power series in rising powers of $`h^2`$, and consequently were used in ref. in the small $`h^2`$ or low energy domain. It would be very interesting to make the transition to the large–$`h^2`$ or high energy case directly from this expression by developing large–$`h^2`$ asymptotic expansions of the Mathieu function Fourier coefficients $`c_{2r}^\nu (h^2)`$ (for the Bessel functions the corresponding expansions are known). We know of no publication where such expansions have been given, but one of us (M.–K.) remembers from private communication with the author of ref. that these Stokes–type asymptotic expansions can indeed be obtained. One writes the recurrence relation of the coefficients (cf. , p.106)
$$c_{2\rho +2}+c_{2\rho 2}=\frac{[\lambda (\nu +2\rho )^2]}{h^2}c_{2\rho }$$
(A.2)
(the Mathieu equation being $`y^{\prime \prime }+(\lambda 2h^2\mathrm{cos}2x)y=0`$). For $`|h^2|\mathrm{}`$ this implies
$$c_{2\rho +2}i^{(2\rho +2)/2}$$
Setting
$$c_{2\rho +2}b_{\rho +1},b_\rho =i^\rho \beta _\rho $$
we have
$`b_{\rho +1}+b_{\rho 1}`$ $`=`$ $`{\displaystyle \frac{[\lambda (\nu +\rho )^2]}{h^2}}b_\rho ,`$
$`\beta _{\rho +1}+\beta _{\rho 1}`$ $`=`$ $`{\displaystyle \frac{i[\lambda (\nu +\rho )^2]}{h^2}}\beta _\rho `$ (A.3)
¿From this we deduce that the next approximation to $`c_{2\rho +2}`$ is obtained from
$$\beta _r=1+\frac{i}{h^2}\underset{\rho =0}{\overset{r}{}}\left[(\rho +\nu )^2\lambda \right]$$
(A.4)
The sums on the right hand side can be evaluated. E.g.
$$\underset{\rho =0}{\overset{r}{}}\rho ^2=1^2+2^2+3^2+\mathrm{}+r^2=\frac{1}{6}r(r+1)(2r+1)$$
so that one obtains
$$\beta _r=1+\frac{i}{h^2}\left[\frac{r(r+1)(2r+1)}{6}+2\nu \frac{r(r1)}{2}+\nu ^2\lambda \right]$$
(A.5)
Proceeding in this way one can indeed obtain the desired asymptotic expansion of the coefficients $`c_{2\rho }`$. (In fact the asymptotic expansion of the Bessel function – similar to that of a linear combination of Hankel functions – can be obtained from its recurrence relation in a very similar way).
Appendix B
For the determination of the large–h behaviour of the Floquet exponent $`\nu `$ we make use of results of ref.. A fundamental pair $`y_I,y_{II}`$ of respectively even and odd solutions of the original periodic Mathieu equation with eigenvalue $`\lambda `$ defined by
$$y_I(z)=y_I(z),y_{II}(z)=y_{II}(z)$$
can be chosen to satisfy the following boundary conditions (cf. e.g. , pp.99,100)
$$y_I(0)=1,y_{II}(0)=0,y_{I}^{}{}_{}{}^{}(0)=0,y_{II}^{}{}_{}{}^{}(0)=1$$
¿From its original defining property the Floquet exponent $`\nu `$ can then be shown to be given by (cf. , p.101)
$$\mathrm{cos}\pi \nu =y_I(\pi ;\lambda ,h^2)$$
(B.1)
so that (cf., p. 100)
$$\mathrm{cos}\pi \nu +1=2y_I(\pi /2;\lambda ,h^2)y_{II}^{}{}_{}{}^{}(\pi /2;\lambda ,h^2)$$
(B.2)
The solutions $`y_I(z),y_{II}(z)`$ can be identified with the large–h solutions $`ce,se`$ of ref.(there eqs.(64))in terms of functions $`A(z),\overline{A}(z)`$ as in eq.(92) above with normalization constants $`N_0,N_0^{}`$, i.e. in leading order
$$ce(o)=2N_0A(0),se^{}(0)=4hN_{0}^{}{}_{}{}^{}A(0)$$
from which we deduce in leading order for large $`|h|`$ that
$$N_0=2^{3/2},N_{0}^{}{}_{}{}^{}=2^{5/2}/h$$
Eqs.(65) of ref. give the large–h expansions of $`y_I(\pi /2;\lambda ,h^2)`$ and $`y_{II}^{}{}_{}{}^{}(\pi /2;\lambda ,h^2)`$. Inserting these multiplied by the appropriate normalization constants into eq.(B.2) and retaining the dominant terms for large $`|h|`$ we obtain
$$\mathrm{cos}\pi \nu +1=\frac{\pi e^{4h}}{(8h)^{q/2}}\left[\frac{1+\frac{3(q^2+1)}{64h}}{\mathrm{\Gamma }[\frac{3}{4}\frac{q}{4}]\mathrm{\Gamma }[\frac{1}{4}\frac{q}{4}]}+O(\frac{1}{h^2})\right]$$
(B.3)
in agreement with a result cited in ref.(p.210) from with logarithmic corrections. We, however, see no such logarithmic terms in the simpler formulation of ref.. The relation (B.3) we rediscovered here has practically been unknown, largely in view of the difficulty to extract it from the complicated considerations of ref.. Our derivation above is simple and closes a difficult gap which the author of ref. commented upon with the words: “It is not likely at this stage that an analytic relation will ever be found connecting (our) $`\nu `$ and $`\gamma `$ to (our) $`a^2`$ and $`h^2`$”. Our search of later literature did not uncover other derivations. The main source summarizing more recent developments in the field of the Mathieu equation is ref..
Figure Captions
Figure 1
The function $`q(h)`$ plotted versus $`h`$, which, of course, is valid only away from $`h=0`$. The plot should be compared with graphs in ref. where a similar but less convenient quantity is used.
Figure 2
The absorptivity $`A(l,h)`$ for $`l=0`$.
Figure 3
The absorptivity $`A(l,h)`$ for $`l=1`$.
Figure 4
The absorptivity $`A(l,h)`$ for $`l=2`$.
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# Renormalization group irreversible functions in more than two dimensions
## Abstract
There are two general irreversibility theorems for the renormalization group in more than two dimensions: the first one is of entropic nature, while the second one, by Forte and Latorre, relies on the properties of the stress-tensor trace, and has been recently questioned by Osborn and Shore. We start by establishing under what assumptions this second theorem can still be valid. Then it is compared with the entropic theorem and shown to be essentially equivalent. However, since the irreversible function of the (corrected) Forte-Latorre theorem is non universal (whereas the relative entropy of the other theorem is universal), it needs the additional step of renormalization. On the other hand, the irreversibility theorem is only guaranteed to be unambiguous if the integral of the stress-tensor trace correlator is finite, which happens for free theories only in dimension smaller than four. PACS numbers: 11.10.Gh, 04.62.+v, 11.10.Kk
The search for a function representing the irreversible nature of the coarse-graining transformations of Wilson’s renormalization group (RG) has a long history. After the success of Zamolodchikov’s $`c`$-function in two dimensions ($`2D`$), it was shown that a straight-forward generalization to higher dimension was not possible but, at the same time, it was observed that a related function, the integral of the stress-tensor trace on a constant curvature space, could play a similar role . In an interesting article Forte and Latorre formulated an irreversibility theorem in terms of this quantity. However, an exhaustive analysis of this theorem carried out by Osborn and Shore shows that there were missing terms in that theorem that actually spoil the irreversible character of that function.
In a separate development, we have introduced in Field Theory the relative entropy, a quantity borrowed from probability theory which turns out to be the Legendre transform of $`W(\lambda )W(0)`$ with respect to $`\lambda `$ :
$`S_{\mathrm{rel}}(\lambda )`$ $`=`$ $`W(\lambda )W(0)\lambda {\displaystyle \frac{dW}{d\lambda }}`$ (1)
$`=`$ $`WW_0\lambda f_\lambda ,`$ (2)
where $`f_\lambda `$ is a composite field integrated over the whole space, $`f_\lambda =d^Dx\mathrm{\Phi }(x)`$. As a straightforward consequence of its definition, the relative entropy satisfies a monotonicity theorem,
$`{\displaystyle \frac{dS_{\mathrm{rel}}}{d\lambda }}`$ $`=`$ $`{\displaystyle \frac{dW}{d\lambda }}{\displaystyle \frac{d}{d\lambda }}\left(\lambda {\displaystyle \frac{dW}{d\lambda }}\right)=\lambda {\displaystyle \frac{d^2W}{d\lambda ^2}}`$ (3)
$`=`$ $`\lambda {\displaystyle \frac{d}{d\lambda }}f_\lambda =\lambda (f_\lambda f_\lambda )^20,`$ (4)
which can be interpreted as showing the irreversibility of the RG . As we remarked in previous papers, the relative entropy is not the only monotonic quantity with the RG. For example, from the same equation that shows its monotonicity (4), one can realize that the function $`f_\lambda =dW/d\lambda `$ is monotonic as well.
Let us consider the integral $`d^Dx\mathrm{\Theta }(x)`$, where $`\mathrm{\Theta }`$ is the stress tensor trace. In a homogeneous space the expectation value $`\mathrm{\Theta }`$ is independent of the position and only depends on the coupling constants; hence, the integration is trivial, its only effect being to produce an overall factor. We further consider a field theory with simple scaling behavior, namely, with only one coupling constant such that $`\lambda m^y`$, where $`m`$ is the physical mass of the fundamental particle or some other mass scale. This behavior is very common in critical phenomena. Since $`\mathrm{\Theta }`$ gives the response to a change of the scale $`m`$,
$$\mathrm{\Theta }m\frac{dW}{dm}=y\lambda \frac{dW}{d\lambda }=y\lambda \mathrm{\Phi },$$
(5)
where now $`W`$ is a specific quantity (per unit volume). In other words, the expectation value of $`\mathrm{\Theta }`$ is proportional to the monotonic function $`f_\lambda `$. By substituting for it in the monotonicity equation (4), we can write this equation in the maybe more suggestive form
$$m\frac{d}{dm}(m^y\mathrm{\Theta })=m^yd^Dx\mathrm{\Theta }(x)\mathrm{\Theta }(0)_c,$$
(6)
where the subscript $`c`$ means that one is to take the connected correlation function. The integral of this correlation function may be divergent. If $`m0`$ it converges for $`x\mathrm{}`$. On the other hand, the behavior of the two-point function for $`x0`$ is the same as in the massless ($`\lambda =0`$) theory, thus given by the dimension of $`\mathrm{\Phi }`$, $`d_\mathrm{\Phi }`$. Therefore, the integral is UV convergent if $`2d_\mathrm{\Phi }<D`$, that is, if $`y=Dd_\mathrm{\Phi }>D/2`$. One can also derive an equation for $`\mathrm{\Theta }`$:
$$m\frac{d}{dm}\mathrm{\Theta }=d^Dx\mathrm{\Theta }(x)\mathrm{\Theta }(0)_cy\mathrm{\Theta }.$$
(7)
In Euclidean space the form of the quantities defined above is given by scaling (e.g., $`S_{\mathrm{rel}}m^D`$) and has little physical content. In a curved homogeneous space one can form the dimensionless variable $`u=Rm`$, where $`R`$ is the curvature radius, and dimensionless quantities are non-trivial functions of it. In particular, we have the dimensionless function of Refs. , $`c(u)=R^D\mathrm{\Theta }`$. Introducing a constant curvature space has an additional utility: Eq. (7) can also be obtained starting from the scale Ward identities satisfied by the energy momentum tensor as $`R`$ varies . Let us remark that the derivation in Ref. yields a slightly different equation. It has been polished in Ref. similar but more general than Eq. (7):
$`R{\displaystyle \frac{d}{dR}}\left(R^D\mathrm{\Theta }\right)`$ $`=`$ $`R^D{\displaystyle d^Dx\mathrm{\Theta }(x)\mathrm{\Theta }(0)_c}`$ (9)
$`R^D\beta ^i(_i𝒜+_i\beta ^j\mathrm{\Phi }_j).`$
This equation takes into account the possibility of several couplings and the existence of the trace anomaly $`𝒜`$, such that $`T_a^a=\mathrm{\Theta }+𝒜`$, where $`\mathrm{\Theta }=\beta ^i\mathrm{\Phi }_i`$. We can convert Eq. (9) into Eq. (7) by (i) assuming simple scaling behavior, that is, with only one coupling such that $`\beta =y\lambda `$, the anomaly $`𝒜`$ being independent of it, and by (ii) replacing the derivative with respect to $`R`$ with a derivative with respect to $`m`$.
Therefore, even though in the general case no monotonicity theorem seems to follow from Eq. (9) , in our case it does, namely, the one expressed by Eq. (6). However, the monotonic quantity (with respect to $`m`$ or $`R`$, indistinctly) is not just $`c(u)=R^D\mathrm{\Theta }`$, as proposed in Refs. , but rather $`\stackrel{~}{c}(u)=u^yc(u)`$. They only coincide if $`y=0`$, that is, when the coupling constant is dimensionless. Generally, the functions $`c`$ or $`\stackrel{~}{c}`$, involving the composite field $`\mathrm{\Phi }`$, contain (normal order) UV divergences. We can introduce a UV regulator but, given that it can only be removed by introducing another scale (renormalization point), those functions are not universal.
To define a finite monotonic function from the stress-tensor trace, one has, therefore, to perform a subtraction. Let us define the function
$$f(u)=V_{D1}m^yR^{Dy}\mathrm{\Theta },$$
(10)
where $`V_{D1}=2\pi ^{D/2}/\mathrm{\Gamma }(D/2)`$ is the volume of the unit $`(D1)`$-dimensional sphere. This function is essentially $`\stackrel{~}{c}(u)`$, except for a conventional sign (to make it increasing rather than decreasing) and a normalization factor. It is UV divergent but, assuming the convergence of the integral in Eq. (6), one subtraction suffices to render it finite. The point is that when integrating $`df/dm`$ according to Eq. (6), one has an integration constant, which can be infinite. Therefore, we can define a renormalized value as
$$f_{\mathrm{ren}}(mR):=\underset{\mathrm{\Lambda }\mathrm{}}{lim}[f_\mathrm{\Lambda }(mR)f_\mathrm{\Lambda }(m_0R)],$$
(11)
$`\mathrm{\Lambda }`$ and $`m_0`$ being the UV cutoff and the subtraction point, respectively. In particular, one can set $`m_0=0`$. Alternatively, one can use minimal subtraction, by which one only subtracts the divergent part of $`f`$, which is independent of $`m`$ . In any renormalization scheme we use the freedom afforded by the integration constant of Eq. (6), for example, to make $`f(0)=0`$, which is equivalent to taking $`m_0=0`$.
In contrast, the relative entropy is universal (under the assumption of convergence of the integral in Eq. (6)) because the UV divergences of $`W`$ cancel in the definition of $`S_{\mathrm{rel}}`$, Eq. (2). One can define a dimensionless growing entropy $`S`$, proportional to $`S_{\mathrm{rel}}`$. In terms of the function $`f`$,
$$S(u)=y_0^u𝑑vv^{y1}f(v)u^yf(u).$$
(12)
The renormalization constant of $`f`$ cancels in this formula.
To illustrate the general theory, we will study a free massive scalar field $`\varphi `$, with coupling constant $`m^2`$, in $`D`$-dimensional hyperbolic space $`H^D`$, for $`D=2,3,4`$. Naturally, a free massive scalar field theory is the simplest example of simple scaling one can take. The field expectation value $`\varphi ^2`$ is then the Gaussian model energy $`U(m^2)`$ , while $`S_{\mathrm{rel}}`$ is a real thermodynamic entropy. Some expressions for the quantities in $`D=2`$ have been calculated in Ref. , in terms of the variable $`r=(Rm)^2`$. More extensive calculations of $`\varphi ^2`$ are given by Osborn and Shore . For $`H^2`$,
$$f_{\mathrm{ren}}(r)=\psi (\sqrt{r+1/4}+1/2)+\gamma ,$$
(13)
where $`\psi `$ is the digamma function and $`\gamma `$ is the Euler constant. $`f`$ increases with $`r`$, on account of the properties of $`\psi `$.
For $`D=3`$ we could use as well the results of Osborn and Shore but it is easier to use instead the heat-kernel method , since the $`D=3`$ heat kernel is extremely simple :
$$K(0;t)=\frac{e^t}{(4\pi t)^{3/2}}.$$
(14)
Hence,
$`f_{\mathrm{ren}}(r)`$ $`=`$ $`4\pi {\displaystyle \underset{0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{dt}{(4\pi t)^{3/2}}}\left[e^{(r+1)t}e^t\right]`$ (15)
$`=`$ $`\sqrt{r+1}1.`$ (16)
This function is obviously increasing.
In $`D=4`$, the case used as example in Ref. , $`f`$ has an expression similar to the one for $`D=2`$ . However, it is a particularly interesting case because the integral in Eq. (6) is now divergent, so $`f_\mathrm{\Lambda }^{}(r)`$ must be subtracted too. Consequently, two subtractions on $`f`$ are needed now, that is,
$$f_{\mathrm{ren}}(r)=\underset{\mathrm{\Lambda }\mathrm{}}{lim}[f_\mathrm{\Lambda }(r)f_\mathrm{\Lambda }(r_0)(rr_0)f_\mathrm{\Lambda }^{}(r_0)].$$
Subtraction at $`m_0=0`$ yields
$`f_{\mathrm{ren}}(r)`$ $`=`$ $`{\displaystyle \frac{1}{4}}[(r+2)\psi (\sqrt{r+9/4}+1/2)`$ (18)
$`2(1\gamma )(39\gamma +\pi ^2){\displaystyle \frac{r}{9}}],`$
which decreases for $`r>0`$. The reason is the following. One can compute $`f_\mathrm{\Lambda }^{}(r)`$ and it is indeed positive for sufficiently large $`\mathrm{\Lambda }`$, since it diverges as $`\mathrm{ln}(\mathrm{\Lambda }^2/r)`$. However, the subtraction removes precisely this dominant growing term. Given that the function $`f^{\prime \prime }(r)`$ is negative (besides finite), $`f_{\mathrm{ren}}^{}(r)<f_{\mathrm{ren}}^{}(r_0)`$ if $`r>r_0`$. This could induce one to try to make the subtraction at the highest $`r_0`$ possible. This might be the idea behind the procedure proposed in Ref. , where it is demanded that $`lim_r\mathrm{}f_{\mathrm{ren}}(r)=0`$. However, this prescription implies subtracting from $`f`$ a function that is not a first degree polynomial in $`r`$, unlike in standard renormalization prescriptions, as exposed here (see also ).
Similar but more complicated expressions are obtained for the positive curvature case, the $`D`$-dimensional sphere $`S^D`$. In this case, one must also consider that for $`r=0`$ the zero mode must be removed from the discrete spectrum, as done for $`D=2`$ in Ref. . This subtraction, however, does not spoil positivity of the second term in Eq. (4).
Let us clarify the role of the trace anomaly, $`𝒜`$. It is well known that renormalization of the free action on a curved even-dimensional spacetime demands the presence of a term proportional to the curvature. It absorbs a logarithmic divergence that appears in addition to the logarithmically divergent term proportional to $`m^2`$ present on the plane . Thus, the logarithmic derivative of $`W`$ with respect to the scale $`R`$ has two components: the stress-tensor trace on the plane $`\mathrm{\Theta }`$ plus an additional part, independent of $`m`$ and proportional to $`R^D`$, the trace or conformal anomaly. The alert reader may have noticed that the original form of the $`R`$-monotonicity theorem (9) in Ref. has $`R^DT_a^a`$ in place of $`R^D\mathrm{\Theta }`$, but it does not matter because the difference is a constant. Nevertheless, adding this constant would have been a convenient normalization for the critical value of the monotonic quantity, had it been precisely $`c(u)=R^DT_a^a`$, as proposed in Refs. : it would make it proportional to the conformal central charge. However, since the correct monotonic quantity is rather $`\stackrel{~}{c}(u)=u^yR^D\mathrm{\Theta }`$, adding the conformal anomaly would result in a divergence at the critical point.
Let us say a few words about the flat space limit $`R\mathrm{}`$. To take this limit, the function $`f`$ is no longer appropriate, and one must instead consider a local quantity, such as $`R^{yD}f(mR)=V_{D1}m^y\mathrm{\Theta }`$. Thus, for the massive free field theory in $`D=3`$, $`lim_R\mathrm{}R^{yD}f(mR)=m`$. In contrast, for $`D=2`$, $`R^0f(mR)=f(mR)`$ diverges logarithmically as $`R\mathrm{}`$, as deduced from the corresponding asymptotic expansion . It is because $`R`$ plays the role of an IR cutoff, and $`f(0)`$ on the plane is IR divergent, as well as UV divergent. The solution is to subtract at $`r_00`$ before taking the limit, which will depend on $`m_0`$ and, therefore, one cannot construct a universal quantity. The same problem exists in $`D=4`$, even though in this case one should not give particular value to the point $`m_0=0`$, as remarked above. Let us note, in passing, that the leading terms of the asymptotic expansion of $`f(u)`$ yield the flat space limit and, furthermore, for even dimension, the sub-leading term yields the conformal anomaly .
In conclusion, the monotonicity theorems for the relative entropy or for the stress-tensor trace are contained in Eqs. (4). In field theory, $`(f_\lambda f_\lambda )^2`$ is proportional to the integral of the stress-tensor trace correlation, which only converges if $`y`$, the dimension of the coupling constant $`\lambda `$, satisfies $`y>D/2`$. Therefore, only under this condition is the irreversibility theorem unambiguous. However, even in this case, the function $`f`$ associated to the stress-tensor trace is ambiguous (non-universal), being defined only up to a constant, whereas the relative entropy is unambiguous (universal). After renormalization, the ambiguity of $`f`$ is realized as a dependence on $`m_0`$, which is the renormalization point in the simple scheme used here. Setting $`m_0=0`$ achieves a kind of universality, in the sense that no additional scale remains, but it may not be realizable, as occurs for free field theory on the plane. The case $`y>D/2`$ covers many of the critical models of Statistical Mechanics, e.g., the $`3D`$ Ising model universality class, with $`y=1.59`$ . When $`yD/2`$ (in particular, for bosonic free-field theory in $`D=4`$) the integral in the right-hand side of Eq. (6) is UV divergent and must be renormalized, in general spoiling its positivity, so the irreversibility theorem is itself ambiguous and may only hold in a particular renormalization scheme. Accordingly, $`f`$ needs to be subtracted twice. Of course, the problem of the divergence of that integral also affects the relative entropy, which becomes non universal, requiring one additional subtraction further to those implied in its definition (2). Hence, it is doubtful whether one can assign an unambiguous meaning to RG irreversibility for $`yD/2`$.
Since irreversibility in terms of the stress-tensor trace or in terms of the relative entropy are essentially equivalent, one may wonder which formulation is better. From a physical point of view, the theorem for the relative entropy has more content, being related to important notions in Information Theory , while from a mathematical point of view, $`\mathrm{\Theta }`$ is simpler to calculate and, in fact, to calculate $`S_{\mathrm{rel}}`$ one must calculate it before (as in Eq. (12)). I thank Hugh Osborn for a conversation and for patiently explaining to me tricky points on some calculations in Ref.
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# On the Nature of Andromeda IV1footnote 11footnote 1Based on observations with the NASA/ESA Hubble Space Telescope, obtained at the Space Telescope Science Institute, which is operated by the Association of Universities for Research in Astronomy, Inc. under NASA contract No. NAS5-26555. ,2footnote 22footnote 2Based on observations made with the William Herschel and Isaac Newton telescopes operated on the island of La Palma by the Isaac Newton Group in the Spanish Observatorio del Roque de los Muchachos of the Instituto de Astrofisica de Canarias.
## 1 Introduction
Andromeda IV is an enigmatic object first discovered by van den Bergh (1972) during his photographic search for dwarf spheroidal companions to M31. On deep optical plates, it appears as a faint smudge lying at a projected distance of 40 or $``$9 kpc (assuming D$`{}_{M31}{}^{}=`$784 kpc, Holland (1998); Stanek & Garnavich (1998)) to the south-west of M31 (see Figure 1). Noted immediately as being more compact, bluer and of higher surface brightness than the other dwarf spheroidals he identified, van den Bergh suggested And IV was either a relatively old $`\mathrm{`}`$star cloud’ in the outer disk of M31<sup>3</sup><sup>3</sup>3And IV actually appears listed as the open cluster C188 in the Atlas of the Andromeda Galaxy (Hodge, 1981). or a background dwarf galaxy. If in the disk of M31, And IV would have a size of $``$200 pc and lie at a deprojected distance from the center of 25 kpc ($``$5 disk scalelengths, beyond the optical edge of the stellar disk), thus representing an example of a large diffuse star cluster which has formed in the outer regions of a galaxy where there is very little star formation at the present epoch. Such an object could potentially yield important insight into the nature of the star formation process under the extreme physical conditions of low gas surface density, high gas fraction and long dynamical timescales. On the other hand, And IV could be yet another example of a previously uncatalogued dwarf galaxy lurking in our local environs, perhaps even bound to M31 (eg. Armandroff, Davies & Jacoby (1998); Armandroff, Jacoby & Davies (1999)). Determining the distance, constituent stellar populations and evolutionary state of this puzzling object is therefore of obvious importance.
Ground-based study of And IV has been severely hampered due to the presence of a bright (V$``$10) foreground star lying within 30<sup>′′</sup> (see Figure 1) and by the combination of faintness and crowding of M31 stars along the line of sight. Jones (1993) used CFHT imagery to construct a colour-magnitude-diagram (CMD) to V$``$23 of the And IV region, which he interpreted as representing a young population of stars with a narrow age range, an $`\mathrm{`}`$unusually large’ open cluster in M31. Jones’s data, however, did not allow a proper statistical subtraction of M31 field stars to be carried out and thus the possibility remains of a significant foreground M31 contamination in his And IV CMD. The reported HI detection of And IV is also rather uncertain. In their catalogue of extragalactic HI observations, Huchtmeier & Richter (1989) list an HI radial velocity for And IV of $``$375 km/s. This measurement can be traced back to the early work of Emerson (1974), who used the Cambridge Half-Mile telescope to map the disk of M31 with a spatial resolution of 1.5′ $`\times `$ 2.2′. Emerson (1974) notes that And IV appears projected near a faint outer HI arm of M31 (see also Unwin (1980)) and that the measured velocity of the HI along this arm ($``$375 km/s ) is consistent with that expected at that location in the disk based on the major axis rotation curve. The case for the association of this gas with And IV, as opposed to merely M31’s disk, would appear to have no stronger foundation than this, and therefore must be regarded as weak.
Detailed study of the structure and stellar populations of And IV requires the high resolution imaging capability of Hubble Space Telescope. We were awarded 5 orbits of Cycle 6 HST/WFPC2 time to observe And IV and the surrounding M31 field. We present these data here, along with supporting ground based observations (H$`\alpha `$ imagery and optical long-slit spectroscopy). This new dataset clearly resolves the true nature of And IV, for the first time, as a small background dwarf irregular galaxy unassociated with M31. A companion paper will present an analysis of the field stellar populations in the outskirts of M31, derived from the same HST/WFPC2 dataset (Ferguson, et al 2000, in preparation).
## 2 Observations and Reductions
### 2.1 HST/WFPC2 Observations
HST/WFPC2 images were taken of And IV and the surrounding M31 field over five orbits on October 31, 1998 (GO #6734). The 3 WF cameras provide an L-shaped FOV of 150″ by 150″ with 0.1″ pixels, while the PC provides a square 34″ by 34″ FOV with 0.046″ pixels. Total exposure times amounted to 6000s and 6100s in the F555W (WFPC2 broadband V) and F814W (WFPC2 broadband I) filters respectively. The proximity of a very bright star ($``$10th mag) to And IV imposed stringent roll angle constraints on our observations in order to avoid excessive scattered light and bright-object artifacts in our primary region of interest. The optimal orientation of the camera placed And IV largely on the WF3 chip and the bright star on the edge of WF4, a strategy which rendered the WF4 chip essentially useless but kept the region around And IV free of diffraction spikes and other artifacts. The WF3-FIX aperture was centered on $`\alpha _{2000}=`$00<sup>h</sup>42<sup>m</sup>30.1<sup>s</sup> and $`\delta _{2000}=+`$4034′32.7″. Both the WF2 and the PC imaged the surrounding field populations of M31.
Images were processed through the standard STScI pipeline. The frames were split into images of each individual CCD and the vignetted regions of the chips set to zero. The data were taken undithered and measurements of several bright stars confirmed no shifts were required to align the individual frames before combining. Images taken through a given filter were combined using the cosmic-ray rejection algorithm CRREJ in IRAF <sup>4</sup><sup>4</sup>4IRAF is distributed by the National Optical Astronomy Observatories, which are operated by the Association of Universities for Research in Astronomy, Inc., under cooperative agreement with the National Science Foundation.. Corrections were applied for warm pixels using pixel lists provided by STScI that were generated closest to the date of the observations. Figure 3 shows a colour representation of the combined F555W and F814W WF3 images; AndIV is clearly visible as a diffuse blue concentration towards the lower left of the chip.
### 2.2 HST/WFPC2 Photometry
Photometry was performed using the IRAF implementation of the crowded-field photometry package DAOPHOT/ALLSTAR (Stetson, 1987). Stars were detected on each chip using a DAOFIND threshold of 5$`\sigma _{bkgd}`$, a value found to maximise the number of real stars detected while maintaining a low level of spurious sources. Given that our field is relatively crowded ($``$7000$``$8000 stars per WF chip), there were few truly isolated bright stars suitable for characterizing the point-spread function (PSF). Our approach was instead to select a set of $``$30 of the brightest stars on each chip, subtract out all the other stars using a first-guess PSF and then use the remaining $`\mathrm{`}`$clean’ stars to construct a more refined PSF model. This was done for each filter/chip combination. The PSF models were based on Moffat functions with $`\beta =`$1.5 and a look-up table of residuals, and these shapes were held constant across each chip. ALLSTAR photometry was carried out using the parameter values recommended by Cool & King (1995) to maximize the performance. We also experimented with using Tiny Tim PSFs for the photometry instead of PSFs built from the images themselves, and found that the results were generally very similar.
Aperture corrections were derived in the same manner to that used to construct PSF models, namely by isolating the brightest stars on each chip and subtracting out the rest. Total magnitudes were measured for the remaining stars using an aperture of radius 0.5″ (Holtzman et al., 1995b) and aperture corrections were defined in the sense $`m_{PSF}m_{0.5\mathrm{}}`$. As the focus and PSF shape is a function of position on the WFPC2, aperture corrections were allowed to vary linearly across each chip. Corrections were also applied for geometric distortion and charge transfer efficiency using the recommendations of Holtzman et al. (1995a) and Whitmore, Heyer & Casertano (1999) respectively. Before transforming the instrumental magnitudes to the standard system, we corrected our data for the effects of extinction. We adopted a foreground Galactic reddening of E(B$``$V)$`=`$0.08$`\pm `$0.02 towards M31<sup>5</sup><sup>5</sup>5Note that Schlegel, Finkbeiner & Davis (1998) quote the slightly lower value of E(B$``$V)$`=`$0.062 for the Galactic reddening towards M31, derived from the mean dust emission detected by DIRBE in surrounding annuli.. (Burstein & Heiles, 1984) and used Table 12(b) of Holtzman et al. (1995b) to derive the corresponding extinctions in the HST bandpasses. As And IV lies well beyond the main HI disk of M31, in a region where the mean N$`{}_{HI}{}^{}`$3$`\times `$10<sup>20</sup> cm<sup>-2</sup>, (Emerson, 1974; Unwin, 1980; Sofue & Kato, 1981), the reddening due to M31’s disk is low for all reasonable gas-to-dust ratios and hence we apply no correction. As And IV appears projected on a faint outer HI arm, there does exist the possibility of differential reddening across the face of the object (Emerson, 1974; Unwin, 1980). Unfortunately, we have no way to estimate the magnitude of this effect but it should be kept in mind. Transformation to standard Johnson-Cousins V- and I-band magnitudes was carried out via the iterative approach described in Holtzman et al. (1995b), and adopting the values presented in their Table 7.
Our photometry reaches to V$``$27.5, I$``$26.5. Typical 1-$`\sigma `$ photometric errors for stars at 25th magnitude are $`\sigma _V`$0.06 and $`\sigma _I`$0.08. Only stars for which ALLSTAR photometry was deemed high quality ($`\chi <`$2 and $``$0.2$`<`$sharpness$`<`$0.2) are used in our final analysis. The completeness level at faint magnitudes is not a serious concern for the analysis that we present here, and we postpone a detailed discussion of this issue to a future paper. Past experience with WFPC2 photometry leads us to expect $``$80% completeness above V$``$26 and I$``$25.5 (Dohm-Palmer et al., 1997; Cole et al., 1999). Our analysis in this paper focuses only on those stars detected in the WF3 chip.
### 2.3 Ground-based Follow-Up Observations
We also obtained complementary ground-based imaging and spectroscopy of the And IV region. Deep narrow-band observations were obtained via the ING service program in August 1999 using the INT 2.5m telescope equipped with the Wide Field Camera at the f/3.3 prime focus. The Wide Field Camera consists of 4 thinned EEV 4k$`\times `$2k CCDs, each covering an area of $``$23′$`\times `$11′ on the sky with 0.33″ pixels. Our images therefore include a large portion of the surrounding M31 field. Exposures of 3$`\times `$800s were taken through a narrow-band H$`\alpha `$ filter ($`\lambda _c=`$6568Å, $`\mathrm{\Delta }\lambda =`$95Å) and a 300s exposure through broadband Sloan r$``$ for the continuum subtraction. Conditions were photometric, but poor seeing prevailed ($``$2$``$2.5″). The images were reduced in the standard manner (see for example, Ferguson, Wyse, Gallagher & Hunter (1996)) and observations of the spectrophotometric standard Feige 110 from the list of Massey, Strobel, Barnes & Anderson (1988) were used for the photometric calibration. The average sensitivity of the H$`\alpha `$ continuum-subtracted image, taken to be 1$`\sigma `$ of the sky background, is determined to be $``$2.7$`\times `$10<sup>-17</sup> erg s<sup>-1</sup> cm$`{}_{}{}^{2}/\mathrm{}\mathrm{}`$. Figure 3 shows portions of both the unsubtracted and the continuum-subtracted H$`\alpha `$ images centered on And IV, with the WFPC2 field-of-view overlaid on the continuum-subtracted image. Several faint emission line sources are clearly visible in the vicinity of And IV.
We obtained long-slit optical spectroscopy of several of these objects in September 1999 using the WHT 4.2m and the ISIS double-beam spectrograph. The R300B grating was used with a 4k$`\times `$2k thinned EEV CCD in the blue arm to cover the range $``$3700Å$``$5300Å with 3.4Å resolution; in the red arm, we used the R316R grating and a 1K$`\times `$1K TEK CCD to cover $``$5700Å$``$7200Å with 3Å resolution. The slit length and width were 3.7′ and 1″ respectively. Conditions were excellent with sub-arcsecond seeing. Most observations were made at very low airmass to minimize the effects of differential atmospheric refraction. In addition, the slit was rotated close to the parallactic angle whenever the airmass exceeded 1.2. Total exposure times ranged from 1800$``$6300s, depending on the brightness of the source. Spectra of CuAr and CuNe lamps were made during the night to provide a wavelength calibration and observations of spectrophotometric standards from the list of Massey, Strobel, Barnes & Anderson (1988) were made for the flux calibration.
Given the faintness of the emission line sources, special care was required during the processing of the spectra. Individual 2-dimensional spectra from a given slit position were aligned to the nearest integer pixel and the sky removed by fitting a low order polynomial fit with the IRAF BACKGROUND task. The spectra were subsequently averaged together using a cosmic-ray rejection algorithm and apertures ranging from 3-4″ were used to extract one-dimensional spectra of each emission-line object. Care was taken to ensure that red and blue spectra for a given object corresponded to the same physical apertures (both size and shape). A region of $``$20″ centered on the brightest continuum emission from And IV was also extracted. The one-dimensional spectra were wavelength and flux calibrated with residuals of $``$0.5Å and 0.02-0.04 mag respectively. In Figure 4, we show the calibrated spectra for the three brightest nebulae (#3,4,6) in the vicinity of And IV as well as for the underlying continuum emission.
## 3 Results
### 3.1 Stellar Populations
#### 3.1.1 Resolved Emission
The (V,V$``$I) and (I,V$``$I) colour-magnitude diagrams (CMDs) of 6388 stars detected in the WF3 chip are shown in Figure 5. These CMDs reveal a prominent red giant branch (RGB) with a significant intrinsic width, a red clump, a weak blue plume and a possible blue horizontal branch. The striking downturn of the RGB at red colours (V$``$I$``$2) has previously been seen in the outskirts of M31 (eg. Holland, Fahlman & Richer (1996)) and indicates the presence of a metal-rich component to the stellar population. Based on standard star count models (Gilmore, 1981; Wyse & Gilmore, 1989), we expect $``$20 stars redder than the RGB and $``$1 star bluer than the RGB between 21$``$V$``$27 along this line of sight towards M31. Galactic field stars are therefore expected to contaminate our field by a negligible amount ($``$0.3%).
It is difficult to tell from visual inspection of Figure 3 alone what fraction of And IV has resolved on our deep images. In order to isolate the signature of And IV stars from those of the M31 disk, we constructed CMDs for stars lying within a box of 40$`{}_{}{}^{\prime \prime }\times `$60<sup>′′</sup> centered on the brightest emission from And IV and for stars lying on the remaining area of WF3, which we will refer to as the $`\mathrm{`}`$M31 field’ (see Figure 6). This physical region was selected as it completely encompasses the area of $`\mathrm{`}`$diffuse’ emission which defines And IV in Figures 1 and 2. Taking account of the slighly differing areas (0.9:1), we find that the number of stars detected in each region of the chip is largely consistent with a uniform distribution of stars, and indicates only a $``$10% enhancement in star counts (at least to the limits of our photometry, V$``$27.5) in the neighbourhood of And IV. This is a strong indication that only a fraction of And IV has resolved.
Comparison of the morphology of the CMDs in the different areas of the chip reveals some puzzling differences however. The boxes in Figure 6 indicate regions of the CMDs that are populated by stars in the vicinity of And IV but not in the M31 field. There are, for example, considerably more faint stars redward of the RGB (V$``$I$``$1.3, V$``$25) on the And IV CMD than on that of the M31 field (159 stars versus 44). Likewise, the region in between the blue plume and the RGB (V$``$25, 0$``$V$``$I$``$0.7) is also more populated on the And IV CMD (45 versus 7 stars) as is the region just above the RGB (V$``$2.5, 0.8$``$V$``$I$``$2.1, 9 versus 1 star). These overdensities are significant, especially when account is made for the slightly differing physical areas that the stars on each CMD are drawn from.
We first investigate whether the stars detected in the vicinity of And IV simply have larger photometric errors than those in the surrounding M31 field, causing them to exhibit broader blue plumes and red giant branches and scatter to both brighter and fainter magnitudes. The increased and more variable sky background in the And IV region of the chip could possibly lead to this effect, as could increased crowding. To test for this, we calculated the mean photometric errors returned by ALLSTAR in 0.5 magnitude bins for stars lying on and off the And IV area. While the magnitude errors are very similar over most of the magnitude range, they start to diverge towards faint magnitudes. Still, the effect is small. At V$``$27, photometric errors for stars near And IV differ by only $``$0.05 mag from stars located elsewhere on WF3. Thus, while increased photometric scatter can partially explain some of the faint stars redward of the RGB, it cannot explain all of them. Furthermore, increased photometric uncertainties seem an unlikely explanation for the excess populations of stars seen in other regions of the CMD.
We then address the issue of whether the $`\mathrm{`}`$excess’ populations of stars seen in the And IV CMD could be the signature of a distinct, partially resolved stellar population lying at a significant distance beyond M31? Indeed, the stars lying within the marked boxes of the CMD in Figure 6 account for roughly half the observed stellar overdensity seen towards And IV. These stars could appear offset from the M31 main sequence and RGB due to internal extinction within either the galaxy itself, or due to variable small-scale extinction in the foreground disk of M31. We constructed RGB and main-sequence luminosity functions (LFs) in 0.5 mag bins for stars lying within the box centered on And IV and those lying elsewhere on the chip (Figure 7). We crudely define RGB stars as those with V$``$I$`>`$0.6 and main-sequence stars as those with V$``$I$`<`$0.6 (this definition will also allow the inclusion of blue horizontal branch stars into the $`\mathrm{`}`$main sequence’ sample) and normalise the counts in each region for their slightly different areas. Comparison of the main-sequence LFs indicates a genuine excess of blue stars in the vicinity of And IV with respect to the M31 field over the entire magnitude range probed. On the other hand, inspection of the bottom panel of Figure 7 reveals no obvious excess of red stars towards And IV, but does show an intriguing $`\mathrm{`}`$bump’ on the tail of the And IV RGB LF beyond the red clump, at I $``$ 25$``$25.5. This feature also appears in the V-band RGB LF at V $``$ 26$``$26.5, but it is *not* present in the M31 field LF. Quantitative study of this bump is difficult. In this magnitude range, incompleteness is an important factor. In fact, it is probably more of a serious issue in the immediate area of And IV, where the background is higher, than it is elsewhere on the chip; the fact that the bump remains prominent in both passbands suggests it is real. The colour of the feature, V$``$$``$ 1, is suggestive of a population of red giants and we are tempted to speculate that the bump represents a detection of the tip of the red giant branch, or possibly a slightly more luminous extended asymptotic giant branch, in a stellar system located at some distance behind M31. We will return to this issue in more detail in Section 4.
From the analysis of number counts, CMDs and LFs for stars detected in our WFPC2 images, we therefore conclude that there is evidence for only a small enhancement in stellar density towards And IV. Furthermore, a significant fraction of this excess population appears displaced from the main sequence and red giant branch of M31 field stars. While this may imply that And IV is not a star cluster lying in M31’s disk, we are not able to rule out the possibility of a highly-skewed mass function which would allow And IV to still be associated with M31 but not to resolve to the same extent. We therefore turn our attention to other aspects of the WFPC2 data, as well as to other data, in order to derive additional constraints on the nature of the object.
#### 3.1.2 Unresolved Emission
One of the outputs of the ALLSTAR PSF-fitting photometry package is a residual image where all detected and photometered stars have been fit with the adopted PSF and subtracted out. Figure 8 shows the residual WF3 F555W image, which has been smoothed with a box of 10 pixels ($``$ 1<sup>′′</sup>) to remove the residuals of subtracted stars and to increase the signal-to-noise of the faintest emission. A large fraction of unresolved light very clearly remains around And IV; the light distribution appears to be somewhat centrally-concentrated with a regular structure. By comparing the emission remaining in the residual image with that in the original image, we find that only 40% of the F555W light in the And IV region of WF3 has resolved, as compared with 70-80% of the light elsewhere on the chip. For the F814W image, these numbers are 50% and 80$``$90% respectively. This confirms our earlier conclusion that only a small fraction of And IV has resolved compared to the M31 field. From the smoothed image, we are able to define a more accurate centre for And IV, which we report here as $`\alpha _{2000}=`$00<sup>h</sup>42<sup>m</sup>32.3<sup>s</sup> and $`\delta _{2000}=+`$4034′18.7″.
The unresolved emission of And IV is of moderately low surface brightness and very blue colour. Using the residual F555W and F814W images, we measure $`\overline{\mu _V}`$ 24.0, $`\overline{\mu _I}`$ 23.4 and $`\overline{VI}`$ 0.64 within a 30$`{}_{}{}^{\prime \prime }\times `$34<sup>′′</sup> box centered on And IV. We have constructed crude surface brightness profiles for And IV using the 10$`\times `$10 pixel boxcar-smoothed images. Elliptical aperture photometry was carried out using a fixed position angle and ellipticity, both of which were determined by eye to best match the poorly-defined And IV isophotes. The upper panel of Figure 9 shows the V and I-band surface brightness profiles derived out to a radius of 20″ (the extent of the bright residual emission on our WFPC2 image) from the center of And IV. The profiles are observed to decline slowly and smoothly with increasing galactocentric radius. A model of an exponential disk with $`\mu _0=`$23.3 and $`\alpha ^1=`$11″ is overplotted; such a light profile appears to provide a good match to the V-band light of And IV. The lower panel of Figure 9 shows the V$``$I radial colour gradient and indicates a gradual trend of bluer colours towards larger radii, although the S/N ratio becomes very low in these regions.
Surface brightness profiles of open clusters are well-fit by King profiles (eg. Mathieu (1984)), reflecting their tidal limitations. On the other hand, small dwarf irregular galaxies are typically characterised by approximately exponential profiles and small colour gradients (eg. Bremnes, Bingelli & Prugniel (1999)). It thus seems that the exponential profile we find here is further evidence against And IV being an open star cluster in M31; instead, it may support the idea that it is a background dwarf galaxy.
### 3.2 Ionized Gas
Our deep H$`\alpha `$ images reveal eight compact emission-line sources in the general vicinity of And IV, five of which lie within 1′ of our new adopted centre (Figure 3). This clustering of emission-line sources represents a significant overdensity compared to the rest of the field contained in our wide-field images, and suggests a connection between at least some of the sources and the faint diffuse continuum emission which defines And IV. Table 1 lists the positions and fluxes of the H$`\alpha `$ sources, their projected distances from the adopted centre of And IV and indicates whether the source appears to be resolved on our images and whether it has a continuum counterpart. The H$`\alpha `$ fluxes have been measured within a circular aperture of radius 5″, and have been corrected for both Galactic extinction using Burstein & Heiles (1984) and \[NII\] contamination using the mean \[NII\]/H$`\alpha `$ ratio from our spectra (see below). The H$`\alpha `$ fluxes range from 1$``$9$`\times `$10<sup>-15</sup> erg s<sup>-1</sup> cm<sup>-2</sup>; if at the distance of M31, these sources would have only modest H$`\alpha `$ luminosities of 9$``$68$`\times `$10<sup>34</sup> erg s<sup>-1</sup> which, for reference, are 10$``$100 times fainter than that of the Orion nebula (Kennicutt, 1984). Assuming they are ionization-bounded, such objects would be powered by Lyman continuum luminosities, Q<sub>0</sub>, of $``$7$``$50$`\times `$10<sup>46</sup> photons s<sup>-1</sup> (Leitherer & Heckman, 1995), which could possibly be provided by either early B stars<sup>6</sup><sup>6</sup>6Unfortunately, predictions of ionizing fluxes for stars of spectral types later than $``$B0 do not yet exist, and the few direct observations of the Lyman continuum spectral region in B stars have produced very ambiguous results (see discussion in Schaerer & de Koter (1997)). The latest spectral type considered by Schaerer & de Koter (1997) is a B0.5 V star, which they predict produces an ionizing flux of log Q<sub>o</sub>(s$`{}_{}{}^{1})`$ 47.8. or the luminous central stars of planetary nebulae (Schaerer & de Koter, 1997; Vacca, Garmany & Shull, 1996; Mendez & Soffner, 1997). A concentration of PNe within such a small area of the sky seems rather unlikely however.
If we assume ionization by a single massive star in M31, the *most* luminous star we could expect to find associated with each nebula would be of type B0.5V with a magnitude of V$``$20 at that distance (Vacca, Garmany & Shull, 1996; Schaerer & de Koter, 1997). We would not expect to find stars much fainter than this, due to the rapid drop-off in ionizing flux as a function of spectral type. Four of the emission-line sources around And IV (#3,4,5,6) lie in the area of the sky covered by our WFPC2 image and so we were able to search our images for luminous and/or very blue stars and compact clusters in the vicinity of each nebula. Objects #3, 4 and 6 are easily recognizable on Figure 3 as being very blue, and in the case of #3 and 4, also extended. As the F555W filter contains the emission lines \[OIII\]$`\lambda \lambda `$4959,5007 and H$`\beta `$, with even a small transmission at H$`\alpha `$, the nebular morphologies seen at these positions is not surprising. We base our search on the results of aperture photometry, as opposed to PSF-fitting photometry, since the possibility exists that the ionizing sources may be clusters that would not be well-fit by the PSF model; indeed such objects have been deliberately excluded from our final stellar photometry lists (see Section 2.2). The brightest blue sources identified lying within 2″ of the position of each nebula are found to have magnitudes V$``$22.4, 22.6 (Object #3), 23.0 (Object #4) and 24.0 (Object #6), all with V$``$$``$0.1. At the distance of M31, these sources would have spectral types ranging from mid to late-B, and as such, would have difficulty in producing the required ionizing fluxes. There are no very blue stars lying within 2″ of Object #5 however there is a luminous red star with V$``$22.4 and V$``$I$``$1.6. Interestingly, this object lies in one of the areas of $`\mathrm{`}`$excess’ stars identified in Figure 6. Given the considerable stellar density in our WFPC2 field, the likely uncertainties in our absolute positions of the emission line sources and the possibility of differential reddening across the face of And IV, it is difficult to draw firm conclusions regarding the properties of the individual stars which are responsible for the ionization of the nebulae. It would appear safe to conclude however, that unless the ionizing stars are highly obscured, they do not lie at the distance of M31. This provides another piece of evidence for And IV lying a significant distance beyond M31.
### 3.3 Gas-Phase Metallicities
Optical spectra of four of the emission-line sources reveal those in the immediate vicinity of And IV (#3,4,5,6) to be high-excitation HII regions, displaying the usual bright \[OII\] and \[OIII\] lines and, in some cases, even marginal detections of the faint temperature-sensitive \[OIII\]4363Å line (see Figure 4). Objects #7,8, which lie further away from the centre of And IV, appear to be a possible symbiotic nova and a high excitation planetary nebula respectively. No spectra were obtained for Objects #1 and 2.
Emission line fluxes were measured via Gaussian fits to the line profiles. The logarithmic extinction at H$`\beta `$, C(H$`\beta `$), was derived from measurements of the Balmer lines, using the equation
$$\frac{\mathrm{I}_\lambda }{\mathrm{I}_{\mathrm{H}\beta }}=\frac{F_\lambda }{F_{H\beta }}10^{C(H\beta )f(\lambda )}$$
where I<sub>λ</sub> is the intrinsic line flux, F<sub>λ</sub> is the observed line flux, and f($`\lambda )`$ is the Galactic reddening function normalized to H$`\beta `$. The reddening function of Seaton (1979), was adopted, as parametrized by Howarth (1983), and assuming R=A<sub>V</sub>/E(B$``$V)=3.1. Intrinsic case B Balmer line ratios were taken from Osterbrock (1989), assuming an electron density of N<sub>e</sub>=100 cm<sup>-3</sup> and an electron temperature T<sub>e</sub>=10<sup>4</sup> K. The values of C(H$`\beta `$) derived from the H$`\alpha `$/H$`\beta `$, H$`\gamma `$/H$`\beta `$ and H$`\delta `$/H$`\beta `$ ratios agreed within the formal errors and no trend was apparent to indicate the presence of Balmer absorption in the underlying continuum. Given its higher accuracy, we adopt the value of C(H$`\beta `$) determined from the H$`\alpha `$/H$`\beta `$ ratio in our analysis. The derived values of the logarithmic extinction translate into E(B$``$V) $``$ 0.0$``$0.11, and indicate that the line-of-sight extinction towards And IV is very low and almost entirely Galactic.
Formal errors in the derived line ratios were determined by summing in quadrature the statistical noise from the photon counts, the uncertainty in the continuum placement (proportional to the width of the line times the rms in the nearby continuum) and the uncertainty in the flux calibration. In addition, the error in C(H$`\beta `$) was accounted for in deriving the reddening-corrected line ratios. Table 2 presents the reddening-corrected line strengths (relative to H$`\beta `$) for the three brightest HII regions as well as some relevant line ratios. Formal errors are indicated in parentheses.
Oxygen and nitrogen abundances for the HII regions were derived using the well-established $`\mathrm{`}`$semi-empirical’ abundance calibrations proposed by McGaugh (1991, 1994) and Thurston, Edmunds & Henry (1996) via the procedures described in Ferguson, Gallagher & Wyse (1998). These calibrations are based on the strengths of the bright oxygen, nitrogen and Balmer lines of hydrogen via the parameter
$$\mathrm{R}_{23}=\frac{\left[OII\right]\lambda 3727+\left[OIII\right]\lambda \lambda 4959,5007}{H\beta }.$$
A single value of R<sub>23</sub> uniquely specifies O/H over most of the range in metallicity, however there is a turnover region (20–50% solar) where the relationship becomes double valued. McGaugh (1994) advocates the use of the \[NII\] $`\lambda `$6584/\[OII\] $`\lambda `$3727 ratio as a way to discriminate between upper and lower branches, noting that it varies monotonically with O/H and is not very sensitive to the ionization parameter since both ions have similar ionization potentials. The division between upper and lower branches is fairly well-defined, with the reddening-corrected log(\[NII\]/\[OII\])$`>`$1 indicating the upper branch and log(\[NII\]/\[OII\])$`<`$1 indicating the lower branch. Measurements of this line ratio in the HII regions under study here place them all securely on the lower, metal-poor branch of the R<sub>23</sub> relation (see Table 2). In most cases, the \[NII\]$`\lambda `$6548 line was too faint to measure accurately so we have assumed the theoretical value of \[NII\]$`\lambda `$6548$`=`$\[NII\]$`\lambda `$6584/2.95 (Mendoza & Zeippen, 1982) in calculating the N/O ratio. The derived oxygen abundances and nitrogen-to-oxygen ratios are low, ranging from 7$``$9% and 9$``$16% the solar value respectively. The dominant uncertainties in these estimates are the uncertainties in the model calibrations themselves, which are estimated to be $`\pm `$0.2 dex for log(O/H) and $`\pm `$0.1 dex for log(N/O) (see Ferguson, Gallagher & Wyse (1998) for a detailed discussion). With this in mind, there would appear to be little evidence for an intrinsic metallicity dispersion amongst these HII regions, although Obj #4 does seem marginally enhanced in N/O compared to the other two objects.
How do these chemical abundances compare to those of the M31 disk at deprojected location of And IV? Assuming a position angle of 35 and an inclination of 77.5 for the M31 disk, we calculate that And IV would have a deprojected radius of 108.4′ or $``$25 kpc if in the disk. The chemical abundance gradient at large radii in M31 is surprisingly poorly constrained, with the most distant measured HII region lying at only $``$16 kpc (Dennefeld & Kunth, 1981). Zaritsky, Kennicutt & Huchra (1994) quote values of 12+log(O/H)$`=`$9.03 at R$`=`$0.4R<sub>25</sub> and a gradient of $``$0.28 dex/R<sub>25</sub> for M31 (normalised to their adopted value of R$`{}_{25}{}^{}=`$77.4′ ) derived from the measurements of Dennefeld & Kunth (1981) and Blair, Kirshner & Chevalier (1982). Simple extrapolation of this gradient predicts a gas-phase oxygen abundance of $``$70% solar at the location of And IV. The HII regions in the vicinity of And IV therefore have metallicities which are roughly an order of magnitude lower than that expected for M31 disk gas at that radius. This finding adds to the mounting evidence that the emission line sources in the vicinity of And IV are not associated with the disk of M31.
### 3.4 Radial Velocities
Our long-slit spectra also provide a measurement of the radial velocity of each emission-line source. The average velocity and standard deviation determined from the observed wavelengths of the well-detected bright H$`\beta `$, \[OIII\]5007 and H$`\alpha `$ lines are reported for each HII region in Table 2 and are in excellent agreement. These velocities have been corrected by $``$15.5 km/s to account for the motion of the Earth around the sun. The mean of these averages is 256$`\pm `$9 km/s, which can be compared to the value of 248$`\pm `$47 km/s derived from the absorption lines in the underlying galactic continuum (see Figure 4). Both velocities differ significantly from the radial velocity of $``$375 km/s expected for M31’s disk at the projected location of And IV (Emerson, 1974), which is also the velocity reported previously in the literature for And IV. Our results therefore not only establish a direct association between the HII regions and underlying galactic continuum emission, but also provide the final piece of evidence that And IV is unassociated with the disk of M31. We adopt a heliocentric radial velocity of 256$`\pm `$9 km/s for And IV.
It is tempting to speculate on whether the spread of velocities that we measure in the And IV HII regions can be considered real. Interestingly, there does appear to be a systematic change of $``$30 km/s in the radial velocity across the face of And IV going from Objects #3 through 4 to 6, all of which were measured with a single slit position. Given the uncertainties in centroiding the emission lines, as well as the rms residuals of the wavelength solution ($``$ 0.5Å), and the moderate resolution of our spectra, the significance of this gradient should be considered marginal at present.
## 4 Discussion
We have presented a set of new observations which constrain the nature of the enigmatic object And IV. We find compelling evidence that And IV is *not* a star cluster lying in the disk of M31. This evidence includes: (i) the fact that the stellar population of And IV does not resolve to the same extent as that of the M31 field population, (ii) the discovery of individual HII regions in the vicinity of And IV which do not appear to be ionized by stars at the distance of M31, (iii) the metallicities of these HII regions are an order of magnitude lower than that expected for the M31 disk at the projected location of And IV, and (iv) the radial velocity of And IV differs by $``$600 km/s from that expected for the SW side of the M31 disk. Furthermore, the large velocity difference between And IV and M31 also makes it very unlikely that And IV is even a bound satellite of M31. The radial velocities of the known Andromeda satellites generally lie within $``$100 km/s of the systemic velocity of M31 (Mateo, 1998), whereas And IV differs from that by more than 500 km/s.
Two questions therefore remain: what is the true nature of And IV and where exactly does it lie? We begin by addressing the second issue, since it has bearing on the first. Given our measured heliocentric velocity, it is possible to use dynamical considerations to place a limit on the distance of And IV. Adopting the linear Virgocentric infall model of Schechter (1980) with parameters $`\gamma =2`$, V<sub>helio</sub>(Virgo)=976 km s<sup>-1</sup>, $`\omega _{\mathrm{}}`$=220 km s<sup>-1</sup> (Binggeli, Tammann & Sandage, 1987) and D<sub>virgo</sub>=15.9 Mpc (i.e. H<sub>o</sub>=75 km s<sup>-1</sup> Mpc<sup>-1</sup>), And IV’s position and heliocentric velocity imply a distance of 7.0 Mpc. Changing the heliocentric velocity of And IV by $`\pm 90`$ km s<sup>-1</sup> (ie. the typical magnitude of peculiar motion velocities) has the effect of changing the derived distance by $`1.2`$ Mpc, ie. by $`\pm 20\%`$.
Next, we return to the intriguing $`\mathrm{`}`$bump’ seen in the RGB LF of the stars in the vicinity of And IV, which we tentatively associate with the tip of the red giant branch (TRGB) population for a distant stellar system. The TRGB magnitude is known to be very stable at M$`{}_{I}{}^{}`$4 over a wide range of ages (2–15 Gyr) and metallicities, and is widely used as a distance indicator for resolved stellar systems (Lee, Freedman & Madore, 1993). In galaxies with a significant intermediate-age population, the presence of luminous asymptotic giant branch stars above the RGB tends to smear out the edge defining the tip; this effect is very likely to be present in And IV and implies that the distance we derive should be considered as a lower limit to the true distance. From visual inspection, we identify the TRGB from the I-band LF to lie at I$`=`$25.0$`\pm `$0.5 and derive a distance modulus of 29$`\pm `$0.5, corresponding to a linear distance of 6.3$`\pm `$1.5 Mpc. The agreement between this distance determination and that from a dynamical argument is extremely encouraging, especially given the significant uncertainties in each. Taking the average of these estimates, we therefore place And IV at a distance of $`6.7\pm 1.5`$ Mpc. Reassuringly, this distance implies that the most luminous blue stars detected in the vicinity of the And IV HII regions would have absolute magnitudes in the range of M$`{}_{V}{}^{}5`$ to $`7`$, and would therefore correspond to luminous OB stars and clusters, broadly consistent with the observed ionization.
We now turn to clarifying the nature of And IV. The properties established in this paper – moderately low surface brightness, very blue colour and low metallicity – are reminiscent of those observed in $`\mathrm{`}`$typical’ low mass dwarf irregular galaxies (eg. Mateo (1998); Miller (1994)). And IV’s extent of $`40`$″ and disk scalelength of $`11`$″ correspond to linear sizes of $`1.3`$ kpc and 360 pc respectively at a distance of 6.7 Mpc, confirming that the galaxy is indeed physically small. Many dwarf irregulars exhibit a small amount of ongoing star formation, with rates ranging from 0.0001$``$0.01 M yr<sup>-1</sup> (Mateo, 1998; Miller, 1994; Hunter, Hawley & Gallagher, 1993). Summing the H$`\alpha `$ flux (corrected for Galactic reddening but not that internal to And IV itself) from the 5 HII regions detected within $``$1′ of And IV, we derive a current star formation rate of 2.6$`\times 10^5`$ (D/Mpc)<sup>2</sup> M yr<sup>-1</sup> using the proportionality between SFR and H$`\alpha `$ luminosity derived by Kennicutt, Tamblyn & Congdon (1994). At our derived distance of 6.7 Mpc, this translates into 0.001 M yr<sup>-1</sup> and is therefore highly consistent with the rates measured in local dwarfs. Yet another constraint is provided by the gas-phase metallicity of And IV. Based on the mean oxygen abundance of $`10\%`$ solar measured for the And IV HII regions, the metallicity-luminosity relation of Skillman, Kennicutt & Hodge (1989) predicts M$`{}_{B}{}^{}=`$15 which, once again, supports the identification of And IV as a dwarf galaxy.
In terms of properties such as star formation rate, metallicity, central surface brightness and inferred luminosity, And IV appears very similar to Local Group dwarf irregulars IC 1613 and Sextans A (Mateo, 1998). High-quality HST CMDs have recently been published for both of these systems (Dohm-Palmer et al., 1997; Cole et al., 1999), and we consider how these diagrams may help us to better understand the nature of the $`\mathrm{`}`$excess’ resolved stars seen towards And IV (Section 3.1.1). One of the most striking features of the Sextans A CMD is the population of massive core helium-burning stars (the so-called $`\mathrm{`}`$blue loop’ stars) seen just redward of the main sequence; this feature is also seen, albeit to a slightly lesser degree, in IC 1613. As the prominence of blue loop stars is greatest at low metallicites, and as (at least) the gas-phase metallicity of And IV and Sextans A/IC 1613 are all similarly low, it is reasonable to expect that such stars are also present in the And IV CMD. The $`\mathrm{`}`$excess’ stars seen lying between the M31 blue plume and RGB on the And IV CMD, but not on that of the M31 field, could be the signature of this component. Furthermore, IC 1613 shows a sizeable population of red supergiants extending to I$``$16.5 (or V$``$18), which corresponds to M$`{}_{V}{}^{}6.5`$ for a distance modulus of 24.27 (Cole et al., 1999). Given that the $`\mathrm{`}`$excess’ luminous red stars identified on the And IV CMD would have absolute magnitudes in the range $`7.5`$ to $`6.5`$ at our derived distance, it appears likely that these stars are red supergiants belonging to And IV.
An interesting question is whether And IV is an isolated dwarf galaxy or whether it belongs to some larger group environment. There are no catalogued galaxy groups in the vicinity of And IV but a search with NED reveals 14 galaxies in the range 22$`{}_{}{}^{h}<\alpha <0^h`$, 20$`{}_{}{}^{}<\delta <60^{}`$ and 150 km/s$`<`$V$`{}_{helio}{}^{}<`$350 km/s. Of these, the most luminous are IC 1727, NGC 784 and UGC 64. We suspect that And IV may be one of the many low-luminosity galaxies inhabiting this environment.
## 5 Conclusions
We have presented deep HST WFPC2 and ground-based observations of the enigmatic object And IV. The true nature of this object – old $`\mathrm{`}`$star cloud’ in the outer disk of M31 or background galaxy – has remained a mystery since it was first discovered by van den Bergh (1972) during his search for dwarf spheroidal companions to M31. From the analysis of our WFPC2 images and complementary H$`\alpha `$ imaging and long-slit optical spectroscopy, we find compelling evidence that And IV is a background galaxy seen through the disk of M31. The moderate surface brightness ($`\overline{\mu _V}`$24), very blue colour (V$``$I$``$0.6), low current star formation rate ($``$0.001 M yr<sup>-1</sup>) and low metallicity ($``$10% solar) reported here are consistent with And IV being a dwarf irregular galaxy, perhaps a more distant analog of Local Group members IC 1613 and Sextans A. Indeed, such objects are very common in the nearby Universe, and perhaps it is not surprising to find one projected behind M31. The distance to And IV is not tightly constrained by our current dataset, but arguments based on both the observed radial velocity and on a tentative detection of the RGB tip suggest it lies in the range 5$``$D$``$8 Mpc, placing it well outside the confines of the Local Group. At this distance, the physical extent of And IV is consistent for what is expected of small dwarf galaxies. And IV may belong to a loose, previously uncatalogued group, containing major members UGC 64, IC1727 and NGC 784.
We are grateful to Rachel Johnson, Mike Irwin and Nial Tanvir for useful advice and discussions. Gerhardt Meurer is thanked for assistance with Virgocentric infall calculation and Piero Rosati for help with the radial velocity cross-correlation. This research has made use of the NASA/IPAC Extragalactic Database (NED) which is operated by the Jet Propulsion Laboratory, California Institute of Technology, under contract with NASA. Support for this work was provided by NASA through grant number GO-067340195A from the Space Telescope Science Institute, which is operated by AURA, Inc. under NASA contract NAS5-26555.
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# Statistical Entropy of a Stationary Dilaton Black Hole from Cardy Formula
## I INTRODUCTION
Much effort has been concentrated on the statistical mechanical description of the Bekenstein-Hawking black hole entropy \- in terms of microscope states both in string theory and in “quantum geometry” . Strominger calculated the entropy of black holes whose near-horizon geometry is locally $`AdS_3`$ from the asymptotic growth of states. Carlip derived the central extension of the constraint algebra of general relativity by using Brown-Henneaux-Strominger’s approach and manifestly covariant phase space methods \- . He found that a natural set of boundary conditions on the (local) Killing horizon leads to a Virasoro subalgebra with a calculable central charge and the standard Cardy formula gives the Bekenstein-Hawking entropies of some black holes. Those works show a suggestion that the asymptotic behavior of the density of states may be determined by the algebra of diffeomorphism at horizon. Solodukhin obtained same result by a analysis of the Liouville theory near the horizon obtained from dimensional reduction of Einstein gravity. Das, Ghosh, and Mitra studied the statistical entropy of a Schwarzschild black string in five dimensions by counting the black string states which from a representation of the near-horizon conformal symmetry with a central charge. Recently, we extended Carlip’s investigation for vacuum case to a case including a cosmological term and electromagnetic fields and calculated the statistical entropies of the Kerr-Newman black hole and the Kerr-Newman-AdS black hole by using standard Cardy formula.
On the other hand, the quantum correction to entropy of the black hole is an interesting topic-. Recently, Kaul and Majumdar computed the lowest order corrections to the Bekenstein-Hawking entropy in a particular formulation of the “quantum geometry” program of Ashtekar et al. They showed that the leading corrections is a logarithmic term, i. e., the entropy is
$$S\frac{A_H}{4}\frac{3}{2}\mathrm{ln}\frac{A_H}{4}+const.+\mathrm{},$$
(1)
where $`A_H`$ is the event horizon area. Carlip also calculated the quantum corrections to black hole entropy by the Cardy formula and found that the entropy can be expressed as
$$SS_0\frac{3}{2}\mathrm{ln}S_0+\mathrm{ln}c+const.+\mathrm{},$$
(2)
where $`S_0`$ is standard Bekenstein-Hawking entropy and $`c`$ is a central charge of a Virasoro subalgebra. Carlip pointed out that if the central charge is the sense of being independent of the horizon area (Carlip thinks that this can be done by adjust the periodicity $`\beta `$ ), then the factor of $`3/2`$ in logarithmic term will always appear.
We all know that four dimensional dilaton charged black hole obtained in the low-energy effective field theory describing strings have qualitatively different properties from those that appear in the ordinary Einstein gravity. Therefore, it is worth to investigate whether or not the Carlip’s conclusion (the asymptotic behavior of the density of states may be determined by the algebra of diffeomorphism at horizon) and Kaul and Majumdar’s result (the leading corrections to the entropy is a logarithm of the horizon area with a factor $`3/2`$) are valid for the static and stationary dilaton black hole.
We begin in Section II by using the covariant phase techniques to extend Carlip’s investigation for vacuum case $`𝐋_{a_1a_2\mathrm{}a_n}=\frac{1}{16\pi G}ϵ_{a_1a_2\mathrm{}a_n}R`$ to a case for gravity coupled to a Maxwell field and a dilaton, i.e., the Lagrangian n-form is described by $`𝐋_{a_1a_2\mathrm{}a_n}=ϵ_{a_1a_2\mathrm{}a_n}\left[R2(\varphi )^2e^{2\alpha \varphi }F^2\right].`$ A constraint algebra is obtained. In Sec. III, the standard Virasoro subalgebra with corresponding central charges is constructed for stationary dilation black hole. The statistical entropy of the black hole is then calculated by using standard Cardy formula. In Sec. IV, a new Cardy formula is obtained and then the first-order quantum correction to the entropy is studied. The last section devotes to discussion and summary.
## II Algebra of diffeomorphism on the Killing Horizon
Let $`\xi ^a`$ be any smooth vector fields on a spacetime manifold $`𝐌`$, i. e., $`\xi ^a`$ is the infinitesimal generator of a diffeomorphism, Lee, Wald, and Iyer showed that the Lagrangian $`𝐋`$, equation of motion n-form $`𝐄`$, symplectic potential (n-1)-form $`𝚯`$, Noether current (n-1)-form $`𝐉`$, and Noether charge (n-2)-form $`𝐐`$ satisfy following relations
$`\delta 𝐋`$ $`=`$ $`𝐄\delta \varphi +d𝚯,`$ (3)
$`𝐉[\xi ]`$ $`=`$ $`𝚯[\varphi ,_\xi \varphi ]\xi 𝐋,`$ (4)
$`𝐉`$ $`=`$ $`d𝐐,`$ (5)
here and hereafter the “central dot” denotes the contraction of the vector field $`\xi ^a`$ into the first index of the differential form. Hamilton’s equation of motion is given by
$$\delta H[\xi ]=_C\omega [\varphi ,\delta \varphi ,_\xi \varphi ]=_C[\delta 𝐉[\xi ]d(\xi 𝚯[\varphi ,\delta \varphi ])].$$
(6)
By using Eq. (5) and defining a (n-1)-form $`𝐁`$ as
$$\delta _C\xi 𝐁[\varphi ]=_C\xi 𝚯[\varphi .\delta \varphi ],$$
(7)
the Hamiltonian can be expressed as
$$H[\xi ]=_C(𝐐[\xi ]\xi 𝐁[\varphi ]).$$
(8)
The Poisson bracket forms a standard “surface deformation algebra”
$$\{H[\xi _1],H[\xi _2]\}=H[\{\xi _1,\xi _2\}]+K[\xi _1,\xi _2],$$
(9)
where the central term $`K[\xi _1,\xi _2]`$ depends on the dynamical fields only through their boundary values.
The four dimensional low-energy Lagrangian obtained from string theory is
$$𝐋_{abcd}=ϵ_{abcd}\left[R2(\varphi )^2e^{2\alpha \varphi }F^2\right],$$
(10)
where $`ϵ_{abcd}`$ is the volume element, $`\varphi `$ is the dilaton scalar field, $`F_{ab}`$ is the Maxwell field associated with a $`U(1)`$ sub-group of $`E_8\times E_8`$ or $`Spin(32)/Z_2`$, and $`\alpha `$ is a free parameter which governs the strength of the coupling of the dilaton to the Maxwell field. The reason we set the remaining gauge fields and antisymmetric tensor field $`H_{\mu \nu \rho }`$ to zero is that the metrics of stationary and static dilaton black holes are almost obtained form the Lagrangian (10). We know from Lagrangian (10) that the equations of motion $`𝐄`$ for dynamical fields $`A_\mu `$, $`\varphi `$, and $`g_{\mu \nu }`$ can be respectively given by
$$_\mu (e^{2\alpha \varphi }F^{\mu \nu })=0,$$
(11)
$$^2\varphi +\frac{1}{2}e^{2\alpha \varphi }F_{\mu \nu }F^{\mu \nu }=0,$$
(12)
$$R_{\mu \nu }\frac{1}{2}g_{\mu \nu }R=2_\mu \varphi _\nu \varphi g_{\mu \nu }(\varphi )^2+2e^{2\alpha \varphi }F_{\beta \nu }F_\mu ^\beta \frac{1}{2}g_{\mu \nu }e^{2\alpha \varphi }F_{\mu \nu }F^{\mu \nu }.$$
(13)
The symplectic potential (n-1)-form is
$$𝚯_{bcd}[g,_\xi g]=4ϵ_{abcd}\left\{\frac{1}{2}(_e^{[e}\xi ^{a]}+R_e^a\xi ^e)\xi ^e_e\varphi ^a\varphi e^{2\alpha \varphi }F^{af}\left[F_{ef}\xi ^e+(\xi ^eA_e)_{;f}\right]\right\}.$$
(14)
From Eqs. (4) and (14) we have
$`𝐉_{bcd}`$ $`=`$ $`2ϵ_{abcd}\{_e^{[e}\xi ^{a]}2e^{2\alpha \varphi }F^{af}(\xi ^eA_e)_{;f}+[R_e^a{\displaystyle \frac{1}{2}}\delta _e^aR2_e\varphi ^a\varphi +\delta _e^a(\varphi )^2`$ (16)
$`2e^{2\alpha \varphi }F^{af}F_{ef}+{\displaystyle \frac{1}{2}}\delta _e^ae^{2\alpha \varphi }F^2]\xi ^e\}`$
$`=`$ $`2ϵ_{abcd}\left[_e^{[e}\xi ^{a]}2e^{2\alpha \varphi }F^{af}(\xi ^eA_e)_{;f}\right]`$ (17)
$`=`$ $`2ϵ_{abcd}\left[_e^{[e}\xi ^{a]}+4_f(e^{2\alpha \varphi }^{[f}A^{a]}A_e\xi ^e)\right],`$ (18)
in the second and third lines, we used the equations of motion (13) and (11). Eqs. (5) and (18) show that
$$𝐐_{cd}=ϵ_{abcd}\left[^a\xi ^b+4e^{2\alpha \varphi }A_e\xi ^e^aA^b\right].$$
(19)
For a stationary dilaton black hole, the dilaton scalar field, the electromagnetic potential $`A_a`$, and the Killing vector can be respectively expressed as
$`\varphi `$ $`=`$ $`\varphi (r,\theta ),`$ (20)
$`A_a`$ $`=`$ $`(A_0(r,\theta ),A_1(r,\theta ),A_2(r,\theta ),A_3(r,\theta )),`$ (21)
$`\chi _H^a`$ $`=`$ $`\chi _H^{(t)}+\chi _H^{(\phi )}=(1,0,0,\mathrm{\Omega }_H),`$ (22)
where the vector $`\chi _H^{(t)}`$ correspond to time translation invariance, $`\chi _H^{(\phi )}`$ to rotational symmetry, and $`\mathrm{\Omega }_H=(g_{t\phi }/g_{\phi \phi })_H`$ is the angular velocity of the black hole.
As Carlip did in Ref. we define a “stretched horizon” $`\chi ^2=ϵ`$, where $`\chi ^2=g_{ab}\chi ^a\chi ^b`$, $`\chi ^a`$ is a Killing vector. The result of the computation will be evaluated at the event horizon of the black hole by taking $`ϵ`$ to zero. Near the stretched horizon, one can introduce a vector orthogonal to the orbit of $`\chi ^a`$ by $`_a\chi ^2=2\kappa \rho _a,`$ where $`\kappa `$ is the surface gravity. The vector $`\rho ^a`$ satisfies conditions
$`\chi ^a\rho _a={\displaystyle \frac{1}{\kappa }}\chi ^a\chi ^b_a\chi _b=0,`$ everywhere (23)
$`\rho ^a\chi ^a,`$ $`\text{at the horizon}.`$ (24)
To preserve “asymptotic” structure at horizon, we impose Carlip’s boundary conditions
$$\delta \chi ^2=0,\chi ^at^b\delta g_{ab}=0,\delta \rho _a=\frac{1}{2\kappa }_a(\delta \chi ^2)=0,at\chi ^2=0,$$
(25)
where $`t^a`$ is a any unit spacelike vector tangent to boundary $`𝐌`$ of the spacetime $`𝐌`$. And the infinitesimal generator of a diffeomorphism is taken as
$$\xi ^a=\rho ^a+𝒯\chi ^a,$$
(26)
where functions $``$ and $`𝒯`$ obey the relations
$`={\displaystyle \frac{1}{\kappa }}{\displaystyle \frac{\chi ^2}{\rho ^2}}\chi ^a_a𝒯,`$ everywhere (27)
$`\rho ^a_a𝒯=0,`$ $`\text{at the horizon}.`$ (28)
For a one-parameter group of diffeomorphism such that $`D𝒯_\alpha =\lambda _\alpha 𝒯_\alpha `$, ( $`D\chi ^a_a`$ ), one introduces an orthogonality relation
$$_C\widehat{ϵ}𝒯_\alpha 𝒯_\beta \delta _{\alpha +\beta }.$$
(29)
The technical role of the condition (29) is to guarantee the existence of generators $`H[\xi ]`$. By using the other future-directed null normal vector $`N^a=k^a\alpha \chi ^at^a,`$ with $`k^a=\frac{1}{\chi ^2}\left(\chi ^a\frac{|\chi |}{\rho }\rho ^a\right)`$ and a normalization $`N_a\chi ^a=1`$, the volume element can be expressed as
$$ϵ_{abcd}=\widehat{ϵ}_{cd}(\chi _aN_b\chi _bN_a)+\mathrm{}\mathrm{},$$
(30)
the omitted terms do not contribute to the integral.
Form the right hand of Eq. (7)
$$_C\xi ^b𝚯_{bcd}=4_Cϵ_{abcd}\xi ^a\left\{\frac{1}{2}(_e^{[e}\xi ^{b]}+R_e^b\xi ^e)\xi ^e_e\varphi ^be^{2\alpha \varphi }F^{fb}\left[F_{ef}\xi ^e+(\xi ^eA_e)_{;f}\right]\right\},$$
(31)
we know that the first two terms in the right hand of Eq. (31) can be treated as Carlip did in Ref. . At the horizon, by using Eqs. (20), (21), (22) and (25)- (30) we obtain
$`{\displaystyle _C}ϵ_{abcd}\xi _2^a\xi _1^e^b\varphi _e\varphi =0,`$ (32)
and
$`{\displaystyle _C}ϵ_{abcd}e^{2\alpha \varphi }\xi ^aF^{bf}\left[F_{ef}\xi ^e+(\xi ^eA_e)_{;f}\right]`$ (33)
$`=`$ $`{\displaystyle _C}ϵ_{abcd}e^{2\alpha \varphi }\xi ^aF^{bf}\delta _\xi A_f`$ (34)
$`=`$ $`{\displaystyle _C}\widehat{ϵ}_{cd}e^{2\alpha \varphi }\left[{\displaystyle \frac{|\chi |}{\rho }}𝒯\rho _b+\left({\displaystyle \frac{\rho }{|\chi |}}+t\rho \right)\chi _b\right]F^{bf}\delta _\xi A_f`$ (35)
$`=`$ $`0.`$ (36)
Therefore we know that the last three terms in Eq. (31) also gives no contribution to $`K[\xi _1,\xi _2]`$.
By applying Eqs. (21), (22), (26), and (30), we can show that, at the horizon, $`_Cϵ_{abcd}e^{2\alpha \varphi }A_e\xi ^e^aA^b0.`$ Hence, from Eq. (19) we find
$$_CQ_{cd}=_Cϵ_{abcd}^a\xi ^b.$$
(37)
Denoting by $`\delta _\xi `$ the variation corresponding to diffeomorphism generated by $`\xi `$, for the Noether current we have $`\delta _{\xi _2}𝐉[\xi _1]=d[\xi _2(𝚯[\varphi ,_{\xi _1}\varphi ]\xi _1𝐋)].`$ Substituting it into Eq. (6) and using Eq. (14) we obtain
$`\delta _{\xi _2}H[\xi _1]`$ $`=`$ $`{\displaystyle _C}\left(\xi _2𝚯[\varphi ,_{\xi _1}\varphi ]\xi _1𝚯[\varphi ,_{\xi _2}\varphi ]\xi _2\xi _1𝐋\right)`$ (38)
$`=`$ $`{\displaystyle _C}ϵ_{abcd}\left[\xi _2^a_e(^e\xi _1^b^b\xi _1^e)\xi _1^a_e(^e\xi _2^b^b\xi _2^e)\right]`$ (42)
$`4{\displaystyle _C}ϵ_{abcd}e^{2\alpha \varphi }\left\{\xi _2^aF^{fb}\left[F_{ef}\xi _1^e+(\xi _1^eA_e)_{;f}\right]\xi _1^aF^{fb}\left[F_{ef}\xi _2^e+(\xi _2^eA_e)_{;f}\right]\right\}`$
$`{\displaystyle _C}ϵ_{abcd}\left[4R_e^b(\xi _1^a\xi _2^e\xi _2^a\xi _1^e)+\xi _2^a\xi _1^b𝐋\right]`$
$`4{\displaystyle _C}ϵ_{abcd}\left(\xi _2^a\xi _1^e\xi _1^a\xi _2^e\right)^b\varphi _e\varphi .`$
At the horizon, applying Eqs. (20), (21), (22) and (25)- (30) we see that
$`{\displaystyle _C}ϵ_{abcd}(\xi _2^a\xi _1^e\xi _1^a\xi _2^e)^b\varphi _e\varphi `$ (43)
$`=`$ $`{\displaystyle _C}\widehat{ϵ}_{cd}\left({\displaystyle \frac{1}{\kappa }}{\displaystyle \frac{\chi ^2}{\rho ^2}}\right)\left[{\displaystyle \frac{|\chi |}{\rho }}\rho _b\rho ^e\left({\displaystyle \frac{\rho }{|\chi |}}+t\rho \right)\chi _b\chi ^e\right](𝒯_2D𝒯_1𝒯_1D𝒯_2)^b\varphi _e\varphi `$ (44)
$`=`$ $`0,`$ (45)
$`{\displaystyle _C}ϵ_{abcd}\xi _2^a\xi _1^b𝐋`$ (46)
$`=`$ $`{\displaystyle _C}\widehat{ϵ}_{cd}𝐋\left[{\displaystyle \frac{|\chi |}{\rho }}𝒯_2\rho _b+\left({\displaystyle \frac{\rho }{|\chi |}}+t\rho \right)_2\chi _b\right](𝒯_1\chi ^b+_1\rho ^b)`$ (47)
$`=`$ $`{\displaystyle _C}\widehat{ϵ}_{cd}𝐋\left[{\displaystyle \frac{|\chi |}{\rho }}𝒯_2_1\rho ^2+\left({\displaystyle \frac{\rho }{|\chi |}}+t\rho \right)_2𝒯_1\chi ^2\right]`$ (48)
$`=`$ $`0,`$ (49)
and
$`{\displaystyle _C}ϵ_{abcd}R_e^b(\xi _1^a\xi _2^e\xi _2^a\xi _1^e)`$ (50)
$`=`$ $`{\displaystyle _C}\widehat{ϵ}_{cd}R_e^b\left({\displaystyle \frac{1}{\kappa }}{\displaystyle \frac{\chi ^2}{\rho ^2}}\right)\left[{\displaystyle \frac{|\chi |}{\rho }}\rho _b\rho ^e\left({\displaystyle \frac{\rho }{|\chi |}}+t\rho \right)\chi _b\chi ^e\right](𝒯_1D𝒯_2𝒯_2D𝒯_1)`$ (51)
$`=`$ $`0.`$ (52)
Substituting Eqs. (43), (36), (49) and (50) into Eq. (38) we find
$`\delta _{\xi _2}H[\xi _1]`$ $`=`$ $`{\displaystyle _C}ϵ_{abcd}\left[\xi _2^a_e(^e\xi _1^b^b\xi _1^e)\xi _1^a_e(^e\xi _2^b^b\xi _2^e)\right].`$ (53)
We can interpret the left side of Eq. (38) the variation of the boundary term $`J`$ since the “bulk” part of the generator $`H[\xi _1]`$ on the left side vanishes on shell. On the other hand, the change in $`J[\xi _1]`$ under a surface deformation generated by $`J[\xi _2]`$ can be precisely described by Dirac bracket $`\{J[\xi _1],j[\xi _2]\}^{}`$ . Thus we have
$$\{J[\xi _1],J[\xi _2]\}^{}=_Cϵ_{abcd}\left[\xi _2^a_e(^e\xi _1^b^b\xi _1^e)\xi _1^a_e(^e\xi _2^b^b\xi _2^e)\right].$$
(54)
Inserting Eqs. (26), (27) and (30) into (54) we obtain
$`\{J[\xi _1],J[\xi _2]\}^{}`$ $`=`$ $`{\displaystyle _C}\widehat{ϵ}_{cd}\left[{\displaystyle \frac{1}{\kappa }}(𝒯_1D^3𝒯_2𝒯_2D^3𝒯_1)2\kappa (𝒯_1D𝒯_2𝒯_2D𝒯_1)\right].`$ (55)
For any one-parameter group of diffeomorphism satisfying conditions (26) and (27), it is also easy to check that
$$\{\xi _1,\xi _2\}^a=(𝒯_1D𝒯_2𝒯_2D𝒯_1)\chi ^a+\frac{1}{\kappa }\frac{\chi ^2}{\rho ^2}D(𝒯_1D𝒯_2𝒯_2D𝒯_1)\rho ^a.$$
(56)
The Hamiltonian (8) consists of two terms, but Eqs (43) and (36) and discussion about $`\xi 𝚯`$ in Ref. show that the second terms make no contribution. Then, we have
$$J[\{\xi _1,\xi _2\}]=_C\widehat{ϵ}_{cd}\left[2\kappa (𝒯_1D𝒯_2𝒯_2D𝒯_1)\frac{1}{\kappa }D(𝒯_1D^2𝒯_2𝒯_2D^2𝒯_1)\right].$$
(57)
On shell Eq. (9) can be expressed as
$$\{J[\xi _1],J[\xi _2]\}^{}=J[\{\xi _1,\xi _2\}]+K[\xi _1,\xi _2].$$
(58)
Therefore, we know that from Eqs. (55) and (57) the central term is
$$K[\xi _1,\xi _2]=_C\widehat{ϵ}_{cd}\frac{1}{\kappa }(D𝒯_1D^2𝒯_2D𝒯_2D^2𝒯_1).$$
(59)
It is interesting to note that the constraint algebra (58) with Eqs. (55), (57), and (59) has same form as that for the vacuum case. In next section, we will study statistical-mechanical entropies of some stationary dilaton black holes by using the constraint algebra and conformal field theory methods.
## III Statistical Entropy of stationary dilaton black hole
In order to construct a standard Virasoro subalgebra from constraint algebra (55) and (57)-(59), as Cadoni, Mignemi and Carlip did in references we define a new generator $`𝑑vJ`$ in which the function $`v`$ takes period $`T`$. Form stationary conditions (22) we know that a one-parameter group of diffeomorphism satisfying Eqs. (29) and (56) can be taken as
$$𝒯_n=\frac{T}{2\pi }exp\left[in(\frac{2\pi }{T}v+C_\alpha (\phi \mathrm{\Omega }_Hv))\right],$$
(60)
where $`C_\alpha `$ is a arbitrary constant. We should note that one-parameter group (60) is also valid for static black hole since it is a special case of the stationary black hole with $`\mathrm{\Omega }_H=0`$. Substituting Eq. (60) into central term (59) and using condition (29) we obtain
$$K[𝒯_m,𝒯_n]=\frac{iA_H}{8\pi }\frac{2\pi }{\kappa T}m^3\delta _{m+n,0},$$
(61)
where $`A_H=_C\widehat{ϵ}_{cd}`$ is the area of the event horizon. Eq. (58) thus takes standard form of a Virasoro algebra
$$i\{J[𝒯_m],J[𝒯_n]\}=(mn)J[𝒯_{m+n}]+\frac{c}{12}m^3\delta _{m+n,0},$$
(62)
with central charge
$$\frac{c}{12}=\frac{A_H}{8\pi }\frac{2\pi }{\kappa T}.$$
(63)
The boundary term $`J[𝒯_0]`$ can easily be obtained by using Eqs (5), (19), and (60), which is given by
$$J[𝒯_0]=\mathrm{}=\frac{A_H}{8\pi }\frac{\kappa T}{2\pi }.$$
(64)
From standard Cardy’s formula
$$\rho (\mathrm{})exp\left\{2\pi \sqrt{\frac{c}{6}\left(\mathrm{}\frac{c}{24}\right)}\right\},$$
(65)
we know that the number of states with a given eigenvalue $`\mathrm{}`$ of $`J[𝒯_0]`$ grows asymptotically for large $`\mathrm{}`$ as
$$\rho (\mathrm{})exp\left[\frac{A_H}{4}\sqrt{2\left(\frac{2\pi }{\kappa T}\right)^2}\right].$$
(66)
Only if we take the period $`T`$ as the periodicity of the Euclidean black hole, i.e.,
$$T=\frac{2\pi }{\kappa },$$
(67)
the statistical entropy of the stationary dilaton black hole
$$S_0ln\rho (\mathrm{})=\frac{A_H}{4},$$
(68)
coincides with the standard Bekenstein-Hawking entropy.
## IV Logarithmic corrections to black hole entropy
Now lets us consider the first-order quantum correction to the entropy. In oder to do that, we should first derive the logarithmic corrections to the Cardy formula.
In references , Carlip showed that the number of states is
$$\rho (\mathrm{})=𝑑\tau e^{2\pi i\mathrm{}\tau }e^{2\pi i\mathrm{}_0\frac{1}{\tau }}e^{\frac{2\pi ic}{24}\tau }e^{\frac{2\pi ic}{24}\frac{1}{\tau }}\stackrel{~}{Z}(1/\tau ),$$
(69)
where $`\stackrel{~}{Z}(1/\tau )`$ approaches to a constants, $`\rho (\mathrm{}_0)`$, for large $`\tau `$. So the integral (69) can be evaluated by steepest descent provided that the imaginary part of $`\tau `$ is large at the saddle point.
The integral takes the form
$$I[a,b]=𝑑\tau e^{2\pi ia\tau +\frac{2\pi ib}{\tau }}f(\tau ).$$
(70)
The argument of the exponent is extremal at $`\tau _0=\sqrt{\frac{b}{a}}`$, and expanding around $`\tau _0`$, one has
$$I[a,b]𝑑\tau e^{4\pi ia\sqrt{ab}+\frac{2\pi ib}{\tau _0^3}(\tau \tau _0)^2}f(\tau _0)=\left(\frac{b}{4a^3}\right)^{1/4}e^{4\pi i\sqrt{ab}}.$$
(71)
Comparing Eqs. (69) with (70) we know
$$a=\frac{c}{24}\mathrm{},b=\frac{c}{24}\mathrm{}_0.$$
(72)
Therefore, for large $`\mathrm{}`$, if we let $`c_{eff}=c24\mathrm{}_0,`$ the number of states can be expressed as
$$\rho _{cq}(\mathrm{})\left[\frac{c_{eff}}{96\left(\mathrm{}\frac{c}{24}\right)^3}\right]^{1/4}exp\left\{2\pi \sqrt{\frac{c_{eff}}{6}\left(\mathrm{}\frac{c}{24}\right)}\right\}\rho (\mathrm{}_0).$$
(73)
The exponential part in (73) gives the Carlip’s result (C.3) in Appendix C in Ref. , the factor before the exponent devotes the logarithmic correction to black hole entropy.
By Using the central charge (63), eigenvalue (64), constraint condition of the period (67), and new Cardy formula (73), we know that the statistical entropy including first-order quantum correction is given by
$`S`$ $`=`$ $`{\displaystyle \frac{A_H}{4}}{\displaystyle \frac{3}{2}}\mathrm{ln}{\displaystyle \frac{A_H}{4}}+\mathrm{ln}c+const.,`$ (74)
$`=`$ $`{\displaystyle \frac{A_H}{4}}{\displaystyle \frac{1}{2}}\mathrm{ln}{\displaystyle \frac{A_H}{4}}+const..`$ (75)
The first line has two logarithmic terms and agrees with Carlip’s results (2) . However, after we take $`T=\frac{2\pi }{\kappa }`$, the second shows that the factor of the logarithmic term becomes $`\frac{1}{2}`$, which is different from Kaul and Majumdar’s result $`\frac{3}{2}`$.
## V summary and discussion
We extend Carlip’s investigation in Ref. to four dimensional low-energy string theory by the covariant phase techniques. With Carlip’s boundary conditions, a standard Virasoro subalgebra with corresponding central charge for stationary dilaton black hole is constructed at a Killing horizon. We find that only we take $`T`$ as the periodicity of the Euclidean black hole, $`T=\frac{2\pi }{\kappa }`$, the statistical entropy of the stationary dilaton black hole yielded by standard Cardy formula agrees with its Bekenstein-Hawking entropy. Therefore, Carlip’s conclusion—the asymptotic behavior of the density of states may be determined by the algebra of diffeomorphism at horizon—is valid for stationary dilaton black holes obtained from the low-energy effective field theory with Lagrangian (10).
When we consider first-order quantum correction the entropy contains extra logarithmic terms which agrees with Carlip’s results (2) . However, from above discussions we know that in order to get the Bekenstein-Hawking entropy we have to take $`T=\frac{2\pi }{\kappa }`$. That is to say, we can not set central charge $`c`$ to be a universal constant, independent of area of the event horizon, by adjusting periodicity $`T`$ as Carlip suggested in Ref. . Therefore, the factor of the logarithmic term is $`\frac{1}{2}`$, which is different from Kaul and Majumdar’s result, $`\frac{3}{2}`$.
From the derivation given in the section IV we know that the new Cardy formula (73) is valid for general black hole whether or not the black hole is dilatonic. Hence, the factor of the logarithmic term will be $`\frac{1}{2}`$ as long as the spacetime is such that (59) is obey. This means that the discrepancy between Carlip’s approach and that of Kaul and Majumdar is not just for the dilaton black hole, but for any black hole which respects (59), where T is the periodicity of the Euclidean black hole.
###### Acknowledgements.
This work was supported in part by the National Natural Science Foundation of China under grant No. 19975018, Theoretical Physics Foundation of China under grant No. 19947004 and National Foundation of China through C. N. Yang.
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# Eliminating the mean-field shift in multicomponent Bose-Einstein condensates
The promise of Bose-Einstein condensates and atom lasers as sources for atom interferometry based sensors results from their high brightness and coherence which leads to an increase in the signal-to-noise ratio as compared to conventional atom optics. Additionally, further enhancement of sensitivity might be achieved by taking advantage of nonlinear effects which occur in quantum-degenerate atomic fields. For example, feedback between optical and/or matter-wave fields can result in nonlinear instabilities. Such positive feedback between optical and atomic intensity gratings has already led to the design of matter-wave amplifiers. The utility of such instabilities for nonlinear interferometry lies in the associated increase in sensitivity that can be achieved by operating just above or just below the instability threshold. This allows a small perturbation to produce a large change in the properties of the system.
The transition to high density atomic samples carries a price, however, in that atom-atom collisions may introduce unwanted nonlinear (density dependent) phase-front distortions which limit sensitivity. While collisions are often viewed as a source of decoherence in quantum optics, in the regime of coherent atomic matter waves atom-atom interactions lead to coherent nonlinear wave-mixing and can therefore be manipulated using techniques inspired from nonlinear optics. In this Letter we discuss how the density-dependent phase shift can be eliminated in a trapped multi-component condensate. With applications in atom interferometry in mind we consider specifically the case of a narrow atomic waveguide such as those recently microfabricated on glass chips , as these devices hold great promise for the development of integrated atom-interferometric devices.
At temperature $`T0`$ a single-component Bose-Einstein condensate is characterized by a scalar order parameter whose evolution is governed by a nonlinear Schrödinger equation (NLSE). At densities low enough that three-body collisions can be neglected the NLSE contains a cubic nonlinearity whose form is determined by the two-body collision potential. In the limit of $`s`$-wave scattering, the potential is of the form $`V=(4\pi \mathrm{}^2na/m)\delta (𝐫_{12})`$, where $`m`$ is the particle mass, $`n`$ the atomic number density, $`𝐫_{12}`$ the relative position of the atoms, and $`a`$ the $`s`$-wave scattering length. This leads to a density-dependent phase shift in the evolution of the condensate wave function. In the language of nonlinear optics this process is known as self-phase modulation. Since the exact density of the condensate is usually not perfectly known (especially in cases where it has been divided by a ‘beam-splitter’) this is a source of uncertainty that limits the accuracy with which precision measurements can be done.
There is at first sight nothing that can be done to eliminate this shift short of reducing the condensate density to a point where it becomes negligible, or of taking advantage of Feshbach resonances, in which case three-body collisions appear to become a serious problem. The situation is quite different, however, for multi-component condensates. In addition to a self-phase modulation proportional to its own density, each condensate component experiences in that case an additional cross-phase modulation, i.e. a phase shift proportional to the density of the other component. As we will demonstrate, it is possible to engineer the environment of the BEC so that the phase shifts associated with self- and cross-phase modulation cancel each other, thus eliminating the density-dependent mean-field shift from the condensate evolution. As the property which governs the evolution of the condensate phase is the chemical potential, the elimination of the mean-field shift is equivalent to having a chemical potential which is independent of the number of atoms in the BEC. In addition, we show that it is possible for at least one branch of the quasiparticle spectrum to be density-independent as well. Measurements which excite only the ‘collisionless’ branch are therefore insensitive to nonlinear phase shifts and hence of considerable interest in atom interferometry.
To illustrate how this works we consider the simplest possible situation that can lead to the required effect, namely a two-component condensate described by the many-body Hamiltonian
$`\widehat{H}`$ $`=`$ $`{\displaystyle \underset{j=a,b}{}}{\displaystyle d^3r\widehat{\psi }_j^{}(𝐫)\left[\frac{\mathrm{}^2}{2m}^2+U_j(𝐫)\right]\widehat{\psi }_j(𝐫)}`$ (1)
$`+`$ $`{\displaystyle \frac{1}{2}}{\displaystyle \underset{j=a,b}{}}\mathrm{}g_j{\displaystyle d^3r\widehat{\psi }_j^{}(𝐫)\widehat{\psi }_j^{}(𝐫)\widehat{\psi }_j(𝐫)\widehat{\psi }_j(𝐫)}`$ (2)
$`+`$ $`\mathrm{}g_x{\displaystyle d^3r\widehat{\psi }_a^{}(𝐫)\widehat{\psi }_b^{}(𝐫)\widehat{\psi }_b(𝐫)\widehat{\psi }_a(𝐫)}.`$ (3)
Here, $`U_j`$ includes both the external trapping potential and the internal atomic energy for component $`j`$, and the constants $`g_a`$, $`g_b`$, and $`g_x`$ give the strengths of the nonlinearities due to atom-atom collisions. In the language of nonlinear optics, $`g_a`$ and $`g_b`$ determine the self-phase modulation of the two components, whereas $`g_x`$ governs cross-phase modulation between them.
The specific physical system that we have in mind consists of a condensate in a one-dimensional atomic waveguide with tight transverse confinement in the $`x`$-$`y`$ plane, i.e., $`U_j(𝐫)=U_j(x,y)`$, while propagation along the $`z`$-dimension is free. The two components are two internal states (e.g., the Zeeman sublevels) of the same atomic species. We assume that they may convert into each other through a weak linear coupling. The ground state wave functions $`\psi _j(𝐫)`$ are hence found by minimizing the Hartree energy functional while holding the total number of atoms fixed (although the individual atom numbers may vary). This linear coupling is either sufficiently weak, or is turned off adiabatically after the steady state has been established, that we do not explicitly include it in the Hamiltonian.
As ideally one would like to have a ‘single-mode’ wave guide, we assume that the transverse confinement is sufficiently strong that all atoms are ‘frozen’ in the ground state $`\phi _j(x,y)`$ of the transverse potential $`U_j(x,y)`$. This leads us to introduce the atomic annihilation operators $`\widehat{\varphi }_j(z,t)`$ for atoms in the transverse ground state as
$$\widehat{\varphi }_j(z,t)=𝑑x𝑑y\phi _j^{}(x,y)\widehat{\psi }_j(𝐫,t).$$
(4)
From the Hamiltonian (3) the Heisenberg equation of motion for $`\widehat{\varphi }_j(z)`$ is found to be
$`{\displaystyle \frac{d}{dt}}\widehat{\varphi }_j(z)`$ $`=`$ $`i\left[{\displaystyle \frac{\mathrm{}}{2m}}{\displaystyle \frac{^2}{z^2}}+E_j\right]\widehat{\varphi }_j(z)`$ (5)
$``$ $`i\left[V_j\widehat{\varphi }_j^{}(z)\widehat{\varphi }_j(z)+V_x\widehat{\varphi }_k^{}(z)\widehat{\varphi }_k(z)\right]\widehat{\varphi }_j(z),`$ (6)
where $`kj`$ and
$`V_j`$ $`=`$ $`g_j{\displaystyle 𝑑x𝑑y|\phi _j(x,y)|^4};j=a,b`$ (7)
$`V_x`$ $`=`$ $`g_x{\displaystyle 𝑑x𝑑y|\phi _a(x,y)|^2|\phi _b(x,y)|^2}.`$ (8)
We choose the energy reference such that $`E_a=\delta /2`$ and $`E_b=\delta /2`$. From these expressions it is immediately apparent that the effective phase modulation constants along the $`z`$-dimension can be modified by appropriately varying the transverse trapping potential $`U_j(x,y)`$.
Assuming that any perturbation present is too weak to excite the transverse degrees of freedom, the ground state wave function and quasiparticle spectrum of low-lying excitations are found by decomposing the boson field operator $`\widehat{\varphi }_j(z,t)`$ around the ground state Hartree wave function as
$$\widehat{\varphi }_j(z,t)=\left[\varphi _j(z)+\delta \widehat{\varphi }_j(z,t)\right]e^{i\omega _0t},$$
(9)
where the perturbation operators $`\delta \widehat{\varphi }_j(z,t)`$ satisfy the boson commutation relations:
$$[\delta \widehat{\varphi }_j(z,t),\delta \widehat{\varphi }_k^{}(z^{},t)]=\delta _{jk}\delta (zz^{}).$$
(10)
Substituting (9) into (6) and keeping only the leading-order terms leads to a time-independent Gross-Pitaevskii equation from which the ground-state wave function $`\varphi _j(z)`$ and chemical potential $`\mathrm{}\omega _0`$ can be determined.
We work in the condensate rest frame and seek plane wave solutions corresponding to a uniform beam of atoms moving along the wave guide. The densities of the two components are then given by
$`\rho _a|\varphi _a|^2`$ $`=`$ $`{\displaystyle \frac{(V_bV_x)\rho \delta }{V_a+V_b2V_x}},`$ (11)
$`\rho _b|\varphi _b|^2`$ $`=`$ $`{\displaystyle \frac{(V_aV_x)\rho +\delta }{V_a+V_b2V_x}},`$ (12)
where we have assumed a fixed total density $`\rho =\rho _a+\rho _b`$. The relative phase between the two components is arbitrary, as it is either random, or fixed by the linear coupling whose strength we have assumed to be negligible. The frequency of phase rotation for the ground state (chemical potential divided by $`\mathrm{}`$) is determined to be
$$\omega _0=\frac{2(V_aV_bV_x^2)\rho +(V_bV_a)\delta }{2(V_a+V_b2V_x)}.$$
(13)
In the case $`V_a=V_b`$, Eqs. (12) and (13) reduce to Eqs. (8) and (12) of Ref. , respectively.
The phase-rotation frequency $`\omega _0`$ of the condensate ground state contains two contributions, one being proportional to the total density $`\rho `$, and the other to the detuning $`\delta `$. Both terms result from the combined effects of cross- and self-modulation. In contrast to the case of a single-component condensate, Eq. (13) shows that the cross and self-phase modulation contributions can conspire to cancel the density-dependent term, provided only that the condition
$$V_x^2=V_aV_b\mathrm{and}V_a+V_b2V_x0$$
(14)
is met. When the solutions given by Eq. (12) give negative densities no homogeneous ground state exists. Thus we must verify whether or not these conditions can be fulfilled for positive densities $`\rho _a`$ and $`\rho _b`$. We consider the most common situation where all two-body interactions are repulsive, $`V_a,V_b,V_x>0`$, and assume without loss of generality that $`V_b>V_x`$, which implies from Eq. (14) that $`V_a<V_x`$. It is then easily shown that the requirement of $`\rho _a>0`$ and $`\rho _b>0`$ yields
$$\frac{V_x}{V_b}<\frac{\delta }{(V_bV_x)\rho }<1,$$
(15)
which can be achieved by an appropriate choice of the detuning $`\delta `$ and/or of the condensate density $`\rho `$.
To show that these conditions can be met in a real system by modifying the trapping potential, we consider the two-component <sup>87</sup>Rb condensate, where the components $`a`$ and $`b`$ correspond to the hyperfine Zeeman sublevels $`|F=2,m_f=1`$ and $`|F=1,m_f=1`$, respectively. The phase modulation constants are in the ratio $`g_a:g_x:g_b=0.97:1.00:1.03`$. We consider the case of harmonic trapping potentials $`U_a=m\omega ^2(x^2+y^2)/2`$ and $`U_b=m\omega ^2[(xx_0)^2+y^2]/2`$, the offset $`x_0`$ in their centers $`x_0`$ being a control variable that can be changed via a bias magnetic field. We have then
$`V_{a,b}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \xi ^2}}g_{a,b}`$ (16)
$`V_x`$ $`=`$ $`{\displaystyle \frac{1}{2\pi \xi ^2}}e^{x_0^2/2\xi ^2}g_x`$ (17)
where $`\xi =\sqrt{\mathrm{}/m\omega }`$ is the extension of the trap ground state. Hence to have condition (14) satisfied, one can choose $`x_0=0.03\xi `$.
When the condition (14) is satisfied, we obtain the remarkable result that the evolution of the condensate phase is independent of density. We emphasize that by itself, this is of course not sufficient to guarantee measurements unperturbed by nonlinear phase shifts can be carried out. Indeed, the detection of any signal relies on some departure of the state of the condensate away from its ground state. Assuming that this change is small, we can describe it in a linearized approach following the Bogoliubov treatment. We proceed by expanding the Hamiltonian (3) with the help of (4) and (9), keeping only quadratic terms in the operators $`\delta \widehat{\varphi }_j`$, yielding
$`\widehat{H}`$ $``$ $`{\displaystyle \underset{j=a,b}{}}{\displaystyle 𝑑z\delta \widehat{\varphi }_j^{}(z)\left[\frac{\mathrm{}^2}{2m}\frac{^2}{z^2}+\mathrm{}V_j|\varphi _j|^2\right]\delta \widehat{\varphi }_j(z)}`$ (18)
$`+`$ $`{\displaystyle \underset{j=a,b}{}}{\displaystyle \frac{\mathrm{}}{2}}V_j{\displaystyle }dz[\varphi _j^2\delta \widehat{\varphi }_j^{}(z)\delta \widehat{\varphi }_j^{}(z)+H.c.]`$ (19)
$`+`$ $`\mathrm{}V_x{\displaystyle }dz[\varphi _a^{}\varphi _b^{}\delta \widehat{\varphi }_a(z)\delta \widehat{\varphi }_b(z)+\varphi _a\varphi _b^{}\delta \widehat{\varphi }_a^{}(z)\delta \widehat{\varphi }_b(z)`$ (21)
$`+H.c.].`$
No first-order contributions in $`\delta \widehat{\varphi }(z)`$ appear in Eq. (21) as a consequence of the fact that the condensate spinor $`\varphi _j(z)`$ satisfies the time-independent nonlinear Schrödinger equation.
The effective Hamiltonian Eq.(21) can be diagonalized via a generalized Bogoliubov transformation. As there is no confining potential in the $`z`$-direction, we expand the operators $`\delta \widehat{\varphi }_j(z)`$ onto plane waves as
$$\delta \widehat{\varphi }_j(z)=(2\pi )^{1/2}𝑑ke^{ikz}\widehat{c}_j(k).$$
(22)
As a consequence of momentum conservation, only the operators $`\widehat{c}_j(k)`$, $`\widehat{c}_j^{}(k)`$, $`\widehat{c}_j(k)`$, and $`\widehat{c}_j^{}(k)`$ are coupled. In order to diagonalize the Hamiltonian (21) we introduce the annihilation operators for quasiparticles with well-defined momentum $`k`$ according to
$$\widehat{b}_\mu (k)=\underset{j}{}\left[u_{\mu j}(k)\widehat{c}_j(k)+v_{\mu j}(k)\widehat{c}_j^{}(k)\right].$$
(23)
Invariance under rotation of the coordinate axes clearly requires that $`u_{\mu j}(k)=u_{\mu j}(k)`$ and $`v_{\mu j}(k)=v_{\mu j}(k)`$. The coefficients $`u_{\mu j}(k)`$ and $`v_{\mu j}(k)`$ are determined by the requirements that the operators $`\widehat{b}_\mu (k)`$ and $`\widehat{b}_\mu ^{}(k)`$ obey the bosonic commutation relations
$$[\widehat{b}_\mu (k),\widehat{b}_\nu (k^{})]=0$$
(24)
and
$$[\widehat{b}_\mu (k),\widehat{b}_\nu ^{}(k^{})]=\delta _{\mu \nu }\delta (kk^{}),$$
(25)
and that the Hamiltonian (21) takes the form
$$\widehat{H}=\underset{\mu }{}𝑑k\mathrm{}\omega _\mu (k)\widehat{b}_\mu ^{}(k)\widehat{b}_\mu (k),$$
(26)
where the $`\omega _\mu (k)`$ are thus the frequencies of the elementary modes for small collective excitations of the condensate.
For a fixed momentum $`k`$ the coefficients $`u_{\mu j}(k)`$ and $`v_{\mu j}(k)`$ may be represented as the matrix elements of the $`2\times 2`$ matrices $`𝐔`$ and $`𝐕`$. In order to satisfy the commutation relations (24) and (25) it is sufficient that
$$\mathrm{𝐔𝐔}^{}\mathrm{𝐕𝐕}^{}=𝐈,\text{ }\mathrm{𝐔𝐕}^T=\mathrm{𝐕𝐔}^T.$$
(27)
The functions $`u_{\mu j}(k)`$ and $`v_{\mu j}(k)`$ can then be determined by substituting Eq.(23) into the commutator $`[b_\mu (k),H]=\mathrm{}\omega _\mu (k)b_\mu (k)`$, which guarantees that the Hamiltonian is of the form (26). This yields:
$$\omega _\mu \sigma _\mu =𝐌\sigma _\mu $$
(28)
where $`\sigma _\mu (u_{\mu a},u_{\mu b},v_{\mu a},v_{\mu b})^T`$ and the matrix $`𝐌`$ is given by
$`𝐌=\left(\begin{array}{cccc}H_a& V_x\varphi _a^{}\varphi _b& V_a(\varphi _a^{})^2& V_x\varphi _a^{}\varphi _b^{}\\ V_x\varphi _a\varphi _b^{}& H_b& V_x\varphi _a^{}\varphi _b^{}& V_b(\varphi _b^{})^2\\ V_a\varphi _a^2& V_x\varphi _a\varphi _b& H_a& V_x\varphi _a\varphi _b^{}\\ V_x\varphi _a\varphi _b& V_b\varphi _b^2& V_x\varphi _a^{}\varphi _b& H_b\end{array}\right)`$
where $`H_j=\mathrm{}k^2/(2m)+V_j\rho _j`$. The eigenfrequencies $`\omega _\mu (k)`$ whose corresponding eigenvectors satisfy Eqs. (27) are
$`\omega _\mu (k)`$ $`=`$ $`\{{\displaystyle \frac{\mathrm{}k^2}{2m}}[({\displaystyle \frac{\mathrm{}k^2}{2m}}+V_a\rho _a+V_b\rho _b)`$ (29)
$`\pm `$ $`[(V_a\rho _aV_b\rho _b)^2+4V_x^2\rho _a\rho _b]^{1/2}]\}^{1/2}.`$ (30)
These solutions therefore yield the two branches of the quasiparticle excitation spectrum for the two-component condensate.
It is straightforward to show that one of the branches does not depend on the density $`\rho `$ when the condition (14) is satisfied. In this case the eigenfrequencies are
$`\omega _+(k)`$ $`=`$ $`\sqrt{{\displaystyle \frac{\mathrm{}k^2}{2m}}\left[{\displaystyle \frac{\mathrm{}k^2}{2m}}+2(V_a\rho _a+V_b\rho _b)\right]}`$ (31)
$`\omega _{}(k)`$ $`=`$ $`{\displaystyle \frac{\mathrm{}k^2}{2m}},`$ (32)
hence $`\omega _{}`$ corresponds to the ‘collisionless’ branch.
The existence of a collisionless branch is related to the invariance of the system under translation along the $`z`$-axis. It is easy to show that the corresponding eigenvector in this branch is given by
$`\sigma _{}(k)=(V_x\varphi _b,V_a\varphi _a,\mathrm{\hspace{0.17em}0},\mathrm{\hspace{0.17em}0})^T`$
which is $`k`$-independent. The condensate wave function is a plane wave, but the choice of inertial frame in which this plane wave is at rest $`(k=0)`$ is arbitrary. The cancellation of self- and cross-phase modulation is achieved by having a uniform density along the $`z`$-axis, with a constant ratio of the components $`a`$ and $`b`$. If a fraction of the total atoms are boosted into a moving frame while still maintaining their internal superposition state, then the relative densities of the two components are not changed. This leads to the possibility of imparting kinetic energy onto the system without affecting the balance required to eliminate density-dependent phase shifts. As the coefficients $`v_j`$ are found to be zero, it is clear that the quasiparticles created by $`b_{}^{}(k)`$ correspond simply to condensate atoms boosted into a different momentum eigenstate. This remarkable property, reflected in the collisionless branch of (32) should be particularly useful for atom interferometry, where the ability to use Bragg-pulses to induce transitions between the condensate state and the collisionless branch could lead to the design of beam-splitters which maintain mean-field-free conditions.
In conclusion, we have shown that for certain values of collisional constants, the nonlinear phase shift in a two-component Bose condensate can be completely eliminated. Such an effect results from the interplay between self- and cross-phase modulation between the two components. While we have explicitly shown how it is possible to adjust the values of these collisional constants in a one-dimensional atomic waveguide, the same method also applies to two-dimensional systems with strong confinement in the third dimension. We expect that these results will play an important role in atom interferometry where uncontrollable nonlinear phase shift limits applications in precision measurement.
Finally, the question remains to determine whether the elimination of the mean-field shifts can be achieved in a more typical 3-d trap as well. In that case, the ground state must be determined from a full three-dimensional time-independent set of Gross-Pitaevskii equations. It will be interesting to see if the chemical potential can be made density-independent in this case by appropriately engineering the trapping potentials of the different components. Because the existence of the collisionless branch of the excitation spectrum appears to be related to translational invariance, future studies will be required to determine whether such a branch exists in the case of a 3-dimensional trap as well.
This work is supported in part by Office of Naval Research Research Contract No. 14-91-J1205, National Science Foundation Grant PHY98-01099, the Army Research Office and the Joint Services Optics Program.
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# On the existence of a variational principle for deterministic cellular automaton models of highway traffic flow
## 1 Introduction
The publication of the Nagel-Schreckenberg highway traffic flow cellular automaton (CA) model has attracted much interest. Since then, many papers describing various CA models of traffic flow have been published . Most results concerning the properties of traffic flow models have been obtained with the help of either numerical simulations or various extensions of the mean-field approximation. Only few exact results are known. In the case of the Fukui-Ishibashi (FI) model , which is a natural extension of Wolfram’s Rule 184 , Fukś has recently derived the exact expression of the average car flow as a function of time. This result shows that, in the infinite lattice size limit, the average car flow is a monotonous increasing function of time. That is, the FI model (and Rule 184 to which it reduces if the speed limit is equal to 1) obeys a variational principle. The purpose of this paper is to find out to what extent such a principle remains valid for other deterministic CA models of traffic flow.
## 2 Rules and configurations representations
In CA models of traffic flow on a circular one-lane highway, the road is represented by a lattice of $`L`$ cells with periodic boundary conditions. Each cell is either empty (in state 0) or occupied by a car (in state 1). Since the number of cars traveling a circular highway is conserved, such a system evolves according to a number-conserving CA rule. We recently established a necessary and sufficient condition for a one-dimensional $`q`$-state ($`q2`$) deterministic CA rule to be number-conserving, and studied a few illustrative examples. These rules may be seen as deterministic evolution rules of one-dimensional closed systems of interacting particles. Among these rules, some are such that all particles always move in the same direction. All deterministic CA models of highway traffic flow are members of this family of rules.
Although most papers on CA models of traffic flow deal with two-state CA, some number-conserving $`q`$-state CA rules, with $`q>2`$, can also be good candidates . In this case, a cell could, for instance, either represent a longer segment of the highway capable of accommodating a maximum of $`q1`$ cars or the unit segment of a $`(q1)`$-lane highway.
In standard CA modeling it is usually assumed that the configuration at time $`t+1`$ is entirely determined by the configuration at time $`t`$. Within this restrictive framework, models like are not standard CA models since the road configuration at time $`t+1`$ depends on road configurations at time $`t`$ and $`t1`$.
As we shall see, the family of potential models of traffic flow in terms of deterministic one-dimensional two-state standard cellular automata is rather rich, and, in this paper, only models of this type will be considered.
In our discussion of traffic rules, we will not represent CA rules by their rule tables, but make use of a representation which clearly exhibits the particle motion. This motion representation, or velocity rule, which has been introduced in , may be defined as follows. If the integer $`r`$ is the rule’s radius, list all the $`(2r+1)`$-neighborhoods of a given particle represented by a 1 located at the central site of the neighborhood. Then, to each neighborhood, associate an integer $`v`$ denoting the velocity of this particle, that is the number of sites this particle will move in one time step, with the convention that $`v`$ is positive if the particle moves to the right and negative if it moves to the left. For example, Rule 184, defined by
$`f_{184}(0,0,0)`$ $`=0,`$ $`f_{184}(0,0,1)`$ $`=0,`$ $`f_{184}(0,1,0)`$ $`=0,`$ $`f_{184}(0,1,1)`$ $`=1,`$
$`f_{184}(1,0,0)`$ $`=1,`$ $`f_{184}(1,0,1)`$ $`=1,`$ $`f_{184}(1,1,0)`$ $`=0,`$ $`f_{184}(1,1,1)`$ $`=1,`$
is represented by the radius-1 velocity rule:
$$\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶}\mathrm{𝟷}.$$
(1)
The symbol $``$ represents either a 0 or a 1, i.e., either an empty or occupied site. This representation, which clearly shows that a car can move to the next-neighboring site on its right if, and only if, this site is empty, is shorter and more explicit.
When discussing road configurations evolving according to various illustrative rules, the knowledge of both car positions and velocities will prove necessary. Therefore, although we are dealing with two-state CA rules, we will not represent the state of a cell by its occupation number (i.e., 0 or 1), but by a letter in the alphabet $`\{𝚎,\mathrm{𝟶},\mathrm{𝟷},\mathrm{},𝚟_{\mathrm{max}}\}`$ indicating that the cell is either empty (i.e., in state $`𝚎`$) or occupied by a car with a velocity equal to $`𝚟`$ (i.e., in state $`𝚟\{\mathrm{𝟶},\mathrm{𝟷},\mathrm{},𝚟_{\mathrm{max}}\}`$). Note that $`𝚟`$ is the velocity with which the car is going to move at the next time step. Configurations of cells of this type will be called velocity configurations or configurations for short.
## 3 The deterministic FI model of traffic flow
The simplest deterministic CA model of traffic flow is the FI model which might be defined as follows. If $`d_i(t)`$ is the distance, at time $`t`$, between car $`i`$ and car $`i+1`$ (cars are moving to the right), velocities are updated in parallel according to the subrule:
$$v_i(t+1)=\mathrm{min}(d_i(t)1,v_{\mathrm{max}})$$
(2)
where $`v_i(t)`$ is the velocity of car $`i`$ at time $`t`$; then cars move according to the subrule:
$$x_i(t+1)=x_i(t)+v_i(t+1),$$
(3)
where $`x_i(t)`$ is the position of car $`i`$ at time $`t`$. The model contains two parameters: the speed limit $`v_{\mathrm{max}}`$, which is the same for all cars, and the car density $`\rho `$.
For the sake of simplicity, in our discussion of this model, it is sufficient to consider the case $`v_{\mathrm{max}}=2`$. The corresponding radius-2 velocity rule is:
$$\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟷}\mathrm{𝟷},\mathrm{𝟷𝟶𝟶}\mathrm{𝟸}.$$
(4)
Figures 1 and 2 show examples of spatiotemporal patterns for car densities $`\rho `$, respectively, equal to 0.26 and 0.50. Empty cells are white while cells occupied by a car with velocity $`v`$ equal to 0, 1, and 2 are, respectively, light grey, dark grey and black.
According to velocity rule (4), and our choice of configurations representation, a cell occupied by a car with velocity $`v`$ must be preceded by, at least, $`v`$ empty cells. Each configuration is, therefore, a concatenation of the following four types of tiles:
| 2 | e | e |
| --- | --- | --- |
| 1 | e |
| --- | --- |
The first tile, which corresponds to cars moving at the speed limit $`v_{\mathrm{max}}=2`$, will be called a perfect tile, the next two tiles, corresponding to cars with a velocity less than 2 (here 1 and 0) will be called defective tiles, and the last tile will be called a free empty cell, that is, an empty cell which is not part of either a perfect or defective tile.
Figure 1 shows that only the first configurations contain defective tiles. After a few time steps, these tiles progressively disappear, and the last configurations contain only perfect tiles and free empty cells. Hence, all cars move at $`v_{\mathrm{max}}=2`$, and the system is said to be in the free-moving phase. In Figure 2, at the beginning, the same process of annihilation of defective tiles takes place but, in this case, all defective tiles do not eventually disappear. A few cars move at $`v_{\mathrm{max}}`$, while other cars have either a reduced speed ($`v=1`$) or are stopped ($`v=0`$). This regime is called the jammed phase. For $`v_{\mathrm{max}}>2`$, FI models exhibit similar qualitative features.
To analyze the annihilation process of defective tiles in FI models, we need to define what we call a local jam.
Definition 1 In deterministic FI models of traffic flow, a local jam is a sequence of defective tiles preceded by a perfect tile and followed by either a perfect tile or free empty cells.
From this definition, it follows that:
Proposition 1 In the case of deterministic FI models of traffic flow, the number of cars which belong to a local jam is a non-increasing function of time.
This result is a direct consequence of the fact that, by definition, a local jam is preceded by a car which is free to move, and, according to whether a new car joins the local jam from behind or not, the number of cars in the local jam remains unchanged or decreases by one unit. Note that the jammed car just behind the free-moving car leading the local jam becomes itself free to move at the next time step.
In order to establish the variational principle, let analyze more precisely how the structure of the most general local jam changes in one time step. A local jam consisting of $`n`$ defective tiles is represented below:
$$\mathrm{}v_1\underset{v_1}{\underset{}{ee\mathrm{}e}}v_2\underset{v_2}{\underset{}{ee\mathrm{}e}}\mathrm{}\mathrm{}v_n\underset{v_n}{\underset{}{ee\mathrm{}e}}v_{\mathrm{max}}\underset{v_{\mathrm{max}}}{\underset{}{ee\mathrm{}e}}\mathrm{},$$
where, for $`i=1,2,\mathrm{},n`$, $`0v_i<v_{\mathrm{max}}`$. At the next time step, a car with velocity $`v_i`$ located in cell $`k`$ moves to cell $`k+v_i`$. Hence, if the local jam is followed by $`v_0`$ free empty cells, where $`v_00`$, we have to distinguish two cases:
(i) If $`v_0+v_1<v_{\mathrm{max}}`$, then the number of jammed cars remains unchanged but the leftmost jammed car, whose velocity was $`v_1`$, is replaced by a jammed car whose velocity is $`v_0+v_1`$:
$`\mathrm{}v_{\mathrm{max}}\stackrel{v_{\mathrm{max}}+v_0}{\stackrel{}{ee\mathrm{}e}}v_1ee\mathrm{}ev_2ee\mathrm{}e\mathrm{}\mathrm{}`$ $`v_nee\mathrm{}ev_{\mathrm{max}}ee\mathrm{}e\mathrm{}`$
$`\mathrm{}e(v_0+v_1)ee\mathrm{}ev_2ee\mathrm{}e\mathrm{}\mathrm{}v_n`$ $`ee\mathrm{}ev_{\mathrm{max}}ee\mathrm{}e\mathrm{}`$
(ii) If $`v_0+v_1v_{\mathrm{max}}`$, then the local jam loses its leftmost jammed car:
$`\mathrm{}ev_1ee\mathrm{}ev_2ee\mathrm{}e\mathrm{}\mathrm{}`$ $`v_nee\mathrm{}ev_{\mathrm{max}}ee\mathrm{}e\mathrm{}`$
$`\mathrm{}v_{\mathrm{max}}\underset{v_{\mathrm{max}}+v_0^{}}{\underset{}{ee\mathrm{}e}}v_2ee\mathrm{}e\mathrm{}\mathrm{}v_n`$ $`ee\mathrm{}ev_{\mathrm{max}}ee\mathrm{}e\mathrm{}`$
and at the next time step, the local jam is followed by $`v_0+v_1v_{\mathrm{max}}=v_0^{}<v_0`$ free empty cells.
If we partition the lattice in tiles sequences whose end points are perfect tiles, then, between two consecutive perfect tiles, either there is a local jam, and the above proof shows that the number of free empty cells between the two perfect tiles cannot increase, or there is no local jam, and the number of free empty cells between the two perfect tiles remains unchanged. We can, therefore, state:
Proposition 2 In the case of deterministic FI models of traffic flow, the number of free empty cells is a non-increasing function of time.
Remark 1 If a configuration contains no perfect tiles, then it does not contain free empty cells. Such a configuration belongs, therefore, to the limit set, and as we shall see below, the system is in its steady state. On a circular highway, if a configuration contains only one perfect tile, the above reasoning applies without modification.
Remark 2 The above proof shows that at each time step, the rightmost jammed car of a local jam moves one site to the left. Local jams can only move backwards.
If $`L`$ is the lattice length, $`N`$ the number of cars, and $`N_{\mathrm{fec}}(t)`$ the number of free empty cells at time $`t`$, we have
$$N_{\mathrm{fec}}(t)=LN\underset{i=1}{\overset{N}{}}v_i(t),$$
since, at time $`t`$, car $`i`$ is necessarily preceded by $`v_i(t)`$ empty cells. Dividing by $`L`$ we obtain
$$\frac{N_{\mathrm{fec}}(t)}{L}=1\frac{N}{L}\frac{N}{L}\frac{1}{N}\underset{i=1}{\overset{N}{}}v_i(t).$$
Hence, for all $`t`$,
$$\rho _{\mathrm{fec}}(t)=1\rho \rho v_t,$$
(5)
where $`\rho _{\mathrm{fec}}(t)`$ is the density of free empty cells at time $`t`$, and $`v_t`$ the average car velocity at time $`t`$. This last result shows that, when time increases, since the density of free empty cells cannot increase, the average car flow $`\rho v_t`$ cannot decrease. Deterministic FI models of highway traffic flow obeys, therefore, the following variational principle:
Proposition 3 In the case of deterministic FI models of traffic flow, for a given car density $`\rho `$, the average car flow is a non-decreasing function of time and reaches its maximum value in the steady state.
The annihilation process of defective tiles stops when there are either no more defective tiles or no more free empty cells. Since a perfect tile consists of $`v_{\mathrm{max}}+1`$ cells, if the car density $`\rho `$ is less than $`\rho _c=1/(v_{\mathrm{max}}+1)`$, there exist enough free empty cells to annihilate all the defective tiles, and all cars become eventually free to move. If $`\rho >\rho _c`$, there are not enough free empty cells to annihilate all defective tiles, and eventually some cars are not free to move at $`v_{\mathrm{max}}`$. The threshold value $`\rho _c`$ is called the critical density. At the end of the annihilation process, all subsequent configurations belong to the limit set, and the system is either said to be in equilibrium or in the steady state.
Note that the existence of a free-moving phase for a car density less than $`\rho _c=1/(v_{\mathrm{max}}+1)`$, can be seen as a consequence of Relation (5). When $`t`$ goes to infinity, according to whether $`\rho _{\mathrm{fec}}`$ is positive or zero, this relation implies
$$v_{\mathrm{}}=\mathrm{min}(\frac{1\rho }{\rho },v_{\mathrm{max}}).$$
If the system is finite, its state becomes eventually periodic in time, and the period is equal to the lattice size or one of its submultiples.
Since local jams move backwards and free empty cells move forwards, equilibrium is reached after a number of time steps proportional to the lattice size.
Remark 3 In the case of Rule 184, the existence of a free-moving phase for a particle density $`\rho `$ less than the critical density $`\rho _c=\frac{1}{2}`$, obviously implies that the average velocity $`v_t`$, as a function of time $`t`$, is maximum when $`t\mathrm{}`$. For $`\rho >\frac{1}{2}`$, this property is still true since the dynamics of the holes (empty sites) is governed by the conjugate of rule Rule 184 (i.e., Rule 226),<sup>1</sup><sup>1</sup>1If $`f`$ is an $`n`$-input two-state deterministic CA rule, its conjugate, denoted $`Cf`$, is defined by
$$Cf(x_1,x_2,\mathrm{},x_n)=1f(1x_1,1x_2,\mathrm{},1x_n).$$
which describes exactly the same dynamics as Rule 184 but for holes moving to the left. Therefore, for all values of the particle density, the average velocity takes its maximum value in the steady state.
In the next section, we shall examine to what extent the variational principle, valid for all deterministic FI traffic flow models, remains valid for more general deterministic CA models of highway traffic flow.
## 4 Variational principle
In order to extend Proposition 3 to more general deterministic CA models of traffic flow, we have first to characterize, among the class of number-conserving deterministic two-state CA rules, which rules may be considered as acceptable CA traffic rules.
### 4.1 Unidirectional motion
The first obvious condition to be satisfied by a traffic rule is that all cars should move in the same direction, that is, in the motion representation, all velocities should have the same sign.
Example 1 For instance, the rule:
$$\mathrm{𝟶𝟷𝟷}\mathrm{𝟷},\mathrm{𝟷𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟶}\mathrm{𝟷}$$
(6)
is not acceptable since a particle can move either to the right or to the left. Many number-conserving CA rules are not unidirectional. In the case of (6), it can be shown that each particle performs a non-Gaussian pseudo-random walk. The rule being deterministic, the randomness comes from the randomness of the initial configuration.
### 4.2 Existence of a free-moving phase
The second natural condition, which should be satisfied by a deterministic CA traffic rule, is that, when the car density is sufficiently low, each car should eventually be able to move at the speed limit $`v_{\mathrm{max}}`$.
Example 2 A numerical simulation of a one-dimensional system of particles evolving according to the rule:
$$\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟷𝟶𝟷}\mathrm{𝟷},\mathrm{𝟶𝟷𝟶}\mathrm{𝟷},\mathrm{𝟷𝟷𝟶𝟶}\mathrm{𝟸}$$
(7)
shows that, for a particle density $`\rho <\frac{1}{2}`$, all particles eventually move at $`v=1`$. For $`\rho >\frac{1}{2}`$, the asymptotic average velocity decreases monotonically and goes to zero for $`\rho =1`$. This rule is, however, not an acceptable traffic rule. The velocity of a particle can be equal to $`v_{\mathrm{max}}=2`$ only if there is another particle located just behind (i.e., on its left) whose velocity is zero. This implies that, while $`v_{\mathrm{max}}=2`$, the average velocity cannot be larger than 1. The regime in which all particles move at the same velocity $`v=1`$ is not a true free-moving phase.
Example 3 It is often necessary to look closely at the velocity rule to tell if a system of particles evolving according to such a rule exhibits a free-moving phase at low density. Consider, for instance, the rule:
$`\mathrm{𝟷}\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟶𝟷𝟷𝟷}\mathrm{𝟶},\mathrm{𝟶𝟶𝟷}\mathrm{𝟶},`$
$`\mathrm{𝟷}\mathrm{𝟷𝟶}\mathrm{𝟷},\mathrm{𝟶𝟷𝟷𝟶}\mathrm{𝟷}.`$ (8)
The rule elements $`\mathrm{𝟷}\mathrm{𝟷𝟶}\mathrm{𝟷}`$ and $`\mathrm{𝟶𝟷𝟷𝟶}\mathrm{𝟷}`$ indicate that, if the particle density is too low, no motion can probably take place. Actually, a simple numerical simulation shows that below $`\rho _{\mathrm{min}}=\frac{1}{3}`$, the asymptotic average velocity is always zero. More precisely, it can be shown that, in the steady state, if $`\rho _v`$ denotes the fraction of all particles whose velocities is equal to $`v`$, there exist three regimes:
(i) If $`\rho \frac{1}{3}`$, then $`\rho _0=\rho `$.
(ii) If $`\frac{1}{3}\rho \frac{1}{2}`$, $`\rho _0`$ and $`\rho _1`$ satisfy the relations:
$$\rho _0+\rho _1=\rho ,\text{and}3\rho _0+2\rho _1=1,$$
that is,
$$\rho _0=12\rho ,\text{and}\rho _1=3\rho 1.$$
(iii) If $`\frac{1}{2}\rho 1`$, $`\rho _0`$ and $`\rho _1`$ satisfy the relations:
$$\rho _0+\rho _1=\rho ,\text{and}\rho _0+2\rho _1=1,$$
that is,
$$\rho _0=2\rho 1,\text{and}\rho _1=1\rho .$$
A system of particles evolving according to (8) does not exhibit a free-moving phase. It is only when $`\rho `$ is exactly equal to $`\frac{1}{2}`$ that all particles move at $`v_{\mathrm{max}}=1`$. As a traffic rule (8) should be discarded.
Example 4 The existence of a free-moving phase for a system of particles evolving according to a given number-conserving deterministic CA rule does not, however, guarantee that such a rule is a reasonable traffic rule. For example, in the case of the velocity rule:
$`\mathrm{𝟷𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟷𝟶𝟷}\mathrm{𝟷},\mathrm{𝟷𝟶𝟷𝟷}\mathrm{𝟷},`$
$`\mathrm{𝟷𝟷𝟶𝟶}\mathrm{𝟸},\mathtt{\hspace{0.17em}1010}\mathrm{𝟸},\mathrm{𝟷𝟶𝟶}\mathrm{𝟸},`$ (9)
all velocities have the same sign, and a simple numerical simulation shows that, below a critical density equal to $`\frac{1}{2}`$, there exists a free-moving phase in which all the particles move with the velocity $`v_{\mathrm{max}}=2`$. Its flow diagram<sup>2</sup><sup>2</sup>2Traffic engineers call fundamental diagram the plot of the asymptotic average flow $`\rho v_{\mathrm{}}`$ versus the density $`\rho `$. has even a nice tent-shape. But, as a traffic rule, (9) shows that drivers are anticipating the motion of cars on their right . Since all drivers move according to the same rule, no collisions occur. But, to make traffic rules somewhat more realistic, most authors consider essential to add some randomization, as random braking, to the basic deterministic model. In the case of rule (9) this would lead to collisions, whose number would increase with the braking probability. Moreover, at the critical density, systems evolving according to a deterministic car traffic rule exhibit a second-order phase transition, and it has been recently shown that, for traffic flows evolving according to FI traffic rules, random braking is the symmetry-breaking field conjugate to the order parameter defined as $`m=v_{\mathrm{max}}v_{\mathrm{}}`$. For all these reasons, we should not accept as a traffic rule any deterministic CA rule in which a particle can move to an occupied site “knowing” that the particle located at that site will also move.
In the light of the preceding examples, we shall adopt the following definition of a free-moving phase
Definition 2 A system of particles evolving according to a unidirectional deterministic CA rule exhibits a free-moving phase if there exists a number $`0<\rho _c<1`$, called the critical density, such that, starting from a random configuration with a particle density $`\rho <\rho _c`$, the system evolves to an equilibrium state in which, with probability one, all configurations consists of perfect tiles and free empty cells. If the maximum velocity at which a particle can move is $`v_{\mathrm{max}}`$, a perfect tile consists of a cell in state $`v_{\mathrm{max}}`$ preceded by $`v_{\mathrm{max}}`$ empty cells.
Example 5 The structure of the last element of the velocity rule:
$$\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟷𝟶}\mathrm{𝟶},\mathrm{𝟶𝟶𝟷𝟶}\mathrm{𝟷},\mathrm{𝟷𝟷𝟶𝟷}\mathrm{𝟷},\mathrm{𝟷𝟷𝟶𝟶}\mathrm{𝟸}$$
(10)
shows that the average velocity can never be equal to $`v_{\mathrm{max}}=2`$. A system of particles evolving according to such a rule cannot, therefore, exhibit a free-moving phase in the sense of the above definition. This type of result is general: If a particle can move at $`v_{\mathrm{max}}`$ if, and only if, the site located immediately behind it has to be occupied by another particle, then no free-moving phase in the sense of Definition 2 can exist.
Example 6 The velocity rule:
$$\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟶𝟷}\mathrm{𝟷},\mathrm{𝟷𝟶𝟶𝟶}\mathrm{𝟸}$$
(11)
describes the behavior of overcautious drivers who avoid occupying a site located just behind another car. The perfect tile:
| 2 | e | e | e |
| --- | --- | --- | --- |
has not the structure required by Definition 2. It contains an extra empty cell. While it could be reasonable to consider models of traffic flows in which some drivers could have an overcautious behavior, if all drivers behave in the same way, then all cars will stop for a density less than the maximum car density $`\rho =1`$.<sup>3</sup><sup>3</sup>3In CA traffic models, a cell can accomodate at most one car. We will, therefore, not consider such a driving strategy acceptable for deterministic CA traffic rules.
### 4.3 Local jams
The existence of a free-moving phase implies the existence of a mechanism making possible the annihilation of all local jams. For deterministic FI models, the jammed car just behind the free-moving car leading the local jam becomes free to move at the next time step. If, at low density, we want local jams to gradually disappear, this condition, which was automatically satisfied in the case of deterministic FI traffic flow models, should be required for models evolving according to unidirectional number-conserving deterministic CA rules to be acceptable traffic rules.
Definition 3 A system of particles evolving according to a unidirectional velocity rule is a deterministic traffic flow model if a sequence of defective tiles cannot be preceded by free empty cells. A sequence of defective tiles preceded by a perfect tile is called a local jam.
Proposition 1 can then be extended to all deterministic traffic flow CA models as defined above.
Proposition 4 In a deterministic CA model of traffic flow, the number of cars belonging to a local jam is a non-increasing function of time.
Example 7 A system of particles evolving according to the velocity rule:
$$\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟶𝟷𝟶𝟷}\mathrm{𝟶},\mathrm{𝟷𝟷𝟶}\mathrm{𝟷},\mathrm{𝟶𝟷𝟶𝟶}\mathrm{𝟸}$$
(12)
is not a traffic flow model. All configurations are concatenations of the following three types of tiles and free empty cells:
| 2 | e | e |
| --- | --- | --- |
| 0 | 1 | e |
| --- | --- | --- |
| 0 | e |
| --- | --- |
The first tile is the perfect tile while the other two tiles are the only defective tiles. Since to move at $`v_{\mathrm{max}}=2`$ a particle needs not only to have two empty sites in front but also one empty site behind, the existence of a defective tile which contains two particles makes that, as shown in the example below, a sequence of defective tiles might be preceded by free empty cells.
time $`t`$: $`\mathrm{}eeee2ee01e0e2eeee\mathrm{}`$
time $`t+1`$: $`\mathrm{}eeeeee01e01eee2ee\mathrm{}`$
To avoid this behavior a defective tile should consist of a cell occupied by one particle with velocity $`v<v_{\mathrm{max}}`$ preceded by $`n_v`$ empty cells, where $`vn_v<v_{\mathrm{max}}`$.
For FI models, the variational principle followed from the particular structure of defective tiles, that is, the number $`n_v`$ of empty cells preceding a cell in state $`v`$ was always equal to $`v`$. All models in which defective tiles have this structure will, therefore, verify Proposition 3, and we can state:
Proposition 5 If a deterministic CA traffic rule is such that local jams have the following structure
$$\mathrm{}v_1\underset{v_1}{\underset{}{ee\mathrm{}e}}v_2\underset{v_2}{\underset{}{ee\mathrm{}e}}\mathrm{}\mathrm{}v_n\underset{v_n}{\underset{}{ee\mathrm{}e}}v_{\mathrm{max}}\underset{v_{\mathrm{max}}}{\underset{}{ee\mathrm{}e}}\mathrm{}$$
where, for all $`i\{1,2,\mathrm{},n\}`$, $`0v_i<v_{\mathrm{max}}`$, then, the average car flow $`\rho v_t`$ is a non-decreasing function of time $`t`$, and reaches its maximum value in the steady state.
Example 8 The velocity rule
$$\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟷𝟶}\mathrm{𝟷},\mathrm{𝟶𝟷𝟶𝟷}\mathrm{𝟷},\mathrm{𝟶𝟷𝟶𝟶}\mathrm{𝟸}$$
(13)
is a nontrivial CA traffic rule satisfying Proposition 5. The speed limit $`v_{\mathrm{max}}`$ is equal to $`2`$. As shown in Figure 3, at low density, all cars move at $`v_{\mathrm{max}}`$. In this regime, the configurations belonging to the limit set consist of perfect three-cell tiles of the type
| 2 | e | e |
| --- | --- | --- |
in a sea of free empty cells. The critical density $`\rho _c`$ is equal to $`\frac{1}{3}`$.
Figure 4 shows a spatiotemporal pattern for a car density higher than the critical density.
All configurations, except the initial configuration chosen at random, are concatenations of perfect tiles, defective tiles, of the following two types:
| 1 | e |
| --- | --- |
and free empty cells. Eventually all free empty cells disappear. In both cases, except for the initial configuration, all local jams are such that Proposition 5 applies.
Example 9 There are many nontrivial traffic rules to which Proposition 5 applies. The following traffic rule represents another example:
$`\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟷}\mathrm{𝟷},\mathrm{𝟶𝟷𝟶𝟶𝟷}\mathrm{𝟸},`$
$`\mathrm{𝟷𝟷𝟶𝟶}\mathrm{𝟸},\mathrm{𝟶𝟷𝟶𝟶𝟶}\mathrm{𝟹}`$ (14)
Remark 4 If, as in Example 3, we denote by $`\rho _v`$ the fraction of all particles whose velocities are equal to $`v`$ in the steady state, for all velocity rules to which Proposition 5 applies, we have
$`{\displaystyle \underset{v=0}{\overset{v_{\mathrm{max}}}{}}}\rho _v`$ $`=\rho `$ (15)
$`{\displaystyle \underset{v=0}{\overset{v_{\mathrm{max}}}{}}}(v+1)\rho _v`$ $`=1.`$ (16)
Subtracting (15) from (16), we find that the asymptotic average velocity is given by
$$v_{\mathrm{}}=\frac{1}{\rho }\underset{v=0}{\overset{v_{\mathrm{max}}}{}}v\rho _v=\frac{1\rho }{\rho }.$$
The last expression is exactly the result one obtains using a mean-field argument. The fact that mean-field arguments lead to exact results for CA traffic rules to which Proposition 5 applies has already been found for the values of the critical exponents of phase transitions in FI traffic flow models . It might be conjectured that, for all CA traffic rules to which Proposition 5 applies, critical exponents will be equal to their mean-field values.
Example 10 Proposition 5 does not apply to a system of particles evolving according to the rule:
$`\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟶}\mathrm{𝟸}`$ (17)
since all jammed cars have a zero velocity even when there exists a free empty cell in front. All configurations are concatenations of the following three types of tiles:
| 2 | e | e |
| --- | --- | --- |
| 0 | e |
| --- | --- |
and free empty cells. The first tile, which corresponds to cars moving at the speed limit $`v_{\mathrm{max}}=2`$, is the perfect tile, and the other two tiles, corresponding to cars with a zero velocity are the only defective tiles. It is easy to verify, however, that Proposition 4 apply. Since all jammed cars are stopped cars, it follows that, for a given car density, the average velocity $`v_t`$ is a monotonous non-increasing function of time.
Example 11 For systems evolving according to the rule:
$`\mathrm{𝟷𝟷}\mathrm{𝟶},\mathrm{𝟷𝟶𝟷}\mathrm{𝟶},\mathrm{𝟶𝟷𝟶𝟶𝟷}\mathrm{𝟷},`$
$`\mathrm{𝟷𝟷𝟶𝟶𝟷}\mathrm{𝟸},\mathrm{𝟷𝟶𝟶𝟶}\mathrm{𝟹},`$ (18)
due to the existence of defective tiles of the following structures:
| 1 | e | e |
| --- | --- | --- |
| 0 | e |
| --- | --- |
Proposition 5 does not apply, and while Proposition 4 applies, for a given car density, the average velocity $`v_t`$ is not a monotonous non-increasing functions of time, since, as shown below, the velocity of a jammed car, equal to 1 at time $`t`$, may become equal to 0 at time $`t+1`$:
time $`t`$: $`\mathrm{}3eee1ee3eeeeee\mathrm{}`$
time $`t+1`$: $`\mathrm{}ee0e3eeee3eee\mathrm{}.`$
This example shows that, while a deterministic CA rule may be an acceptable deterministic CA traffic rule, the variational principle, which in its stronger form states that, for a given car density, the average flow is a monotonous non-decreasing function of time, does not apply. The question which could be addressed is then: Is, however, the variational principle still valid under a weaker form?
Since we assumed that a system evolving according to a deterministic CA rule should exhibit a free-moving phase at low car density, by definition of the free-moving phase, all cars move at $`v_{\mathrm{max}}`$. Therefore, for $`\rho <\rho _c`$, the average car flow takes its maximum value in the steady state. Above the critical density, it is only if the asymptotic average flow is a well-defined function of car density that the question makes sense. If this is the case, then, since Proposition 4 applies to all traffic flow deterministic CA models, it may be reasonably conjectured that, also for $`\rho >\rho _c`$, the average car flow takes its maximum value in the steady state.
Remark 5 This paper deals with the existence of a variational principle for number-conserving deterministic CA models of highway traffic flow. The variational principle, as stated in Proposition 5, may be valid, however, for one-dimensional closed systems of particles evolving according to more general number-conserving deterministic CA rules. For instance, in the case of Rule (7), presented Example 2, all configurations are concatenations of the following three tiles:
| 0 | 2 | e | e |
| --- | --- | --- | --- |
| 1 | e |
| --- | --- |
and free empty cells. The tiles’ structure shows that Relation (5) is valid. If the particle density $`\rho `$ is less than $`\frac{1}{3}`$, there are enough free empty cells to annihilate all cells of the first and third type; and all configurations of the limit set are concatenations of tiles of the second type and free empty cells. If $`\rho >\frac{1}{3}`$, there are not enough free empty cells to annihilate all cells of the first and third type; and all configurations of the limit set contain tiles of all types but no free empty cells. Although such a system does not exhibit a free-moving phase according to Definition 2, cells of the second type play the role of “perfect tiles” while the other two types of cells can be regarded as “defective tiles”.
When Relation (5) is satisfied, it is necessary to verify that the density of free empty cells is a non-increasing function of time to ensure the validity of the variational principle under its stronger form. For instance, while $`\rho _{\mathrm{fec}}(t)`$ cannot increase with time $`t`$ for a system of particles evolving according to Rule (7), this is no more the case for a system of particles evolving according Rule (10), presented in Example 5, and which is obtained by replacing the element $`\mathrm{𝟶𝟷𝟶}\mathrm{𝟷}`$ of Rule (7) by the two elements $`\mathrm{𝟷𝟶𝟷𝟶}\mathrm{𝟶}`$ and $`\mathrm{𝟶𝟶𝟷𝟶}\mathrm{𝟷}`$. As for Rule (7), all configurations are concatenations of the same above three tiles and free empty cells, which implies that Relation (5) is satisfied. But, due the structure of the element $`\mathrm{𝟷𝟶𝟷𝟶}\mathrm{𝟶}`$ of Rule (10), as shown below (three-particle system evolving on a size-10 lattice):
time $`t`$ $`\underset{¯}{ee}02ee1e\underset{¯}{ee}`$
time $`t+1`$ $`\underset{¯}{ee}1e\underset{¯}{e}1e0\underset{¯}{ee}`$
the number of free empty cells ($`\underset{¯}{e}`$) increases from 4 at time $`t`$ to 5 at time $`t+1`$ resulting in a decrease of the average velocity.
## 5 Conclusion
As illustrated by a number of different examples, many number-conserving deterministic CA rules cannot be considered reasonable deterministic traffic rules, and we tried to list the essential properties characterizing a traffic rule in order to give a general definition of deterministic CA models of highway traffic flow. Various notions such as free-moving phase, perfect and defective tiles, and local jam play an important role in our discussion. We have then shown that, within the framework of this definition, a variety of deterministic CA models of traffic flow obey a variational principle which, in its stronger form, states that, for a given car density, the average car flow is a non-decreasing function of time. This result has been established for traffic flow models whose configurations exhibits local jams of a given structure. If local jams have a different structure, while this variational principle may still apply to systems evolving according to some particular rules, it will not apply in general. However, if the asymptotic average car flow is a well-defined function of car density, since we have proved that for all traffic flow models the number of jammed cars of a local jam cannot increase, we conjectured that it will apply under a weaker form which states that, for a given car density, the average car flow takes its maximum value in the steady state.
The variational principle also applies, even in its stronger form, to many number-conserving deterministic CA rules which cannot be considered reasonable traffic rules.
Acknowledgements
The author is grateful to Henryk Fukś and Andrés Moreira for their very good suggestions. This work has been done during a stay at the Centro de Modelamiento Matemático de la Universidad de Chile in Santiago thanks to FONDAP-CONICYT. The unflagging interest of Eric Goles has been of great help.
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# Nonequilibrium Quantum Dynamics of Second Order Phase Transitions
## I Introduction
A system can interact directly with an environment to make its coupling parameters depend explicitly on time. Even the effective coupling parameters of a subsystem of a closed system, though conserved as a whole, may depend implicitly on time through an interaction with the rest of the system. The characteristic feature of these open systems is that their effective coupling parameters depend explicitly or implicitly on time. Therefore the genuine understanding of these systems requires the real-time processes from their initial conditions. Of a particular interest are the systems undergoing phase transitions. When a system cools down through an interaction with an environment, it may undergo a phase transition and its coupling parameters depend explicitly on time. Similarly, matter fields in the expanding early Universe undergo phase transitions through the interaction with gravity.
Phase transitions are one of the most physically important phenomena in nature and have wide applications from condensed matter physics, particle physics and even to cosmology. The Kibble mechanism explains formation of topological defects in symmetry breaking phase transitions . The kinetic process by Zurek has revealed a new feature of symmetry breaking phase transitions . The dynamics of phase transitions and formation of topological defects since then have become an important tool in variety of phenomena in condensed matter physics . It is also widely accepted that symmetry breaking phase transitions in the early Universe are inevitable for the structure formation of the present Universe . There has been an attempt through laboratory experiments to investigate the process of structure formation in the early stage of Universe . In QCD the quark-antiquark condensate breaks chiral symmetry when the temperature of quark-gluon plasma is lowered .
However, the most difficult and subtle facet of phase transitions is to understand its dynamics and the formation process of topological defects. As emphasized above, systems do become out of equilibrium in general during phase transitions because their coupling parameters depend explicitly on time through the interaction with an environment (heat bath). Hence the nonequilibrium dynamics of phase transitions should differ significantly from the equilibrium dynamics. In finite temperature field theory one obtains the effective action for the system in a thermal equilibrium or quasi-equilibrium by calculating its quantum fluctuations about a vacuum . But as the phase transition proceeds, the fluctuations grow and the stability is lost. The effective action, when extrapolated to the phase transition regime in a literal sense, has a complex value, the imaginary part of which is related with the decay rate of the false vacuum . In this sense finite temperature field theory can not be directly applied to study symmetry breaking phase transitions.
Schwinger and Keldysh introduced the closed time-path integral to treat properly the quantum evolution of quantum fields out of equilibrium from their initial thermal equilibrium . Since then the closed time-path integral method has been developed and applied to many related problems . Recently the closed time-path integral method has been employed to study the nonequilibrium dynamics of second order phase transitions . Another method is the functional Schrödinger-picture approach, in which the evolution of quantum states is found for explicitly time-dependent Hamiltonians . The large $`N`$-expansion method and the mean-field or Hartree-Fock method are used in conjunction with either the Schwinger-Keldish or functional Schrödinger method. Still another methods are the time-dependent variational principle , the generating function for correlation functions , and thermal field theory .
The purpose of this paper is two-fold. In the first part of this paper we elaborate and establish the recently introduced Liouville-von Neumann (LvN) approach so that it can readily be applied to nonequilibrium dynamics. In the second part we apply the approach to the systems undergoing time-dependent second order phase transitions. The LvN approach is a canonical method that unifies the functional Schrödinger equation for the quantum evolution of pure states and the LvN equation for the quantum description of mixed states of either equilibrium or nonequilibrium. One of the advantages is that one can make use of the well-known techniques of quantum mechanics and quantum many-particle systems. It is based on the observation by Lewis and Riesenfeld that the quantum LvN equation, which is originally used to define the density operator for a mixed state, can also be used to find the exact pure states of a time-dependent harmonic oscillator. This observation makes it possible to find not only the mixed state but also the pure state of a time-dependent system. This LvN approach has been developed and applied to quantum fields in an expanding Friedmann-Robertson-Walker universe and to open boson and fermion systems .
In this paper we particularly focus on the model systems whose coupling parameters change signs during the evolution and emphasize the role these systems playing in the second order phase transition. In the case of a time-dependent harmonic oscillator or an ensemble of such oscillators with a positive time-dependent frequency squared, the Fock space consists of the number states which are the exact quantum states of the Schrödinger equation . The density operator is constructed in terms of the annihilation and creation operators that satisfy the LvN equation . Hence the coherent, thermal and coherent-thermal states are found rather straightforwardly according to the standard technique of quantum mechanics. We further show that the same construction of the Fock space and density operator still holds for the time-dependent harmonic oscillator with a sign changing frequency squared. By studying some analytically solvable models we investigate how the instability grows. Another technical strong point of the LvN approach is that it can also be used for a time-dependent anharmonic oscillator. At the leading order the LvN approach is equipped with the time-dependent annihilation and creation operators, the vacuum state of which is already the time-dependent Gaussian state that minimizes the Dirac action . The LvN approach thus provides one with a nonperturbative quantum description for the time-dependent anharmonic oscillator, too. We find the coherent, thermal and coherent-thermal states for the anharmonic oscillator with a polynomial potential and study the dynamics from their effective actions.
As field models for the second order phase transition, we consider first a free massive and then a self-interacting scalar field, the mass of which changes the sign in a finite time period through an external interaction. By studying analytically the exactly solvable model of the free scalar field, we show how the instability of long wavelength modes grows in time. The two-point thermal correlation function is expressed in terms of the classical solution for each mode that is already found in terms of a well-known function. The domain sizes are evaluated analytically by the steepest descent method and are shown to grow as $`t^{1/4}`$ during the quench and as the Cahn-Allen scaling relation $`t^{1/2}`$ after the completion of quench. Remarkably, the Cahn-Allen scaling relation shows a time-lag given by the cubic power of the quench period, which is absent in the instantaneous quench model . However, both the instantaneous and the finite quench models have the same scaling relation given by the classical Cahn-Allen equation , confirming the result of the instantaneous quench model . The free scalar field model describes only the stage of spinodal instability from the unstable local maximum to the spinodal line. In the self-interacting field model, the back-reaction of an interacting term contributes positively to the frequencies of long wavelength modes and shuts off the exponential growth of instability after crossing the spinodal line at sufficiently later times.
The organization of this paper is as follows. It mainly consists of two parts: in the first part from Sec. II to Sec. V the LvN approach is elaborated to be applicable to phase transitions, and in the second part of Secs. VI and VII it is then applied to the second order phase transitions. In Sec. II we review the LvN approach to time-dependent quantum systems. In Sec. III we study the time-dependent harmonic oscillator and find the exact Fock space and the density operator. The density matrix is found and the nonequilibrium quantum dynamics is studied for the coherent, thermal and coherent-thermal states. In Sec. IV the LvN approach is applied to time-dependent inverted oscillators as a toy model for the second order phase transition. In Sec. V we extend the LvN approach to time-dependent anharmonic oscillators with polynomial potentials. The LvN approach leads to the nonperturbative Gaussian state at the leading order. We also show that the coherent state of the LvN approach recovers the result from the time-dependent mean-field or Hartree-Fock method. In Sec. VI we study an exactly solvable model of a free scalar field, which has a time-dependent mass coupling parameter and undergoes smoothly the second order phase transition for a finite quench period. The two-point thermal correlation function is analytically evaluated during the quench and after the completion of quench, and the scaling relations for the domain size are found. In Sec. VII we study a self-interacting scalar field that undergoes the time-dependent second order phase transition.
## II Liouville-von Neumann (LvN) Approach
In this section we briefly review but emphasize the underlying assumptions of the LvN approach to time-dependent quantum systems introduced in Ref. . A time-dependent system can not remain in the initial equilibrium or the instantaneous quasi-equilibrium because the density operator
$$\rho _\mathrm{H}=\frac{1}{Z_\mathrm{H}(t)}e^{\beta \widehat{H}(t)},$$
(1)
where $`\widehat{H}(t)`$ is a time-dependent Hamiltonian operator and $`Z_\mathrm{H}`$ is the partition function, does not satisfy the quantum LvN equation. This means that though the system starts in the initial thermal equilibrium, its final state can be far away from the initial one. Even an initial pure state evolves toward a final one which differs drastically from the initial one, and leads, for instance, to particle production . The Fock or Hilbert space of such time-dependent system of a quantum field transforms unitarily inequivalently so that even the initial vacuum state evolves to a superposition of particle states at final times. Thus the nonequilibrium system described by time-dependent quantum Hamiltonian follows the evolution of a mixed state that is out of equilibrium and is characterized by time-dependent processes. To properly describe the nonequilibrium evolution we shall adopt two assumptions.
First, from a microscopic point of view we assume that the nonequilibrium system obeys the quantum law, i.e. the time-dependent Schrödinger or Tomonaga-Schwinger equation
$$i\mathrm{}\frac{}{t}|\mathrm{\Psi }(t)=\widehat{H}(t)|\mathrm{\Psi }(t).$$
(2)
Here $`\widehat{H}(t)`$ is the time-dependent Hamiltonian of the system. This assumption is physically well-grounded since all the individual constituents of the system should obey the quantum law and the system as a whole should still obey the quantum law provided that all the interactions among the individuals are properly taken into account.
Second, from a statistical point of view it is assumed that even the nonequilibrium system obeys the time-dependent quantum Liouville-von Neumann (LvN) equation
$$i\mathrm{}\frac{}{t}\widehat{\rho }(t)+[\widehat{\rho }(t),\widehat{H}(t)]=0.$$
(3)
The LvN equation has been used to find the density operator for equilibrium systems that are stationary. Now the nonequilibrium system with their time-dependent coupling parameters follows the same equation, so the density operator (1) that is directly defined in terms of the Hamiltonian itself does not necessarily satisfy the equation. In Sec. III we shall see how much the density operator satisfying Eq. (3) deviates from the instantaneous density operator (1) for time-dependent harmonic oscillators.
In the context of quantum mechanics a powerful canonical method was developed by Lewis and Riesenfeld . They observed that any operator $`\widehat{𝒪}(t)`$ satisfying the quantum LvN equation
$$i\mathrm{}\frac{}{t}\widehat{𝒪}(t)+[\widehat{𝒪}(t),\widehat{H}(t)]=0,$$
(4)
can also be used to find the exact quantum states of Eq. (2). In fact, any exact quantum state is given by
$$|\mathrm{\Psi }(t)=\underset{n}{}c_ne^{\frac{i}{\mathrm{}}\gamma _n(t)}|\lambda _n,t.$$
(5)
where
$`\widehat{𝒪}(t)|\lambda _n,t=\lambda _n|\lambda _n,t,`$ (6)
$`\gamma _n(t)={\displaystyle 𝑑t\lambda _n,t|\left(i\mathrm{}\frac{}{t}\widehat{H}(t)\right)|\lambda _n,t}.`$ (7)
Another useful property of LvN approach is the linearity of the LvN equation. In fact, a product $`\widehat{𝒪}_1(t)\widehat{𝒪}_2(t)`$ satisfies Eq. (4) whenever $`\widehat{𝒪}_1(t)`$ and $`\widehat{𝒪}_2(t)`$ satisfy Eq. (4). Therefore it holds that any analytic functional $`F[𝒪(t)]`$ satisfies Eq. (4) provided that $`𝒪(t)`$ satisfies the same equation. In particular, we can still use $`\widehat{𝒪}(t)`$ to define the density operator for the time-dependent system
$$\widehat{\rho }_𝒪(t)=\frac{e^{\beta \widehat{𝒪}(t)}}{Z_𝒪(t)},Z_𝒪(t)=\mathrm{Tr}[e^{\beta \widehat{𝒪}(t)}].$$
(8)
Here $`\beta `$ is a free parameter and will be identified with the inverse temperature only for the equilibrium system, in which the Hamiltonian itself satisfies Eq. (4) and is used to define the standard density operator (1).
## III Time-Dependent Oscillator
As a simple nonequilibrium system we begin with an ensemble of time-dependent harmonic oscillators. This system carries more meaning than being merely a toy model because most of systems with some exception such as the massless $`\mathrm{\Phi }^4`$-theory, either in equilibrium or in nonequilibrium, can be approximated by a quadratic Hamiltonian around zero-force points, stable or unstable. In the technical aspect the time-dependent oscillator can be exactly solved in terms of its classical solution. We now focus on a general oscillator with a time-dependent mass and frequency squared
$$\widehat{H}(t)=\frac{1}{2m(t)}\widehat{p}^2+\frac{m(t)}{2}\omega ^2(t)\widehat{q}^2,$$
(9)
where $`\omega ^2(t)`$ is allowed to change the sign during a phase transition. The LvN approach will be employed below to find the exact Fock space and to construct the coherent, thermal and coherent-thermal states.
### A Fock Space
The key idea of the LvN approach to the Hamiltonian (9) is to require a pair of operators
$`\widehat{a}(t)`$ $`=`$ $`i\left(u^{}(t)\widehat{p}m(t)\dot{u}^{}(t)\widehat{q}\right),`$ (10)
$`\widehat{a}^{}(t)`$ $`=`$ $`i\left(u(t)\widehat{p}m(t)\dot{u}(t)\widehat{q}\right),`$ (11)
to satisfy the LvN equation (4). This results in the classical equation of motion for the complex auxiliary variable $`u`$
$$\ddot{u}(t)+\frac{\dot{m}(t)}{m(t)}\dot{u}(t)+\omega ^2(t)u(t)=0.$$
(12)
Note that these operators depend explicitly on time through $`u(t)`$ and are hermitian conjugate to each other. Further, these operators can be made the annihilation and creation operators with the standard commutation relation for all the times
$$[\widehat{a}(t),\widehat{a}^{}(t)]=1.$$
(13)
The above commutation relation is guaranteed by the wronskian condition
$$\mathrm{}m(t)\left(\dot{u}^{}(t)u(t)\dot{u}(t)u^{}(t)\right)=i.$$
(14)
A comment is in order. Nothing prevents one from using these operators for an inverted time-dependent oscillator as far as Eq. (14) is satisfied. The inverted oscillator will be treated in detail in Sec. IV in the context of second order phase transitions.
¿From the argument in Sec. II, using $`\widehat{a}(t)`$ and $`\widehat{a}^{}(t)`$ one may construct two particular operators that also satisfy Eq. (4): the number and the density operator. By defining the number operator in the usual way
$$\widehat{N}(t)=\widehat{a}^{}(t)\widehat{a}(t),$$
(15)
one finds the Fock space consisting of the time-dependent number states
$$\widehat{N}(t)|n,t=n|n,t.$$
(16)
The vacuum state is the one that is annihilated by $`\widehat{a}(t)`$ and the $`n`$th number state is obtained by applying $`\widehat{a}^{}(t)`$ $`n`$-times to the vacuum state:
$`\widehat{a}(t)|0,t=0,`$ (17)
$`|n,t={\displaystyle \frac{(\widehat{a}^{}(t))^n}{\sqrt{n!}}}|0,t.`$ (18)
In the coordinate representation the number state is given by (see Appendix A)
$$\mathrm{\Psi }_n(q,t)=\left(\frac{1}{2\pi \mathrm{}^2u^{}(t)u(t)}\right)^{1/4}\frac{1}{\sqrt{2^nn!}}\left(\frac{u(t)}{u^{}(t)}\right)^nH_n\left(\frac{q}{\sqrt{2\mathrm{}^2u^{}(t)u(t)}}\right)\mathrm{exp}\left[\frac{i}{2}\frac{m}{\mathrm{}}\frac{\dot{u}^{}(t)}{u^{}(t)}q^2\right],$$
(19)
where $`H_n`$ is the Hermite polynomial.
Equation (11) can be inverted to yield the position and momentum operators
$`\widehat{q}=\mathrm{}\left(u(t)\widehat{a}(t)+u^{}(t)\widehat{a}^{}\right),`$ (20)
$`\widehat{p}=\mathrm{}m(t)\left(\dot{u}(t)\widehat{a}(t)+\dot{u}^{}(t)\widehat{a}^{}\right).`$ (21)
Hence the expectation value of the position and momentum with respect to each number state vanishes
$$n,t|\widehat{q}|n,t=n,t|\widehat{p}|n,t=0.$$
(22)
The only nonvanishing expectation values come from even powers of the position or momentum. The quadratic power of the position and momentum has the expectation values
$`n,t|\widehat{q}^2|n,t`$ $`=`$ $`\mathrm{}^2u^{}(t)u(t)(2n+1),`$ (23)
$`n,t|\widehat{p}^2|n,t`$ $`=`$ $`\mathrm{}^2m^2(t)\dot{u}^{}(t)\dot{u}(t)(2n+1),`$ (24)
$`n,t|(\widehat{q}\widehat{p}+\widehat{p}\widehat{q})|n,t`$ $`=`$ $`\mathrm{}^2m(t)\left(\dot{u}^{}(t)u(t)+u^{}(t)\dot{u}(t)\right)(2n+1).`$ (25)
The Hamiltonian thus has the expectation value
$$H_\mathrm{n}(t)=n,t|\widehat{H}(t)|n,t=\frac{\mathrm{}^2}{2}m(t)\left[\dot{u}^{}(t)\dot{u}(t)+\omega ^2(t)u^{}(t)u(t)\right](2n+1).$$
(26)
The prominent advantages of using $`\widehat{a}^{}(t)`$ and $`\widehat{a}(t)`$ appear when one tries, by using the standard technique of quantum mechanics, to construct other quantum states such as the coherent, thermal and coherent-thermal states.
### B Coherent State
The coherent state is a particularly useful quantum state in quantum field theory for phase transitions or nonequilibrium dynamics. It may be used to obtain the effective potential for the classical background as on order parameter with contributions from quantum fluctuations. By treating $`\widehat{a}(t)`$ and $`\widehat{a}^{}(t)`$ as the annihilation and creation operators, we follow the definition of the coherent state for the time-independent oscillator . The coherent state is defined as an eigenstate of $`\widehat{a}(t)`$:
$$\widehat{a}(t)|\alpha ,t=\alpha |\alpha ,t,$$
(27)
where $`\alpha `$ is a complex constant. The coherent state can be treated algebraically by introducing a displacement operator
$$\widehat{D}(\alpha )=e^{\alpha \widehat{a}^{}(t)+\alpha ^{}\widehat{a}(t)},$$
(28)
which is unitary
$$\widehat{D}(\alpha )\widehat{D}^{}(\alpha )=\widehat{D}^{}(\alpha )\widehat{D}(\alpha )=\widehat{I}.$$
(29)
Moreover, the displacement operator translates by constants the annihilation and creation operators through the unitary transformation
$$\widehat{D}(\alpha )\widehat{a}(t)\widehat{D}^{}(\alpha )=\widehat{a}(t)+\alpha ,\widehat{D}(\alpha )\widehat{a}^{}(t)\widehat{D}^{}(\alpha )=\widehat{a}(t)+\alpha ^{},$$
(30)
and the inverse unitary transformation
$$\widehat{D}^{}(\alpha )\widehat{a}(t)\widehat{D}(\alpha )=\widehat{a}(t)\alpha ,\widehat{D}^{}(\alpha )\widehat{a}^{}(t)\widehat{D}(\alpha )=\widehat{a}(t)\alpha ^{}.$$
(31)
Hence one sees from Eq. (30) that the coherent state results from the unitary transformation of the vacuum state
$`|\alpha ,t`$ $`=`$ $`\widehat{D}^{}(\alpha )|0,t`$ (32)
$`=`$ $`e^{\frac{\alpha ^{}\alpha }{2}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}{\displaystyle \frac{\alpha ^n}{\sqrt{n!}}}|n,t.`$ (33)
It has been known since Schrödinger that the coherent state gives rise to a classical field for the time-independent oscillator . In our time-dependent case, from the expectation values of the position and momentum with respect to the coherent state
$`q_c(t)=\alpha ,t|\widehat{q}|\alpha ,t`$ $`=`$ $`\mathrm{}\left(\alpha u(t)+\alpha ^{}u^{}(t)\right),`$ (34)
$`p_c(t)=\alpha ,t|\widehat{p}|\alpha ,t`$ $`=`$ $`\mathrm{}m(t)\left(\alpha \dot{u}(t)+\alpha ^{}\dot{u}^{}(t)\right),`$ (35)
one sees that $`q_c`$ indeed satisfies the classical equation of motion (12), because $`u(t)`$ has already satisfied Eq. (12) and $`\alpha `$ is a constant. Besides, $`q_c`$ is real and $`p_c=m(t)\dot{q}_c`$, so one may identify $`q_c`$ and $`p_c`$ with the classical position and momentum.
In the above the coherent state has been constructed from the exact Fock space. There is another method to find the coherent state based on the minimization of action . In contrast with the LvN approach in which the operators (11) are required to satisfy the LvN equation, one regards $`u`$ as a free parameter, works on the $`u`$-parameter Fock space and minimizes the Hamiltonian expectation value with respect to the coherent state. To show the method in detail, take the Hamiltonian expectation value with respect to the coherent state
$$H_\mathrm{C}(t)=\alpha ,t|\widehat{H}(t)|\alpha ,t=H_c(t)+H_\mathrm{V}(t),$$
(36)
which consists of the classical part
$$H_c(t)=\frac{1}{2m(t)}p_c^2+\frac{m(t)}{2}\omega ^2(t)q_c^2,$$
(37)
and the vacuum fluctuation part given in Eq. (26) with $`n=0`$
$$H_\mathrm{V}(t)=\frac{\mathrm{}^2}{2}m(t)\left[\dot{u}^{}(t)\dot{u}(t)+\omega ^2(t)u^{}(t)u(t)\right].$$
(38)
By writing the complex $`u`$ in a polar form
$$u(t)=\frac{\xi (t)}{\sqrt{\mathrm{}}}e^{i\theta (t)},$$
(39)
in terms of which Eq. (14) becomes $`\dot{\theta }=1/(2m\xi ^2)`$, and by introducing $`p_\xi =m(t)\dot{\xi }`$, one obtains the effective Hamiltonian
$$H_\mathrm{C}(t)=\frac{1}{2m(t)}p_c^2+\frac{m(t)}{2}\omega ^2(t)q_c^2+\mathrm{}\left[\frac{1}{2m(t)}p_\xi ^2+\frac{m(t)}{2}\omega ^2(t)\xi ^2+\frac{1}{8m(t)\xi ^2}\right].$$
(40)
The last term in the square bracket of Eq. (40) has the same form as the angular momentum for a particle having rotational symmetry in two dimensions, but its origin is rooted on the condition (14) from quantization. Thus the effective Hamiltonian from the coherent state is equivalent to a two-dimensional Hamiltonian, which consists of the classical and quantum fluctuation parts. The variables $`q_c`$ and $`\xi `$ are independent and the Hamilton equations are for $`q_c`$
$`{\displaystyle \frac{dq_c}{dt}}`$ $`=`$ $`{\displaystyle \frac{H_\mathrm{C}}{p_c}}={\displaystyle \frac{1}{m(t)}}p_c,`$ (41)
$`{\displaystyle \frac{dp_c}{dt}}`$ $`=`$ $`{\displaystyle \frac{H_\mathrm{C}}{q_c}}=m(t)\omega ^2(t)q_c,`$ (42)
and for $`\xi `$
$`{\displaystyle \frac{d\xi }{dt}}`$ $`=`$ $`{\displaystyle \frac{(H_\mathrm{C}/\mathrm{})}{p_\xi }}={\displaystyle \frac{1}{m(t)}}p_\xi ,`$ (43)
$`{\displaystyle \frac{dp_\xi }{dt}}`$ $`=`$ $`{\displaystyle \frac{(H_\mathrm{C}/\mathrm{})}{\xi }}=m(t)\omega ^2(t)\xi +{\displaystyle \frac{1}{4m(t)\xi ^3}}.`$ (44)
It is then easy to show that the Hamilton equations (44) equal to the second order equation
$$\ddot{\xi }(t)+\frac{\dot{m}(t)}{m(t)}\dot{\xi }(t)+\omega ^2(t)\xi \frac{1}{4m^2(t)\xi ^3}=0,$$
(45)
and that Eq. (45) is nothing but the equation (12) when $`u`$ has the form (39) and satisfies the condition (14). Hence the minimization of the effective action gives the identical result as the LvN approach. Still another method is the mean field approach, in which the position and momentum are divided into a classical background and a fluctuation part
$$q=q_c+q_f,p=p_c+p_f.$$
(46)
Then the total Hamiltonian is composed of three parts: the classical background and fluctuation parts
$$H(t)=H_c(t)+H_f(t)+H_{int}(t),$$
(47)
where
$$H_f(t)=\frac{1}{2m(t)}p_f^2+\frac{m(t)}{2}\omega ^2(t)q_f^2$$
(48)
is the fluctuation Hamiltonian, and
$$H_{int}(t)=\frac{1}{m(t)}p_cp_f+m(t)\omega ^2(t)q_cq_f$$
(49)
is the interaction Hamiltonian between the classical background and the fluctuation. And then quantize the fluctuation Hamiltonian (48) according to the method in the previous subsection but keep the classical one unquantized. Since the last two terms proportional to $`q_f`$ and $`p_f`$ have the zero expectation value, the expectation value of the total Hamiltonian (47) with respect to the vacuum state of the fluctuation Hamiltonian (48) yields exactly the effective Hamiltonian (40). Therefore, it has been shown that the expectation value of the original Hamiltonian with respect to the coherent state is equivalent to the sum of the classical part (37) and the vacuum expectation value of the fluctuation part (48).
### C Thermal State and Density Matrix
The ensemble of time-dependent oscillators exhibits intrinsically nonequilibrium behaviors, so it does lose a rigorous physical meaning to attribute any thermal property to the density operator (1). However, the harmonic oscillator problem is exactly solvable, so even in the time-dependent oscillator case one may look for a density operator which is quadratic in the position and momentum, and then fix their variable coefficients to satisfy the LvN equation . On the other hand, in the LvN approach we can still use the operators (11) that have already satisfied the LvN equation and define the density operator from them. But there still remains a free parameter to incorporate the initial thermal equilibrium. We study the physical meaning of the density operator and see how the initial thermal equilibrium evolves quantum mechanically .
By noting that $`\widehat{N}(t)`$ satisfies Eq. (4), we define the density operator by
$$\widehat{\rho }_\mathrm{T}(t)=\frac{1}{Z_N}e^{\beta \mathrm{}\omega _0(\widehat{N}(t)+\frac{1}{2})},$$
(50)
where $`\beta `$ and $`\omega _0`$ are free parameters and $`Z_N`$ is the partition function given by
$$Z_N=\underset{n=0}{\overset{\mathrm{}}{}}n,t|e^{\beta \mathrm{}\omega _0(\widehat{N}(t)+\frac{1}{2})}|n,t=\frac{1}{2\mathrm{sinh}(\frac{\beta \mathrm{}\omega _0}{2})}.$$
(51)
It has the same form as the standard density operator, the time-independent annihilation and creation operators now being replaced by the time-dependent ones (11). So Eq. (50) includes the time-independent case as a special case by choosing $`\beta =1/(k_BT)`$ and $`\omega _0`$ the oscillator frequency. In the coordinate representation the density matrix is given by (see Appendix B)
$`\rho _\mathrm{T}(q^{},q,t)`$ $`=`$ $`{\displaystyle \frac{1}{Z_N}}{\displaystyle \underset{n=0}{\overset{\mathrm{}}{}}}\mathrm{\Psi }_n(q^{},t)\mathrm{\Psi }_n^{}(q,t)e^{\beta \mathrm{}\omega _0(n+\frac{1}{2})}`$ (52)
$`=`$ $`\left[{\displaystyle \frac{\mathrm{tanh}(\frac{\beta \mathrm{}\omega _0}{2})}{2\pi \mathrm{}^2u^{}u}}\right]^{1/2}\mathrm{exp}\left[{\displaystyle \frac{i}{4}}{\displaystyle \frac{m}{\mathrm{}}}{\displaystyle \frac{d}{dt}}\mathrm{ln}(u^{}u)(q^2q^2)\right]`$ (54)
$`\times \mathrm{exp}\left[{\displaystyle \frac{1}{8\mathrm{}^2u^{}u}}\left\{(q^{}+q)^2\mathrm{tanh}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}})+(q^{}q)^2\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}})\right\}\right].`$
Now the density matrix (54) can be compared with that for the time-independent oscillator and the density operator (1) for the instantaneous Hamiltonian. For that purpose we restrict our attention to the particular case, in which the mass is constant, $`m(t)=m_0`$, and $`\omega ^2(t)`$ is positive (see Sec. IV for the sign changing case of $`\omega ^2(t)`$) and slowly changing $`|\dot{\omega }(t)/\omega (t)|1`$. In that case we may look for the solution to Eq. (12) of the form
$$u(t)=\frac{1}{\sqrt{2\mathrm{}m_0\mathrm{\Omega }(t)}}e^{i{\scriptscriptstyle \mathrm{\Omega }(t)}},$$
(55)
where
$$\mathrm{\Omega }^2(t)=\omega ^2(t)+\frac{3}{4}\frac{\dot{\mathrm{\Omega }}^2(t)}{\mathrm{\Omega }^2(t)}\frac{1}{2}\frac{\ddot{\mathrm{\Omega }}(t)}{\mathrm{\Omega }(t)}.$$
(56)
The adiabatic (WKB) solution is obtained by approximating $`\mathrm{\Omega }(t)\omega (t)`$. Then the density matrix (54) reduces to the adiabatic one
$`\rho _\mathrm{A}(q^{},q,t)`$ $`=`$ $`\left[{\displaystyle \frac{m_0\omega (t)\mathrm{tanh}(\frac{\beta \mathrm{}\omega _0}{2})}{\pi \mathrm{}}}\right]^{1/2}\mathrm{exp}\left[{\displaystyle \frac{i}{4}}{\displaystyle \frac{m_0}{\mathrm{}}}{\displaystyle \frac{\dot{\omega }(t)}{\omega (t)}}(q^2q^2)\right]`$ (58)
$`\times \mathrm{exp}\left[{\displaystyle \frac{m_0\omega (t)}{4\mathrm{}}}\left\{(q^{}+q)^2\mathrm{tanh}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}})+(q^{}q)^2\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}})\right\}\right].`$
On the other hand, the density operator (1) for the instantaneous Hamiltonian has the matrix representation
$`\rho _\mathrm{H}(q^{},q,t)`$ $`=`$ $`\left[{\displaystyle \frac{m_0\omega (t)\mathrm{tanh}(\frac{\beta \mathrm{}\omega (t)}{2})}{\pi \mathrm{}}}\right]^{1/2}`$ (60)
$`\times \mathrm{exp}\left[{\displaystyle \frac{m_0\omega (t)}{4\mathrm{}}}\left\{(q^{}+q)^2\mathrm{tanh}({\displaystyle \frac{\beta \mathrm{}\omega (t)}{2}})+(q^{}q)^2\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega (t)}{2}})\right\}\right].`$
In the case of the time-independent oscillator with $`\omega (t)=\omega _0`$, the density matrices (54) and (60) reduce further to the standard one . The instantaneous density matrix is compared with the adiabatic one by taking the ratio
$`{\displaystyle \frac{\rho _\mathrm{H}(q^{},q,t)}{\rho _\mathrm{A}(q^{},q,t)}}=\left[{\displaystyle \frac{\mathrm{tanh}(\frac{\beta \mathrm{}\omega (t)}{2})}{\mathrm{tanh}(\frac{\beta \mathrm{}\omega _0}{2})}}\right]^{1/2}\mathrm{exp}\left[{\displaystyle \frac{i}{4}}{\displaystyle \frac{m_0}{\mathrm{}}}{\displaystyle \frac{\dot{\omega }(t)}{\omega (t)}}(q^2q^2)\right]`$ (61)
$`\times \mathrm{exp}\left[{\displaystyle \frac{m_0\omega (t)}{4\mathrm{}}}\left\{(q^{}+q)^2{\displaystyle \frac{\mathrm{sinh}\left(\frac{\beta \mathrm{}}{2}(\omega (t)\omega _0)\right)}{\mathrm{cosh}(\frac{\beta \mathrm{}\omega (t)}{2})\mathrm{cosh}(\frac{\beta \mathrm{}\omega _0}{2})}}(q^{}q)^2{\displaystyle \frac{\mathrm{sinh}\left(\frac{\beta \mathrm{}}{2}(\omega (t)\omega _0)\right)}{\mathrm{sinh}(\frac{\beta \mathrm{}\omega (t)}{2})\mathrm{sinh}(\frac{\beta \mathrm{}\omega _0}{2})}}\right\}\right].`$ (62)
The second factor in Eq. (62) gives rise to a small phase factor because $`\omega (t)`$ is slowly varying. As far as $`\omega (t)`$ remains close to $`\omega _0`$, the instantaneous density matrix (60) is close to the adiabatic one (58). Otherwise, the exact nonequilibrium evolution (58) is far away from the quasi-equilibrium one described by (60). This implies a significant deviation of the nonequilibrium evolution from the equilibrium one as $`\omega (t)`$ differs from $`\omega _0`$ by a large amount.
Being mostly interested in phase transitions, we assume the system to start from an initial thermal equilibrium at early times. This requires that the oscillator have a constant frequency $`\omega _0`$ at early times and $`\beta `$ be fixed to the inverse temperature. As in the case of coherent state, the evolution of the initial thermal state can be found the effective Hamiltonian from this thermal state. From the expectation values (see the appendix of Ref. )
$`\widehat{q}^2_\mathrm{T}`$ $`=`$ $`\mathrm{Tr}\left[\widehat{\rho }_\mathrm{T}(t)\widehat{q}^2\right]=\mathrm{}^2u^{}(t)u(t)\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}}),`$ (63)
$`\widehat{p}^2_\mathrm{T}`$ $`=`$ $`\mathrm{Tr}\left[\widehat{\rho }_\mathrm{T}(t)\widehat{p}^2\right]=\mathrm{}^2m^2(t)\dot{u}^{}(t)\dot{u}(t)\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}}),`$ (64)
one obtains the effective Hamiltonian from the thermal state
$$H_\mathrm{T}(t)=\mathrm{Tr}\left[\widehat{\rho }_\mathrm{T}(t)\widehat{H}(t)\right]=\frac{\mathrm{}^2}{2}m(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _0}{2})\left[\dot{u}^{}(t)\dot{u}(t)+\omega ^2(t)u^{}(t)u(t)\right].$$
(65)
Once again by using the complex parameter $`u`$ in the polar form (39) and by introducing $`p_\xi =m(t)\dot{\zeta }`$, one may rewrite the effective Hamiltonian as
$$H_\mathrm{T}(t)=\mathrm{}\mathrm{coth}(\frac{\beta \mathrm{}\omega _0}{2})\left[\frac{1}{2m(t)}p_\xi ^2+\frac{m(t)}{2}\omega ^2(t)\zeta ^2+\frac{1}{8m(t)\xi ^2}\right]=\mathrm{coth}(\frac{\beta \mathrm{}\omega _0}{2})H_\mathrm{V}(t).$$
(66)
Since $`\mathrm{}\mathrm{coth}(\beta \mathrm{}\omega _0/2)`$ is constant, $`H_\mathrm{T}(t)`$ has the same Hamilton equations (45) as $`H_\mathrm{V}(t)`$. This is identical to the classical equation of motion (12) together with the boundary condition (14).
### D Coherent-Thermal State
A more general density operator that is at most quadratic in the position and momentum was introduced in Ref. . It has the form
$$\widehat{\rho }_{\mathrm{C}.\mathrm{T}}(t)=\frac{1}{Z_{\mathrm{C}.\mathrm{T}}(t)}\mathrm{exp}\left[\beta \left\{\mathrm{}\omega _0\widehat{a}^{}(t)\widehat{a}(t)+\delta \widehat{a}^{}(t)+\delta ^{}\widehat{a}(t)+ϵ_0\right\}\right],$$
(67)
where $`Z_{\mathrm{C}.\mathrm{T}}`$ is a partition function. In fact, the density operator of Eq. (67) can be transformed into that of Eq. (50) by the unitary transformation
$$\widehat{D}^{}(\alpha )\widehat{\rho }_\mathrm{C}(t)\widehat{D}(\alpha )=\widehat{\rho }(t),$$
(68)
where $`\widehat{D}(\alpha )`$ with $`\alpha =\delta /(\mathrm{}\omega _0)`$ and $`ϵ_0=(\mathrm{}\omega _0/2)+(|\delta |^2/\mathrm{}\omega _0)`$. From the expectation values
$`\widehat{q}^2_{\mathrm{C}.\mathrm{T}}`$ $`=`$ $`\mathrm{Tr}\left[\widehat{\rho }_{\mathrm{C}.\mathrm{T}}(t)\widehat{q}^2\right]=q_c^2+\mathrm{}^2u^{}(t)u(t)\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}}),`$ (69)
$`\widehat{p}^2_{\mathrm{C}.\mathrm{T}}`$ $`=`$ $`\mathrm{Tr}\left[\widehat{\rho }_{\mathrm{C}.\mathrm{T}}(t)\widehat{p}^2\right]=p_c^2+\mathrm{}^2m^2(t)\dot{u}^{}(t)\dot{u}(t)\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}}),`$ (70)
follows the effective Hamiltonian
$$H_{\mathrm{C}.\mathrm{T}}(t)=H_c(t)+H_\mathrm{T}(t),$$
(71)
where $`H_c`$ is the classical Hamiltonian. The effective Hamiltonian (71) has almost the same form as Eq. (40) except for the overall factor $`\mathrm{coth}(\beta \mathrm{}\omega _0/2)`$, hence describes the same Hamilton equations (44).
In summary, what we have shown in this section is that the LvN approach provides us with the exact Fock space and various quantum states for time-dependent oscillators. The only auxiliary field necessary for that purpose is a complex solution $`u`$ to the classical equation of motion (12) that satisfies the wronskian condition (14) from quantization. It has further been shown that the LvN approach is equivalent to the minimization principle for the effective action and to the mean field method. However, the advantage of the LvN approach lies in the manifest correspondence with quantum mechanics and quantum many-particle system and in the readiness to apply their standard techniques. For instance, the density matrix (54) has been found, which has many features similar with the standard one.
## IV Inverted Harmonic Oscillator
As a toy model for the second order phase transition, let us consider the time-dependent harmonic oscillator
$$\widehat{H}=\frac{1}{2}\widehat{p}^2+\frac{1}{2}\omega ^2(t)\widehat{q}^2,$$
(72)
where $`\omega ^2(t)`$ has the asymptotic value $`\omega _i^2(>0)`$ far before and $`\omega _f^2(<0)`$ far after the quench. The conspicuous point of the model is the sign change of $`\omega ^2(t)`$. At earlier times before the quench the oscillator executes a stable motion about $`q=0`$, the global minimum, but after the quench the potential is inverted and $`q=0`$ becomes an unstable configuration.
At earlier times far before the phase transition, the complex solution to Eq. (12) satisfying Eq. (14) is given by
$$u_i(t)=\frac{e^{i\omega _it}}{\sqrt{2\mathrm{}\omega _i}}.$$
(73)
According to Eq. (11), the Fock space is now constructed from the annihilation and creation operators
$`\widehat{a}(t)`$ $`=`$ $`{\displaystyle \frac{e^{i\omega _it}}{\sqrt{2\mathrm{}\omega _i}}}\left(i\widehat{p}+\omega _i\widehat{q}\right)=e^{i\omega _it}\widehat{a}_0,`$ (74)
$`\widehat{a}^{}(t)`$ $`=`$ $`{\displaystyle \frac{e^{i\omega _it}}{\sqrt{2\mathrm{}\omega _i}}}\left(i\widehat{p}+\omega _i\widehat{q}\right)=e^{i\omega _it}\widehat{a}_0^{}.`$ (75)
Note that $`\widehat{a}(t)`$ and $`\widehat{a}^{}(t)`$ differ from the standard ones $`\widehat{a}_0`$ and $`\widehat{a}_0^{}`$ only by phase factors. Though the Hamiltonian has the standard representation
$$\widehat{H}_i=\mathrm{}\omega _i\left(\widehat{a}^{}(t)\widehat{a}(t)+\frac{1}{2}\right)=\mathrm{}\omega _i\left(\widehat{a}_0^{}\widehat{a}_0+\frac{1}{2}\right),$$
(76)
the phase factors are necessary for $`\widehat{a}(t)`$ and $`\widehat{a}^{}(t)`$ to satisfy the LvN equation. Hence the vacuum expectation value is given by the well-known result
$$H_\mathrm{V}=\frac{1}{2}\mathrm{}\omega _i,$$
(77)
and the coherent state (33) yields
$$H_\mathrm{C}=\frac{1}{2}p_c^2+\frac{1}{2}\omega _i^2q_c^2+\frac{1}{2}\mathrm{}\omega _i.$$
(78)
Now the density operator (50) reduces to the standard one
$$\widehat{\rho }_{i,\mathrm{T}}=\frac{1}{Z_N}e^{\beta \mathrm{}\omega _i(\widehat{a}_0^{}\widehat{a}_0+\frac{1}{2})},$$
(79)
after identifying $`\omega _0=\omega _i`$ and $`\beta =1/(k_BT)`$, and leads to the Hamiltonian expectation value
$$H_{i,\mathrm{T}}=\frac{1}{2}\mathrm{}\omega _i\mathrm{coth}(\frac{\beta \mathrm{}\omega _i}{2}).$$
(80)
On the other hand, at later times far after the quench, the solution to Eq. (12) is given by
$`u_f(t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\mathrm{}}}}\left[C_1(\omega _i,\omega _f)\mathrm{cosh}(\omega _ft)iC_2(\omega _i,\omega _f)\mathrm{sinh}(\omega _ft)\right]`$ (81)
$`=`$ $`{\displaystyle \frac{1}{2\sqrt{2\mathrm{}}}}\left[\left(C_1(\omega _i,\omega _f)iC_2(\omega _i,\omega _f)\right)e^{\omega _ft}+\left(C_1(\omega _i,\omega _f)+iC_2(\omega _i,\omega _f)\right)e^{\omega _ft}\right],`$ (82)
where $`C_j,j=1,2`$ depend on the intermediate process toward the final state. Remarkably, the vacuum and thermal expectation values vanish:
$$H_{f,\mathrm{V}}=H_{f,\mathrm{T}}=0,$$
(83)
since the kinetic and potential energies contribute equally. Whereas the expectation value with respect to the coherent state (33) and the general density operator (67) takes the form
$$H_{f,\mathrm{C}}=H_{f,\mathrm{C}.\mathrm{T}}=\frac{1}{2}p_c^2\frac{1}{2}\omega _f^2q_c^2.$$
(84)
This implies physically that as the system undergoes the phase transition out of equilibrium, quantum effects vanish and it becomes classical. Large uncertainty has been suggested as one of the criteria on the classicality
$$(\mathrm{\Delta }q)(\mathrm{\Delta }p)=\frac{\mathrm{}}{2}\left[\omega _fC_1^2(\omega _i,\omega _f)\mathrm{cosh}^2(\omega _ft)+\omega _fC_1^2(\omega _i,\omega _f)\mathrm{sinh}^2(\omega _ft)\right].$$
(85)
Though the oscillator starts with the minimum uncertainty given by Eq. (73), its uncertainty increases exponentially and it becomes eventually classical after the completion of quench. When Eq. (82) is substituted into Eq. (54), the density matrix depends on the intermediate process and spreads as $`\sqrt{u_f^{}(t)u_f(t)}`$. Therefore the quintessence of second order phase transitions lies in the whole process how systems evolve out of equilibrium from their initial equilibrium.
To show explicitly how the nonequilibrium dynamics depends on the intermediate processes, we consider two exactly solvable models. The first model, which is the zero mode of the free scalar model in Sec. VI, is an oscillator describing a finite smooth quench with the mass
$$m^2(t)=m_1^2m_0^2\mathrm{tanh}\left(\frac{t}{\tau }\right).$$
(86)
The mass has $`m_i^2=m_0^2+m_1^2`$ at earlier times $`(t=\mathrm{})`$ and has $`m_f^2=m_0^2+m_1^2<0`$ at later times $`(t=\mathrm{})`$. $`\tau `$ measures the quench rate, i.e. the rate of change of mass. The instantaneous quench corresponds to the $`(\tau =0)`$-limit. The solution to Eq. (12) is found
$$u(t)=\frac{e^{m_it}}{\sqrt{2\mathrm{}m_i}}{}_{2}{}^{}F_{1}^{}(\frac{\tau }{2}(im_im_f),\frac{\tau }{2}(im_i+m_f);1i\tau m_i;e^{2t/\tau }).$$
(87)
At earlier times the solution (87) has the correct asymptotic form
$$u_i(t)=\frac{e^{im_it}}{\sqrt{2\mathrm{}m_i}}.$$
(88)
On the other hand, at later times the asymptotic form of Eq. (87) becomes
$`u_f(t)`$ $`=`$ $`\left[{\displaystyle \frac{1}{\sqrt{2\mathrm{}m_i}}}{\displaystyle \frac{(1)\mathrm{\Gamma }(1im_i\tau )\mathrm{\Gamma }(m_f\tau )}{\frac{\tau }{2}(im_im_f)\mathrm{\Gamma }^2(\frac{\tau }{2}(im_im_f)}}\right]e^{m_ft}`$ (90)
$`+\left[{\displaystyle \frac{1}{\sqrt{2\mathrm{}m_i}}}{\displaystyle \frac{(1)\mathrm{\Gamma }(1im_i\tau )\mathrm{\Gamma }(m_f\tau )}{\frac{\tau }{2}(im_i+m_f)\mathrm{\Gamma }^2(\frac{\tau }{2}(im_i+m_f)}}\right]e^{m_ft}.`$
Two points are observed: the initial solution branches into an unstable growing and a decaying mode as expected and the coefficients $`C_1,C_2`$ of Eq. (82) depend on the the mass parameters $`m_i,m_f`$ and the quench rate $`\tau `$. In other words, the final asymptotic state of nonequilibrium evolution depends on the intermediate process.
The next model describes various quench processes and exhibits how the stable mode far before the quench branches into the unstable growing and decaying modes during the quench. Without loss of generality the nonequilibrium phase transition is assumed to take place through the time-dependent frequency (mass) squared
$`\omega ^2(t)=\{\begin{array}{cc}\omega _i^2,\hfill & t_i>t\text{,}\hfill \\ \omega _i^2\left(\frac{t_0t}{t_0t_i}\right)^{(2l_1+1)/(2l_2+1)},\hfill & t_0>t>t_i\text{,}\hfill \\ \omega _i^2\left(\frac{tt_0}{t_0t_i}\right)^{(2l_1+1)/(2l_2+1)},\hfill & t_f>t>t_0\text{,}\hfill \\ \omega _f^2\omega _i^2\left(\frac{t_ft_0}{t_0t_i}\right)^{(2l_1+1)/(2l_2+1)},\hfill & t>t_f\text{,}\hfill \end{array}`$ (91)
where $`l_1`$ and $`l_2`$ are non-negative integers. Here $`t_0t_i`$ adjusts the rate of and $`t_ft_0`$ the duration of the quench. The particular case of $`l_1=l_2=0`$ is used as the finite linear quench model . Before the time $`t_0`$, the system maintains the symmetry about $`q=0`$, the minimum of the potential. But as time goes on after $`t_0`$, $`q=0`$ remains no longer the true minimum of the system and the symmetry is broken. The particular form of the power-law in Eq. (91) is chosen to allow an analytical continuation of $`\omega ^2(t)`$ for changing the sign and to make its derivatives also continuous. Before $`t_i>t`$, the solution is given by Eq. (73). During $`t_0>t>t_i`$, the solution is given by a linear superposition of Hankel functions
$$u(t)=D_1z^\nu H_\nu ^{(2)}(z)+D_2z^\nu H_\nu ^{(1)}(z),$$
(92)
where
$$z=2\nu \omega _i(t_0t_i)\left(\frac{t_0t}{t_0t_i}\right)^{1/2\nu },\nu =\frac{2l_2+1}{2l_1+4l_2+3}.$$
(93)
Here $`H_\nu ^{(2)}`$ and $`H_\nu ^{(1)}`$ are positive and negative frequency solutions, respectively. The constants $`D_1`$ and $`D_2`$ are determined by continuity of $`u(t)`$ and $`\dot{u}(t)`$ across $`t_i`$:
$`D_1`$ $`=`$ $`{\displaystyle \frac{e^{i\omega _it_i}}{\sqrt{2\mathrm{}\omega _i}}}{\displaystyle \frac{\pi }{4z_i^\nu }}\left[iz_i{\displaystyle \frac{d}{dz_i}}H_\nu ^{(1)}(z_i)\left(\omega _ii{\displaystyle \frac{1}{2(t_0t_i)}}\right)H_\nu ^{(2)}(z_i)\right],`$ (94)
$`D_2`$ $`=`$ $`{\displaystyle \frac{e^{i\omega _it_i}}{\sqrt{2\mathrm{}\omega _i}}}{\displaystyle \frac{\pi }{4z_i^\nu }}\left[iz_i{\displaystyle \frac{d}{dz_i}}H_\nu ^{(2)}(z_i)+\left(\omega _ii{\displaystyle \frac{1}{2(t_0t_i)}}\right)H_\nu ^{(1)}(z_i)\right],`$ (95)
where
$$z_i=2\nu \omega _i(t_0t_i).$$
(96)
Beyond the quench time $`t_0`$, it is necessary to do carefully the analytic continuation and to take a suitable Riemann sheet so that $`\omega ^2(t)`$ and its derivatives are to be continuous from $`\omega ^2>0`$ to $`\omega ^2<0`$ across $`t_0`$. The analytic continuation of the solution (92) yields
$`u(t)`$ $`=`$ $`{\displaystyle \frac{D_1}{2}}\stackrel{~}{z}^\nu \left[e^{i3\pi (l_2+\frac{1}{2})}H_\nu ^{(2)}(i\stackrel{~}{z})+e^{i\pi (l_2+\frac{1}{2})}H_\nu ^{(2)}(i\stackrel{~}{z})\right]`$ (98)
$`+{\displaystyle \frac{D_2}{2}}\stackrel{~}{z}^\nu \left[e^{i3\pi (l_2+\frac{1}{2})}H_\nu ^{(1)}(i\stackrel{~}{z})+e^{i\pi (l_2+\frac{1}{2})}H_\nu ^{(1)}(i\stackrel{~}{z})\right],`$
where
$$\stackrel{~}{z}=2\nu \omega _i(t_0t_i)\left(\frac{tt_0}{t_0t_i}\right)^{1/2\nu }.$$
(99)
At later times $`(\stackrel{~}{z}1)`$ during the quench, the solution (98) has the asymptotic form
$$u_f(t)=\sqrt{\frac{1}{2\pi }}\stackrel{~}{z}^{\nu \frac{1}{2}}e^{\stackrel{~}{z}}\left[D_1e^{i\pi \left(3l_2+\frac{\nu }{2}+\frac{3}{2}\right)}+D_2e^{i\pi \left(l_2\frac{\nu }{2}\frac{1}{2}\right)}\right].$$
(100)
Thus the stable mode of the oscillating solution (73) branches into the growing mode (100) which dominates during various quench processes and into the decaying mode which contributes negligibly to the correlation functions. The asymptotic solution (100) also depends on the intermediate processes through $`t_0,t_i,t_f`$ and $`l_1,l_2`$.
## V Time-Dependent Anharmonic Oscillator
In this section we extend the formalism developed in Sec. III to the time-dependent anharmonic oscillator with the Hamiltonian
$$H(t)=\frac{p^2}{2m(t)}+m(t)V(q),$$
(101)
where
$$V(q)=\frac{\lambda _{2n}(t)}{(2n)!}q^{2n}.$$
(102)
Though the potential of a power law is assumed, the formalism can readily be generalized to any polynomial and analytic potential. In the time-independent case $`(\lambda _{2n}=\mathrm{constant})`$, the variational perturbation method has been introduced as one of the powerful methods to find the Hilbert space . The vacuum state in this approach is the Gaussian wave functional that minimizes the effective action. The excited states are then obtained from the vacuum state just as number states of a harmonic oscillator are obtained from the ground state. Though these states can be calculated explicitly in terms a complex solution to the classical equation for motion, they are in fact equivalent to those from the mean-field method.
However, there have been some attempts to go beyond the Gaussian state. In Ref. a scheme was proposed to find the operators, which generalize the annihilation and creation operators and satisfy the LvN equation (4), to all the orders of coupling constant in the time-independent case and to the first order in the time-dependent case. In particular, the generalized annihilation and creation operators for the time-independent oscillator with a quartic potential satisfy a q-deformed algebra rather than the standard commutation relation , from which follows an algebraic construction of excited states and energy spectra beyond the variational Gaussian approximation. It would be interesting to find such an algebraic structure for interacting quantum fields, which may shed some light on the nonperturbative method beyond the mean-field method. Also it would be interesting to compare this scheme with other nonperturbative methods in the time-independent case such as the perturbative expansion method around the Gaussian effective action and the time-dependent variational method . But we shall not pursue further this issue in this paper.
### A Fock Space
In the case of time-dependent anharmonic oscillators, the LvN approach searches for the annihilation and creation operators that are still linear in the position and momentum and satisfy the LvN equation (4):
$`\widehat{A}(t)`$ $`=`$ $`i\left(v^{}(t)\widehat{p}m(t)\dot{v}^{}(t)\widehat{q}\right),`$ (103)
$`\widehat{A}^{}(t)`$ $`=`$ $`i\left(v(t)\widehat{p}m(t)\dot{v}(t)\widehat{q}\right).`$ (104)
One then requires them to satisfy the LvN equation (4), leading to the equation
$$\ddot{v}(t)\widehat{q}+\frac{\dot{m}(t)}{m(t)}\dot{v}(t)\widehat{q}+v(t)\frac{\delta V(\widehat{q})}{\delta \widehat{q}}=0,$$
(105)
and further differentiates functionally with respect to $`\widehat{q}`$ and takes the vacuum expectation value of the resultant equation
$$\ddot{v}(t)+\frac{\dot{m}(t)}{m(t)}\dot{v}(t)+0,t|\frac{\delta ^2V(\widehat{q})}{\delta \widehat{q}^2}|0,tv(t)=0.$$
(106)
Here the vacuum state is annihilated by $`\widehat{A}(t)`$
$$\widehat{A}(t)|0,t=0.$$
(107)
One makes $`\widehat{A}(t)`$ and $`\widehat{A}^{}(t)`$ the annihilation and creation operators, respectively, by imposing the standard commutation relation for all times
$$[\widehat{A}(t),\widehat{A}^{}(t)]=1.$$
(108)
Equation (108) is equivalent to the wronskian condition
$$\mathrm{}m(t)\left(\dot{v}^{}(t)v(t)\dot{v}(t)v^{}(t)\right)=i.$$
(109)
The Fock space consists of the number state obtained by applying $`\widehat{A}^{}(t)`$ $`n`$-times to the vacuum state
$$|n,t=\frac{\left(\widehat{A}^{}(t)\right)^n}{\sqrt{n!}}|0,t.$$
(110)
These number states are excited states and have the coordinate representation (19) now with $`u(t)`$ replaced by $`v(t)`$.
¿From the position and momentum operators expressed in terms of the annihilation and creation operators
$`\widehat{q}=\mathrm{}\left(v(t)\widehat{A}(t)+v^{}(t)\widehat{A}^{}(t)\right),`$ (111)
$`\widehat{p}=\mathrm{}m(t)\left(\dot{v}(t)\widehat{A}(t)+\dot{v}^{}(t)\widehat{A}^{}(t)\right),`$ (112)
follow the vacuum expectation values
$`0,t|\widehat{q}^{2n}|0,t={\displaystyle \frac{(2n)!}{2^nn!}}\left[\mathrm{}^2v^{}(t)v(t)\right]^n,`$ (113)
$`0,t|\widehat{p}^2|0,t=\mathrm{}^2m^2(t)\dot{v}^{}(t)\dot{v}(t).`$ (114)
Then the classical equation of motion (106) becomes
$$\ddot{v}(t)+\frac{\dot{m}(t)}{m(t)}\dot{v}(t)+\frac{\lambda _{2n}(t)}{2^{n1}(n1)!}\left[\mathrm{}^2v^{}(t)v(t)\right]^{n1}v(t)=0.$$
(115)
There is another method to derive Eq. (115). By using the Wick-ordering
$$\widehat{q}^{2n}=\underset{k=0}{\overset{n}{}}\frac{(2n)!\mathrm{}^{2n}}{2^kk!(2n2k)!}[v^{}(t)v(t)]^k:[v(t)\widehat{A}(t)+v^{}(t)\widehat{A}^{}(t)]^{2(nk)}:,$$
(116)
one obtains the Hamiltonian truncated at the quadratic order of $`\widehat{A}(t)`$ and $`\widehat{A}(t)`$
$`\widehat{H}_\mathrm{G}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\mathrm{}^2m(t)\left[(\dot{v}(t)\widehat{A}(t))^2+2\dot{v}^{}(t)\dot{v}(t)\widehat{A}^{}(t)\widehat{A}(t)+(\dot{v}^{}(t)\widehat{A}^{}(t))^2\right]`$ (117)
$`+`$ $`m(t){\displaystyle \frac{\mathrm{}^{2n}\lambda _{2n}(t)}{2^n(n1)!}}[v^{}(t)v(t)]^{n1}\left[(v(t)\widehat{A}(t))^2+2v^{}(t)v(t)\widehat{A}^{}(t)\widehat{A}(t)+(v^{}(t)\widehat{A}^{}(t))^2\right],`$ (118)
where purely $`c`$-number terms are neglected. One then requires $`\widehat{A}^{}(t)`$ and $`\widehat{A}(t)`$ to satisfy the LvN equation (4) for the truncated Hamiltonian $`\widehat{H}_\mathrm{G}(t)`$. Now the LvN equations for $`\widehat{A}^{}(t)`$ and $`\widehat{A}(t)`$ lead exactly to the equation of motion (115).
### B Effective Hamiltonians
Still another method to derive the equations of motion for $`v`$ and $`q_c`$ is the minimization principle for the effective action. For that purpose we consider a complex $`v(t)`$-parameter family of the Fock spaces constructed by the annihilation and creation operators (104) and take the Hamiltonian expectation value with respect to various quantum states such as the vacuum, coherent, thermal and coherent-thermal states. We do not require $`\widehat{A}(t)`$ and $`\widehat{A}^{}(t)`$ to satisfy the LvN equation a priori, but minimize the action to determine the equation of motion for $`v(t)`$.
The first effective Hamiltonian is the vacuum expectation value
$$H_\mathrm{V}(t)=0,t|\widehat{H}(t)|0,t=\frac{\mathrm{}^2}{2}m(t)\dot{v}^{}\dot{v}+m(t)\frac{\lambda _{2n}(t)}{2^nn!}\left[\mathrm{}^2v^{}(t)v(t)\right]^n.$$
(119)
The next state under consideration is the coherent state, which is obtained by applying a displacement operator to the vacuum state,
$$|\alpha ,t=\widehat{D}^{}(\alpha )|0,t=e^{\alpha \widehat{A}^{}(t)\alpha ^{}\widehat{A}(t)}|0,t.$$
(120)
Then the coherent state expectation value leads to the effective Hamiltonian
$$H_\mathrm{C}(t)=\alpha ,t|\widehat{H}(t)|\alpha ,t,$$
(121)
which, with the aid of Eq. (116), is decomposed into
$$H_\mathrm{C}(t)=H_c(t)+H_q(t).$$
(122)
Here $`H_c(t)`$ is the classical part
$$H_c(t)=\frac{p_c^2}{2m(t)}+m(t)\frac{\lambda _{2n}(t)}{(2n)!}q_c^{2n},$$
(123)
and $`H_q(t)`$ denotes all the quantum contributions including the vacuum expectation value
$$H_q(t)=\frac{\mathrm{}^2}{2}m(t)\dot{v}^{}\dot{v}+m(t)\underset{k=1}{\overset{n}{}}\frac{\lambda _{2n}(t)}{2^kk!(2n2k)!}\left[\mathrm{}^2v^{}(t)v(t)\right]^kq_c^{2(nk)},$$
(124)
where
$$q_c=\alpha ,t|\widehat{q}|\alpha ,t,p_c=\alpha ,t|\widehat{p}|\alpha ,t.$$
(125)
The final state is the thermal state defined by the density operator
$$\widehat{\rho }_\mathrm{T}=\frac{1}{Z_N}e^{\beta \omega _0(\widehat{A}^{}(t)\widehat{A}(t)+\frac{1}{2})},$$
(126)
with $`Z_N`$ being the partition function. The density operator (126) leads to the effective Hamiltonian
$$H_\mathrm{T}(t)=\mathrm{Tr}\left[\widehat{\rho }_\mathrm{T}\widehat{H}(t)\right]=\frac{\mathrm{}^2}{2}m(t)\dot{v}^{}\dot{v}+m(t)\frac{\lambda _{2n}(t)}{2^nn!}\widehat{q}^2_\mathrm{T}^n,$$
(127)
where
$$\widehat{q}^2_\mathrm{T}=\mathrm{}^2v^{}v\mathrm{coth}(\frac{\beta \omega _0}{2}).$$
(128)
Likewise, the density operator of the form (67) with $`\widehat{a}^{}(t)`$ and $`\widehat{a}(t)`$ replaced by $`\widehat{A}^{}(t)`$ and $`\widehat{A}(t)`$ leads to the effective Hamiltonian
$$H_{\mathrm{C}.\mathrm{T}}(t)=H_c(t)+\frac{\mathrm{}^2}{2}m(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _0}{2})\dot{v}^{}\dot{v}+m(t)\underset{k=1}{\overset{n}{}}\frac{\lambda _{2n}(t)}{2^kk!(2n2k)!}\widehat{q}^2_\mathrm{T}^kq_c^{2(nk)}.$$
(129)
Note that Eqs. (119) and (122) are also obtained from Eqs. (127) and (129), respectively, by replacing $`\mathrm{}^2v^{}v\mathrm{coth}(\beta \omega _0/2)`$ with $`\mathrm{}^2v^{}v`$ or taking the zero-temperature limit $`(\beta \mathrm{})`$.
We now study the dynamics of the effective Hamiltonians. We mainly focus on the effective Hamiltonian (129) since Eq. (122) is the limiting case of Eq. (129) when $`\beta \mathrm{}`$, i.e., $`T0`$, and Eq. (119) is the limiting case of Eq. (122) when $`q_c=p_c=0`$. The independent variables of the Hamiltonian (129) are $`(q_c,p_c)`$, $`(v,p_v=m\dot{v}^{})`$ and $`(v^{},p_v^{}=m\dot{v})`$. So we obtain the equation of motion for $`q_c`$
$$\ddot{q}_c+\frac{\dot{m}}{m}\dot{q}_c+\frac{\lambda _{2n}(t)}{(2n1)!}q_c^{2n1}+\underset{k=1}{\overset{n}{}}\frac{\lambda _{2n}(t)}{2^kk!(2n2k1)!}\widehat{q}^2_\mathrm{T}^kq_c^{2n2k1}=0.$$
(130)
The equation of motion for $`v^{}`$ is given by
$$\ddot{v}+\frac{\dot{m}}{m}\dot{v}+\underset{k=1}{\overset{n}{}}\frac{\lambda _{2n}(t)}{2^{k1}(k1)!(2n2k)!}q_c^{2n2k}\widehat{q}^2_{\mathrm{T}.}^{k1}v=0,$$
(131)
and the complex conjugate of Eq. (131) is for $`v`$. The equations of motion from the effective Hamiltonian (122) is the limiting case of Eqs. (130) and (131) when $`\widehat{q}^2_{(\mathrm{T})}=\mathrm{}^2v^{}v`$:
$`\ddot{q}_c+{\displaystyle \frac{\dot{m}}{m}}\dot{q}_c+{\displaystyle \frac{\lambda _{2n}(t)}{(2n1)!}}q_c^{2n1}+{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{\lambda _{2n}(t)}{2^kk!(2n2k1)!}}(\mathrm{}^2v^{}v)^kq_c^{2n2k1}=0,`$ (132)
$`\ddot{v}+{\displaystyle \frac{\dot{m}}{m}}\dot{v}+{\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{\lambda _{2n}(t)}{2^{k1}(k1)!(2n2k)!}}q_c^{2n2k}(\mathrm{}^2v^{}v)^{k1}v=0.`$ (133)
And the limiting case $`q_c=0`$ of Eq. (133) is the equation for the effective Hamiltonian (119)
$$\ddot{v}+\frac{\dot{m}}{m}\dot{v}+\frac{\lambda _{2n}(t)}{2^{n1}(n1)!}(\mathrm{}^2v^{}v)^{n1}v=0.$$
(134)
Note that Eq. (134) is identical to Eq. (115) from the LvN approach.
Or, by writing $`v`$ in the polar form
$$v(t)=\frac{\zeta (t)}{\sqrt{\mathrm{}}}e^{i\theta (t)},$$
(135)
and by introducing the momentum $`p_\zeta =m(t)\dot{\zeta }`$, the effective Hamiltonian (129) is rewritten as
$`H_{\mathrm{C}.\mathrm{T}}(t)={\displaystyle \frac{p_c^2}{2m(t)}}+m(t){\displaystyle \frac{\lambda _{2n}(t)}{(2n)!}}q_c^{2n}+\mathrm{}\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}})\left[{\displaystyle \frac{p_\zeta ^2}{2m(t)}}+{\displaystyle \frac{1}{8m(t)\zeta ^2}}\right]`$ (136)
$`+m(t){\displaystyle \underset{k=1}{\overset{n}{}}}{\displaystyle \frac{\lambda _{2n}(t)}{2^kk!(2n2k)!}}\left[\mathrm{}\zeta ^2\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}})\right]^kq_c^{2(nk)}.`$ (137)
Now the equation of motion for the classical field $`q_c`$ is given by
$$\ddot{q}_c+\frac{\dot{m}}{m}\dot{q}_c+\frac{\lambda _{2n}(t)}{(2n1)!}q_c^{2n1}+\underset{k=1}{\overset{n}{}}\frac{\lambda _{2n}(t)}{2^kk!(2n2k1)!}[\mathrm{}\zeta ^2\mathrm{coth}(\frac{\beta \mathrm{}\omega _0}{2})]^kq_c^{2n2k1}=0,$$
(138)
and that for $`\zeta `$ by
$$\ddot{\zeta }+\frac{\dot{m}}{m}\dot{\zeta }\frac{1}{4m^2\zeta ^3}+\underset{k=1}{\overset{n}{}}\frac{\lambda _{2n}(t)}{2^{k1}(k1)!(2n2k)!}q_c^{2n2k}[\mathrm{}\zeta ^2\mathrm{coth}(\frac{\beta \omega _0}{2})]^{k1}\zeta =0.$$
(139)
The phase of $`v(t)`$ is obtained by the integration
$$\theta (t)=\frac{1}{2m(t)\zeta ^2(t)}.$$
(140)
### C Coherent State vs. Hartree-Fock Method
In this subsection we show that the nonequilibrium dynamics obtained from the effective Hamiltonian (129) in the LvN approach recovers exactly the equations of motion from the mean field and Hartree-Fock methods.
First, the effective Hamiltonian from the coherent state can be obtained in another way. By dividing $`q`$ and $`p`$ into a classical background and a quantum fluctuation
$$q=q_c+q_f,$$
(141)
we obtain the expectation value for the Hamiltonian (101) with respect to the thermal state (126)
$$H_{\mathrm{c}.\mathrm{f}.}(t)=\frac{p_c^2}{2m(t)}+\frac{\widehat{p}_f^2_\mathrm{T}}{2m(t)}+m(t)\underset{k=0}{\overset{n}{}}\frac{\lambda _{2n}(t)}{2^kk!(2n2k)!}\widehat{q}_f^2_\mathrm{T}^kq_c^{2(nk)}.$$
(142)
where we have used
$$\widehat{q}^{2n}_\mathrm{T}=\underset{k=0}{\overset{n}{}}\frac{(2n)!}{2^kk!(2n2k)!}q_c^{2(nk)}\left[\mathrm{}^2v^{}v\mathrm{coth}(\frac{\beta \mathrm{}\omega _0}{2})\right]^k,$$
(143)
and
$$\widehat{q}_f^{2n+1}_\mathrm{T}=0=\widehat{p}_f^{2n+1}_\mathrm{T}.$$
(144)
Noting that $`k=0`$ term recovers the classical potential and
$$\widehat{p}_f^2_\mathrm{T}=m^2\mathrm{}^2\dot{v}^{}\dot{v}\mathrm{coth}(\frac{\beta \mathrm{}\omega _0}{2}),$$
(145)
we can show that Eq. (142) coincides with Eq. (129). Therefore the coherent state leads exactly to the result from the mean-field method.
Second, the Hartree-Fock factorization theorem leads to the effective Hamiltonian
$$H_{\mathrm{H}.\mathrm{F}}(t)=\frac{\widehat{p}^2_{\mathrm{H}.\mathrm{F}}}{2m(t)}+m(t)\frac{\lambda _{2n}(t)}{(2n)!}\widehat{q}^{2n}_{\mathrm{H}.\mathrm{F}},$$
(146)
where
$`\widehat{p}^2_{\mathrm{H}.\mathrm{F}}=p_c^2+\widehat{p}_f^2,`$ (147)
$`\widehat{q}^2_{\mathrm{H}.\mathrm{F}}=q_c^2+\widehat{q}_f^2`$ (148)
and
$`\widehat{q}^{2n}_{\mathrm{H}.\mathrm{F}}`$ $`=`$ $`{\displaystyle \underset{k=0}{\overset{2n}{}}}{\displaystyle \frac{(2n)!}{k!(2nk)!}}q_c^{2nk}\widehat{q}_f^k`$ (149)
$`=`$ $`{\displaystyle \underset{k=0}{\overset{n}{}}}{\displaystyle \frac{(2n)!}{2^kk!(2n2k)!}}q_c^{2(nk)}\left[k\widehat{q}_f^2_\mathrm{T}^{k1}\widehat{q}_f^2(k1)\widehat{q}_f^2_\mathrm{T}^k\right]`$ (151)
$`+{\displaystyle \underset{k=0}{\overset{n1}{}}}{\displaystyle \frac{(2n)!}{2^kk!(2n2k1)!}}q_c^{2n2k1}\widehat{q}_f^2_\mathrm{T}^k\widehat{q}_f.`$
for $`n2`$. The thermal expectation value of the equation of motion for $`q_c`$ yields
$$\ddot{q}_c+\frac{\dot{m}}{m}\dot{q}_c+\underset{k=0}{\overset{n}{}}\frac{\lambda _{2n}(t)}{2^kk!(2n2k)!}\widehat{q}_f^2_\mathrm{T}^kq_c^{2n2k1}=0,$$
(152)
and for $`\widehat{q}_f`$
$$\ddot{\widehat{q}}_f+\frac{\dot{m}}{m}\dot{\widehat{q}}_f+\underset{k=0}{\overset{n}{}}\frac{\lambda _{2n}(t)}{2^{k1}(k1)!(2n2k)!}\widehat{q}_f^2_\mathrm{T}^{k1}q_c^{2(nk)}\widehat{q}_f=0.$$
(153)
Therefore it has been shown that Eqs. (152) and (153) are the same as Eqs. (130) and (131) from the coherent state representation.
## VI Free Scalar Field for Phase Transition
As a simple field model for the second order phase transition, we consider a free complex scalar field, the mass of which changes the sign during the quench.<sup>§</sup><sup>§</sup>§The complex scalar field model may be related with the time-dependent Landau-Ginsburg theory provided that the free energy be interpreted as the Hamiltonian in this paper . The system is described by the Lagrangian density
$$(𝐱,t)=\dot{\mathrm{\Phi }}^{}(𝐱,t)\dot{\mathrm{\Phi }}(𝐱,t)\mathrm{\Phi }^{}(𝐱,t)\mathrm{\Phi }(𝐱,t)m^2(t)\mathrm{\Phi }^{}(𝐱,t)\mathrm{\Phi }(𝐱,t).$$
(154)
Here the coupling parameter $`m^2(t)`$ is assumed to begin with an initial positive value before, to change the sign during, and to reach a final negative value after the quench. The $`\mathrm{\Phi }`$ and $`\mathrm{\Phi }^{}`$ are treated as independent fields. The Hamiltonian is given by
$$H(t)=d^3𝐱\left[\mathrm{\Pi }^{}(𝐱,t)\mathrm{\Pi }(𝐱,t)+\mathrm{\Phi }^{}(𝐱,t)\mathrm{\Phi }(𝐱,t)+m^2(t)\mathrm{\Phi }^{}(𝐱,t)\mathrm{\Phi }(𝐱,t)\right],$$
(155)
where
$`\mathrm{\Pi }(𝐱,t)={\displaystyle \frac{\delta (𝐱,t)}{\delta \dot{\mathrm{\Phi }}(𝐱,t)}}=\dot{\mathrm{\Phi }}^{}(𝐱,t),`$ (156)
$`\mathrm{\Pi }^{}(𝐱,t)={\displaystyle \frac{\delta (𝐱,t)}{\delta \dot{\mathrm{\Phi }}^{}(𝐱,t)}}=\dot{\mathrm{\Phi }}(𝐱,t)`$ (157)
are conjugate momenta.
The field and momentum are Fourier-decomposed as
$`\mathrm{\Phi }(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\varphi _𝐤(t)e^{i𝐤𝐱}},`$ (158)
$`\mathrm{\Pi }(𝐱,t)`$ $`=`$ $`\dot{\mathrm{\Phi }}^{}(𝐱,t)={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\dot{\varphi }_𝐤^{}(t)e^{i𝐤𝐱}}{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\pi _𝐤(t)e^{i𝐤𝐱}},`$ (159)
and their conjugates as
$`\mathrm{\Phi }^{}(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\varphi _𝐤^{}(t)e^{i𝐤𝐱}},`$ (160)
$`\mathrm{\Pi }^{}(𝐱,t)`$ $`=`$ $`\dot{\mathrm{\Phi }}(𝐱,t)={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\dot{\varphi }_𝐤(t)e^{i𝐤𝐱}}{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\pi _𝐤^{}(t)e^{i𝐤𝐱}}.`$ (161)
So space integrals of quadratic fields and momenta result in momentum integrals for the decoupled modes
$`{\displaystyle d^3𝐱\mathrm{\Phi }^{}\mathrm{\Phi }}`$ $`=`$ $`{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\varphi _𝐤^{}\varphi _𝐤},`$ (162)
$`{\displaystyle d^3𝐱\mathrm{\Pi }^{}\mathrm{\Pi }}`$ $`=`$ $`{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\pi _𝐤^{}\pi _𝐤},`$ (163)
$`{\displaystyle d^3𝐱\mathrm{\Phi }^{}\mathrm{\Phi }}`$ $`=`$ $`{\displaystyle \frac{d^3𝐤}{(2\pi )^3}𝐤^2\varphi _𝐤^{}\varphi _𝐤}.`$ (164)
One then obtains the Hamiltonian as the sum of infinite number of time-dependent harmonic oscillators
$$H(t)=\frac{d^3𝐤}{(2\pi )^3}\left[\pi _𝐤^{}\pi _𝐤+\left(𝐤^2+m^2(t)\right)\varphi _𝐤^{}\varphi _𝐤\right].$$
(165)
Canonical quantization is prescribed by imposing the commutation relations at equal times
$`[\widehat{\varphi }_𝐤^{}(t),\widehat{\pi }_𝐤(t)]=i\mathrm{}\delta _{𝐤,𝐤^{}},`$ (166)
$`[\widehat{\varphi }_𝐤^{}^{}(t),\widehat{\pi }_𝐤^{}(t)]=i\mathrm{}\delta _{𝐤,𝐤^{}},`$ (167)
and all the other commutators vanish. Following Sec. III, we find the two pairs of the annihilation and creation operators (11) for each $`𝐤`$-mode,
$`\widehat{a}_𝐤(t)`$ $`=`$ $`i\left(\phi _𝐤^{}(t)\widehat{\pi }_𝐤^{}\dot{\phi }_𝐤^{}(t)\widehat{\varphi }_𝐤\right),`$ (168)
$`\widehat{a}_𝐤^{}(t)`$ $`=`$ $`i\left(\phi _𝐤(t)\widehat{\pi }_𝐤\dot{\phi }_𝐤(t)\widehat{\varphi }_𝐤^{}\right),`$ (169)
and
$`\widehat{a}_𝐤^{}(t)`$ $`=`$ $`i\left(\phi _𝐤^{}(t)\widehat{\pi }_𝐤\dot{\phi }_𝐤^{}(t)\widehat{\varphi }_𝐤^{}\right),`$ (170)
$`\widehat{a}_𝐤^{}(t)`$ $`=`$ $`i\left(\phi _𝐤(t)\widehat{\pi }_𝐤^{}\dot{\phi }_𝐤(t)\widehat{\varphi }_𝐤\right),`$ (171)
where $`\phi _𝐤`$ and $`\phi _𝐤^{}`$ satisfy the same classical equation of motion
$$\ddot{\phi }_𝐤(t)+\left(𝐤^2+m^2(t)\right)\phi _𝐤(t)=0.$$
(172)
They further satisfy the standard commutation relations
$$[\widehat{a}_𝐤^{},\widehat{a}_𝐤^{}]=\delta _{𝐤,𝐤^{}},[\widehat{a}_𝐤^{}^{},\widehat{a}_𝐤^{}]=\delta _{𝐤,𝐤^{}}.$$
(173)
The Fock space for each mode can be constructed according to Sec. III. We consider two symmetric states: the vacuum and thermal states. The vacuum is the one annihilated by all the $`\widehat{a}_𝐤`$ and $`\widehat{a}_𝐤^{}`$:
$$\widehat{a}_𝐤(t)|0,t=0,\widehat{a}_𝐤^{}(t)|0,t=0.$$
(174)
By inverting Eqs. (169) and (171) one expresses the fields as
$`\widehat{\varphi }_𝐤=\mathrm{}\left(\phi _𝐤\widehat{a}_𝐤+\dot{\phi _𝐤}^{}\widehat{a}_𝐤^{}\right),`$ (175)
$`\widehat{\varphi }_𝐤^{}=\mathrm{}\left(\phi _𝐤\widehat{a}_𝐤+\dot{\phi _𝐤}^{}\widehat{a}_𝐤^{}\right),`$ (176)
from which follow the vacuum expectation values
$`\widehat{\mathrm{\Phi }}^{}\widehat{\mathrm{\Phi }}_\mathrm{V}={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\right]},`$ (177)
$`\widehat{\mathrm{\Pi }}^{}\widehat{\mathrm{\Pi }}_\mathrm{V}={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\mathrm{}^2\dot{\phi }_𝐤^{}(t)\dot{\phi }_𝐤(t)\right]}.`$ (178)
The initial thermal state defined by the density operator for each mode
$`\widehat{\rho }(t)={\displaystyle \underset{𝐤}{}}\widehat{\rho }_𝐤(t)={\displaystyle \underset{𝐤}{}}\{{\displaystyle \frac{1}{Z_𝐤}}\mathrm{exp}[\beta \mathrm{}\omega _{i,𝐤}(\widehat{a}_𝐤^{}(t)\widehat{a}_𝐤(t)+{\displaystyle \frac{1}{2}})]`$ (179)
$`\times {\displaystyle \frac{1}{Z_𝐤^{}}}\mathrm{exp}[\beta \mathrm{}\omega _{i,𝐤}(\widehat{a}_𝐤^{}(t)\widehat{a}_𝐤^{}(t)+{\displaystyle \frac{1}{2}})]\}`$ (180)
leads to the thermal expectation values
$`\widehat{\mathrm{\Phi }}^{}\widehat{\mathrm{\Phi }}_\mathrm{T}`$ $`=`$ $`\mathrm{Tr}\left[\widehat{\rho }(t)\widehat{\mathrm{\Phi }}^{}\widehat{\mathrm{\Phi }}\right]={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})\right]},`$ (181)
$`\widehat{\mathrm{\Pi }}^{}\widehat{\mathrm{\Pi }}_\mathrm{T}`$ $`=`$ $`\mathrm{Tr}\left[\widehat{\rho }(t)\widehat{\mathrm{\Pi }}^{}\widehat{\mathrm{\Pi }}\right]={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\mathrm{}^2\dot{\phi }_𝐤^{}(t)\dot{\phi }_𝐤(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})\right]}.`$ (182)
One then finds the two-point correlation functions at equal times by taking the expectation value with respect to the vacuum state
$$G_\mathrm{V}(𝐲,𝐱,t)=\widehat{\mathrm{\Phi }}^{}(𝐲,t)\widehat{\mathrm{\Phi }}(𝐱,t)_\mathrm{V}=\frac{d^3k}{(2\pi )^3}\left[\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\right]e^{i𝐤(𝐱𝐲)},$$
(183)
and with respect to the thermal state
$$G_\mathrm{T}(𝐲,𝐱,t)=\widehat{\mathrm{\Phi }}^{}(𝐲,t)\widehat{\mathrm{\Phi }}(𝐱,t)_\mathrm{T}=\frac{d^3k}{(2\pi )^3}\left[\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})\right]e^{i𝐤(𝐱𝐲)},$$
(184)
where
$$\omega _{i,𝐤}=\sqrt{𝐤^2+m^2(\mathrm{})}.$$
(185)
### A Instantaneous Quench
The instantaneous quench model is an analytically solvable one, in which the mass changes as
$`m^2(t)=\{\begin{array}{cc}m_i^2,\hfill & t<0\text{,}\hfill \\ m_f^2,\hfill & t>0\text{.}\hfill \end{array}`$ (186)
Before the quench $`(t<0)`$, the solution to Eq. (172), which also satisfies the condition (14), is given by
$$\phi _{i,𝐤}(t)=\frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}e^{i\omega _{i,𝐤}t},\omega _{i,𝐤}=\sqrt{k^2+m_i^2}.$$
(187)
Then the two-point vacuum correlation function (183) becomes
$`G_{i,\mathrm{V}}(𝐲,𝐱,t)`$ $`=`$ $`{\displaystyle \frac{d^3k}{(2\pi )^3}\frac{\mathrm{}}{2\sqrt{k^2+m_i^2}}e^{i𝐤(𝐱𝐲)}}`$ (188)
$`=`$ $`{\displaystyle \frac{\mathrm{}}{4\pi ^2}}{\displaystyle \frac{m_iK_1(m_i|𝐱𝐲|)}{|𝐱𝐲|}},`$ (189)
where $`K_1`$ is the modified Bessel function. Equation (189) coincides with the result for a massive scalar field in Ref. . Similarly, the two-point thermal correlation function is given by
$`G_{i,\mathrm{T}}(𝐲,𝐱,t)`$ $`=`$ $`G_{i,\mathrm{V}}(𝐲,𝐱,t)`$ (190)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{2\pi ^2}}{\displaystyle \frac{m_i}{\sqrt{|𝐱𝐲|^2+m_i^2}}}{\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}K_1(m_i\sqrt{|𝐱𝐲|^2+(\beta \mathrm{}n)^2)}).`$ (191)
On the other hand, after the quench $`(t>0)`$, the classical equations of motion are classified into two types: the one from the long wavelength modes with $`k^2<m_f^2`$ has the negative frequency squared and exhibits an exponential behavior, and the other from the short wavelength modes with $`k^2>m_f^2`$ still has the positive frequency squared and shows an oscillatory behavior. Each mode moves under a constant frequency squared before the quench time but suddenly experiences a potential step in the case of short wavelengths and a potential barrier in the case of long wavelengths. There is the analogy between Eq. (172) and the scattering problem of quantum mechanics. The solution to Eq. (172) together with the initial asymptotic data (187) is the complex conjugate of the scattering wave function by either the potential step or barrier . The solution to Eq. (172) after the quench should match at the quench time $`t=0`$ continuously with Eq. (187) before the quench. It is rather straightforward to find such solutions for the short wavelength modes $`(k^2>m_f^2)`$
$$\phi _{f_S,𝐤}(t)=\frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}\left[i\frac{\omega _{i,𝐤}}{\omega _{f,𝐤}}\mathrm{sin}(\omega _{f,𝐤}t)+\mathrm{cos}(\omega _{f,𝐤}t)\right],\omega _{f,𝐤}=\sqrt{k^2m_f^2},$$
(192)
and for the long wavelength modes $`(k^2<m_f^2)`$
$$\phi _{f_U,𝐤}(t)=\frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}\left[i\frac{\omega _{i,𝐤}}{\stackrel{~}{\omega }_{f,𝐤}}\mathrm{sinh}(\stackrel{~}{\omega }_{f,𝐤}t)+\mathrm{cosh}(\stackrel{~}{\omega }_{f,𝐤}t)\right],\stackrel{~}{\omega }_{f,𝐤}=\sqrt{m_f^2k^2}.$$
(193)
A few comments are in order. The solution (193) represents an instability due to the phase transition and is obtained by continuing analytically the solution (192). When Eq. (192) is rewritten as
$$\phi _{f,𝐤}(t)=\frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}\left[\left(\frac{\omega _{f,𝐤}+\omega _{i,𝐤}}{2\omega _{f,𝐤}}\right)e^{i\omega _{f,𝐤}t}+\left(\frac{\omega _{f,𝐤}\omega _{i,𝐤}}{2\omega _{f,𝐤}}\right)e^{i\omega _{f,𝐤}t}\right],$$
(194)
the first and second terms correspond to the positive and negative frequencies, respectively, hence the second term explains the particle creation by changing frequency .
After some manipulations of algebra, we obtain the two-point vacuum correlation function after the quench
$`G_{f,\mathrm{V}}(𝐲,𝐱,t)`$ $`=`$ $`G_{i,\mathrm{V}}(𝐲,𝐱,t)`$ (195)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4\pi ^2|𝐱𝐲|}}{\displaystyle _0^{m_f}}𝑑kk\left({\displaystyle \frac{\omega _{i,𝐤}^2+\stackrel{~}{\omega }_{f,𝐤}^2}{\omega _{i,𝐤}\stackrel{~}{\omega }_{f,𝐤}^2}}\right)\mathrm{sin}(k|𝐱𝐲|)\mathrm{sinh}^2(\stackrel{~}{\omega }_{f,𝐤}t)`$ (196)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4\pi ^2|𝐱𝐲|}}{\displaystyle _{m_f}^{\mathrm{}}}𝑑kk\left({\displaystyle \frac{\omega _{i,𝐤}^2\omega _{f,𝐤}^2}{\omega _{i,𝐤}\omega _{f,𝐤}^2}}\right)\mathrm{sin}(k|𝐱𝐲|)\mathrm{sin}^2(\omega _{f,𝐤}t).`$ (197)
Similarly, the two-point thermal correlation function is given by
$`G_{f,\mathrm{T}}(𝐲,𝐱,t)`$ $`=`$ $`G_{i,\mathrm{T}}(𝐲,𝐱,t)`$ (198)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4\pi ^2|𝐱𝐲|}}{\displaystyle _0^{m_f}}𝑑kk\left({\displaystyle \frac{\omega _{i,𝐤}^2+\stackrel{~}{\omega }_{f,𝐤}^2}{\omega _{i,𝐤}\stackrel{~}{\omega }_{f,𝐤}^2}}\right)\mathrm{sin}(k|𝐱𝐲|)\mathrm{sinh}^2(\stackrel{~}{\omega }_{f,𝐤}t)\mathrm{coth}\left({\displaystyle \frac{\beta \mathrm{}\omega _{i,𝐤}}{2}}\right)`$ (199)
$`+`$ $`{\displaystyle \frac{\mathrm{}}{4\pi ^2|𝐱𝐲|}}{\displaystyle _{m_f}^{\mathrm{}}}𝑑kk\left({\displaystyle \frac{\omega _{i,𝐤}^2\omega _{f,𝐤}^2}{\omega _{i,𝐤}\omega _{f,𝐤}^2}}\right)\mathrm{sin}(k|𝐱𝐲|)\mathrm{sin}^2(\omega _{f,𝐤}t)\mathrm{coth}\left({\displaystyle \frac{\beta \mathrm{}\omega _{i,𝐤}}{2}}\right).`$ (200)
The first terms in Eqs. (197) and (200) are the two-point vacuum and thermal correlations (189) and (191), respectively, before the quench. Therefore the remaining two terms describe the effect of the quench. In particular, the second terms are dominant and rooted on the instability during the phase transition, which is missing in the field theoretical approach to equilibrium dynamics. Note that $`\omega _{i,𝐤}^2\omega _{f,𝐤}^2=m_i^2+m_f^2`$ and $`\omega _{i,𝐤}^2+\stackrel{~}{\omega }_{f,𝐤}^2=m_i^2+m_f^2`$, so the amplitudes of $`\mathrm{sin}(k|𝐱𝐲|)`$ decrease as $`1/k^3`$ for very short wavelengths, hence short wavelengths contribute negligibly. However, there is a residual contribution from near the critical wavelength $`k_c=m_f`$, which becomes much smaller than the second terms at later times and will not be considered any more.
We wish to determine the size of domains from the second order phase transition of the instantaneous quench. At later times $`(m_ft1)`$ after the quench the dominant contribution to Eq. (200) comes from the second term, so one has approximately
$$G_{f_U,\mathrm{T}}(r,t)\frac{\mathrm{}}{16\pi ^2r}_0^{m_f}𝑑k\left\{ke^{2\stackrel{~}{\omega }_{f,𝐤}t}\right\}\mathrm{sin}(rk)F_\mathrm{I}(k),$$
(201)
where $`r=|𝐱𝐲|`$, and
$$F_\mathrm{I}(k)=\left(\frac{m_i^2+m_f^2}{\omega _{i,𝐤}\stackrel{~}{\omega }_{f,𝐤}^2}\right)\mathrm{coth}\left(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2}\right).$$
(202)
The function of $`k`$ in the curly bracket of Eq. (201) has a sharp peak at $`k_0=\sqrt{m_f/2t}`$, whereas $`F_\mathrm{I}(k)`$ is a slowly varying function. We employ the steepest descent method (see Appendix C) to obtain
$$G_{f_U,\mathrm{T}}(r,t)G_{f_U,\mathrm{T}}(0,t)\frac{\mathrm{sin}\left(\sqrt{\frac{m_f}{2t}}r\right)}{\sqrt{\frac{m_f}{2t}r}}\mathrm{exp}\left[\frac{m_fr^2}{8t}\right],$$
(203)
where
$$G_{f_U,\mathrm{T}}(0,t)=\frac{\mathrm{}}{16\pi ^2}\sqrt{\frac{\pi }{2e}}\left(\frac{m_f}{2t}\right)^{3/2}e^{2m_ft}F_\mathrm{I}\left(k_0=\sqrt{\frac{m_f}{2t}}\right).$$
(204)
Therefore the size of domains grows according to the classical Cahn-Allen scaling relation
$$\xi _D(t)=\sqrt{\frac{8t}{m_f}}.$$
(205)
The scaling relation (205) for the instantaneous quench confirms the result obtained in Refs. .
### B Finite Smooth Quench
The instantaneous quench does not exhibit the essential spinodal behavior during the quench process. To see the dynamics of the second order phase transition one needs a finite quench period. Such a finite and smooth quench model is described by a field with the mass given by
$$m^2(t)=m_1^2m_0^2\mathrm{tanh}\left(\frac{t}{\tau }\right),(m_0^2>|m_1^2|).$$
(206)
At earlier times $`t=\mathrm{}`$, the mass has the initial value
$$m^2=m_i^2=m_0^2m_1^2>0,$$
(207)
and at later times $`t=\mathrm{}`$, the final value
$$m^2=m_f^2=(m_0^2+m_1^2)<0.$$
(208)
Here $`\tau `$ measures the quench rate: the large $`\tau `$-limit implies that the mass changes slowly from $`m_i^2`$ at $`t=\mathrm{}`$ to $`m_f^2`$ at $`t=+\mathrm{}`$, whereas the small $`\tau `$-limit implies a rapid change of the mass. In particular, the $`(\tau =0)`$-limit corresponds to the instantaneous change from $`m_i^2`$ to $`m_f^2`$ at $`t=0`$. That is, the instantaneous quench is a special case of the finite smooth quench model. To find the Fock space for each mode one needs to solve the classical equation of motion
$$\ddot{\phi }_𝐤(t)+\left(𝐤^2m_1^2m_0^2\mathrm{tanh}\left(\frac{t}{\tau }\right)\right)\phi _𝐤(t)=0.$$
(209)
It should be noted that, as in the instantaneous quench model, long wavelength modes $`(kk_c=m_f)`$ let the frequency change the sign at later times $`(t\tau )`$
$$\omega _𝐤^2(t)=𝐤^2m_1^2m_0^2<0,$$
(210)
and suffer from the spinodal instability. Each long wavelength mode has a different quench time determined by $`\omega _𝐤(t_𝐤)=0`$.
The solutions to Eq. (209) are found separately for the stable modes and unstable modes. The stable modes $`(km_f)`$ have the solutions
$$\phi _𝐤(t)=C_𝐤e^{2p_𝐤t}{}_{2}{}^{}F_{1}^{}(\beta _{+,𝐤},\beta _{,𝐤};\gamma _𝐤;e^{2t/\tau }),$$
(211)
where
$`p_𝐤`$ $`=`$ $`i{\displaystyle \frac{1}{2}}\omega _{i,𝐤},`$ (212)
$`\beta _{\pm ,𝐤}`$ $`=`$ $`i{\displaystyle \frac{\tau }{2}}\left(\omega _{i,𝐤}\pm \omega _{f,𝐤}\right),`$ (213)
$`\gamma _𝐤`$ $`=`$ $`1i\tau \omega _{i,𝐤},`$ (214)
with
$$\omega _{i,𝐤}=\sqrt{k^2+m_i^2},\omega _{f,𝐤}=\sqrt{k^2m_f^2}.$$
(215)
Whereas the unstable modes $`(k<m_f)`$ have the solutions
$$\phi _𝐤(t)=C_𝐤e^{2p_𝐤t}{}_{2}{}^{}F_{1}^{}(\stackrel{~}{\beta }_{+,𝐤},\stackrel{~}{\beta }_{,𝐤};\gamma _𝐤;e^{2t/\tau }),$$
(216)
where
$$\stackrel{~}{\beta }_{\pm ,𝐤}=\frac{\tau }{2}\left(i\omega _{i,𝐤}\pm \stackrel{~}{\omega }_{f,𝐤}\right),$$
(217)
with
$$\stackrel{~}{\omega }_{f,𝐤}=\sqrt{m_f^2k^2}.$$
(218)
At earlier times $`(\tau \tau )`$ before the quench begins, both the solutions (211) and (216) have the same asymptotic form
$$\phi _{i,𝐤}(t)=C_𝐤e^{i\omega _{i,𝐤}t},$$
(219)
so the constant is fixed to satisfy Eq. (14)
$$c_𝐤=\frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}.$$
(220)
#### 1 During the Quench
During the quench process $`(|t|<\tau )`$, the asymptotic forms for the solutions (211) and (216) are unfortunately not available. Instead, one may expand the mass (206) to the linear order
$$m^2(t)=m_1^2m_0^2\left(\frac{t}{\tau }\right),$$
(221)
which is a good approximation as far as $`|t|\tau `$. But we assume $`m_0\tau 1`$ so that the linear quench process is sufficiently long enough to allow the asymptotic analysis. Each mode of the scalar field then has the frequency squared
$$\omega _𝐤^2(t)=𝐤^2m_1^2m_0^2\left(\frac{t}{\tau }\right)=m_0^2\left(\frac{t_𝐤t}{\tau }\right),$$
(222)
where $`t_𝐤`$ is the quench time for the corresponding unstable mode
$$t_𝐤=\frac{\tau }{m_0^2}(𝐤^2m_1^2).$$
(223)
Unless $`|t_𝐤/\tau |1`$, the quench time $`t_𝐤`$ occurs outside the valid regime for Eq. (221). So we restrict our attention to those unstable modes with $`t_𝐤\tau `$. The frequency (222) corresponds to the special case of Eq. (91), in which $`l_1=l_2=0`$, $`t_0=t_𝐤`$, $`t_0t_i=\tau `$ and $`\omega _i=m_0`$. Then the solution in the intermediate regime before the quench, matching with the initial solution (219), is given by Eq. (92):
$$\phi _{m_S,𝐤}(t)=D_1z_𝐤^{1/3}H_{1/3}^{(2)}(z_𝐤)+D_2z_𝐤^{1/3}H_{1/3}^{(1)}(z_𝐤),$$
(224)
where
$$z_𝐤=\frac{2}{3}m_0\tau \left(\frac{t_𝐤t}{\tau }\right)^{1/3}.$$
(225)
The coefficients $`D_1`$ and $`D_2`$ are given by Eq. (95). After the quench time $`(t>t_𝐤)`$ for each mode, the solution (224) is analytically continued for the unstable mode
$`\phi _{m_U,𝐤}(t)={\displaystyle \frac{D_1}{2}}\stackrel{~}{z}_𝐤^{1/3}\left[e^{i\frac{3}{2}\pi }H_{1/3}^{(2)}(i\stackrel{~}{z}_𝐤)+e^{i\frac{1}{2}\pi }H_{1/3}^{(2)}(i\stackrel{~}{z}_𝐤)\right]`$ (226)
$`+{\displaystyle \frac{D_2}{2}}\stackrel{~}{z}_𝐤^{1/3}\left[e^{i\frac{3}{2}\pi }H_{1/3}^{(1)}(i\stackrel{~}{z}_𝐤)+e^{i\frac{1}{2}\pi }H_{1/3}^{(1)}(i\stackrel{~}{z}_𝐤)\right],`$ (227)
where
$$\stackrel{~}{z}_𝐤=\frac{2}{3}m_0\tau \left(\frac{tt_𝐤}{\tau }\right)^{1/3}.$$
(228)
The two-point thermal correlation function
$$G_{m,\mathrm{T}}(r,t)=\frac{\mathrm{}^2}{2\pi ^2}_0^{m_f}𝑑kk^2\frac{\mathrm{sin}(kr)}{kr}\phi _{m_U,𝐤}^{}(t)\phi _{m_U,𝐤}(t)+\frac{\mathrm{}^2}{2\pi ^2}_{m_f}^{\mathrm{}}𝑑kk^2\frac{\mathrm{sin}(kr)}{kr}\phi _{m_S,𝐤}^{}(t)\phi _{m_S,𝐤}(t),$$
(229)
with $`r=|𝐱𝐲|`$, is dominated by the unstable modes (227), since the stable modes (224) oscillate rapidly and do contribute little. By using the asymptotic form for the solution (227) in the regime $`\tau t(\tau /m_0^2)^{1/3}`$
$$\phi _{m_U,𝐤}(t)=\sqrt{\frac{1}{2\pi }}\frac{e^{\stackrel{~}{z}_𝐤}}{\stackrel{~}{z}_𝐤^{1/6}}\left[e^{i\frac{1}{3}\pi }D_1+e^{i\frac{2}{3}\pi }D_2\right].$$
(230)
one has approximately
$$\phi _{m_U,𝐤}^{}(t)\phi _{m_U,𝐤}(t)=\frac{1}{8\mathrm{}}\left\{\frac{1}{m_0}\left(\frac{z_i}{\stackrel{~}{z}_𝐤}\right)^{1/3}\mathrm{sin}^2\left(z_i+\frac{\pi }{4}\right)\right\}e^{2\stackrel{~}{z}_𝐤},$$
(231)
where $`z_i=2m_0\tau /3`$. So the correlation function takes the form
$$G_{m_U,\mathrm{T}}(r,t)\frac{\mathrm{}}{16\pi ^2r}_0^{m_f}𝑑k\left\{ke^{2\stackrel{~}{\omega }_{f,𝐤}\stackrel{~}{t}}\right\}\mathrm{sin}(kr)F_{\mathrm{II}}(k),$$
(232)
where
$$F_{\mathrm{II}}(k)=\frac{1}{m_0}\left(\frac{\frac{2}{3}m_0\tau }{\stackrel{~}{z}_𝐤}\right)^{1/3}\mathrm{sin}^2\left(\frac{2}{3}m_0\tau +\frac{\pi }{4}\right)\mathrm{coth}\left(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2}\right).$$
(233)
The function in the curly bracket is rapidly varying and has a peak at $`k_0=(m_0/2\sqrt{\tau t})^{1/2}`$. Hence after applying the steepest decent method (see Appendix C), we finally obtain
$$G_{m_U,\mathrm{T}}(r,t)G_{m_U,\mathrm{T}}(0,t)\frac{\mathrm{sin}\left(\frac{\sqrt{\tau t}}{m_0}r\right)}{\frac{\sqrt{\tau t}}{m_0}r}\mathrm{exp}\left[\frac{r^2}{8\frac{\sqrt{\tau t}}{m_0}}\right],$$
(234)
where
$$G_{m_U,\mathrm{T}}(0,t)=\frac{\mathrm{}}{64\pi ^2m_0}\left(\frac{\pi m_0^3}{(\tau t)^{3/2}}\right)^{1/2}e^{(\frac{4}{3}m_0t+2\frac{m_1^2}{m_0}\tau )\sqrt{\frac{t}{\tau }}}F_{II}\left(k_0=(\frac{m_0}{2\sqrt{\tau t}})^{1/2}\right).$$
(235)
We thus have shown the scaling relation for the domain size
$$\xi _D(t)=2\left(\frac{2\tau t}{m_0^2}\right)^{1/4}.$$
(236)
The scaling relation (236) confirms, up to a numerical factor, the result in Ref. . However, the power-law is different from the Cahn-Allen scaling relation after the completion of quench.
#### 2 After the Quench
At later times $`(t\tau )`$ after the completion of quench, the solution (211) has the asymptotic form
$`\phi _{f,𝐤}(t)={\displaystyle \frac{e^{2p_𝐤t}}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}}[{\displaystyle \frac{\mathrm{\Gamma }(\gamma _𝐤)\mathrm{\Gamma }(\beta _{,𝐤}\beta _{+,𝐤})}{\mathrm{\Gamma }(\beta _{,𝐤})\mathrm{\Gamma }(\gamma _𝐤\beta _{+,𝐤})}}e^{2\beta _{+,𝐤}t/\tau }{}_{2}{}^{}F_{1}^{}(\beta _{+,𝐤},1\gamma _𝐤+\beta _{+,𝐤};1\beta _{,𝐤}+\beta _{+,𝐤};e^{2t/\tau })`$ (237)
$`+{\displaystyle \frac{\mathrm{\Gamma }(\gamma _𝐤)\mathrm{\Gamma }(\beta _{+,𝐤}\beta _{,𝐤})}{\mathrm{\Gamma }(\beta _{+,𝐤})\mathrm{\Gamma }(\gamma _𝐤\beta _{,𝐤})}}e^{2\beta _{,𝐤}t/\tau }{}_{2}{}^{}F_{1}^{}(\beta _{,𝐤},1\gamma _𝐤+\beta _{,𝐤};1\beta _{+,𝐤}+\beta _{,𝐤};e^{2t/\tau })].`$ (238)
The asymptotic form for Eq. (216) is obtained by replacing $`\beta _{\pm ,𝐤}`$ by $`\stackrel{~}{\beta }_{\pm ,𝐤}`$. From the asymptotic form of the hypergeometric function , we find the asymptotic form for the stable modes
$`\phi _{f_S,𝐤}(t)={\displaystyle \frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}}{\displaystyle \frac{\mathrm{\Gamma }(1i\omega _{i,𝐤}\tau )\mathrm{\Gamma }(i\omega _{f,𝐤}\tau )}{\mathrm{\Gamma }(1i\frac{\tau }{2}(\omega _{i,𝐤}+\omega _{f,𝐤}))\mathrm{\Gamma }(i\frac{\tau }{2}(\omega _{i,𝐤}+\omega _{f,𝐤}))}}e^{i\omega _{f,𝐤}t}`$ (239)
$`+{\displaystyle \frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}}{\displaystyle \frac{\mathrm{\Gamma }(1i\omega _{i,𝐤}\tau )\mathrm{\Gamma }(i\omega _{f,𝐤}\tau )}{\mathrm{\Gamma }(1i\frac{\tau }{2}(\omega _{i,𝐤}\omega _{f,𝐤}))\mathrm{\Gamma }(i\frac{\tau }{2}(\omega _{i,𝐤}\omega _{f,𝐤}))}}e^{i\omega _{f,𝐤}t},`$ (240)
and for the unstable modes
$`\phi _{f_U,𝐤}(t)={\displaystyle \frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}}{\displaystyle \frac{\mathrm{\Gamma }(1i\omega _{i,𝐤}\tau )\mathrm{\Gamma }(\stackrel{~}{\omega }_{f,𝐤}\tau )}{(1)\frac{\tau }{2}(i\omega _{i,𝐤}\stackrel{~}{\omega }_{f,𝐤})\mathrm{\Gamma }^2(\frac{\tau }{2}(i\omega _{i,𝐤}\stackrel{~}{\omega }_{f,𝐤}))}}e^{\stackrel{~}{\omega }_{f,𝐤}t}`$ (241)
$`+{\displaystyle \frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}}{\displaystyle \frac{\mathrm{\Gamma }(1i\omega _{i,𝐤}\tau )\mathrm{\Gamma }(\stackrel{~}{\omega }_{f,𝐤}\tau )}{(1)\frac{\tau }{2}(i\omega _{i,𝐤}+\stackrel{~}{\omega }_{f,𝐤})\mathrm{\Gamma }^2(\frac{\tau }{2}(i\omega _{i,𝐤}+\stackrel{~}{\omega }_{f,𝐤}))}}e^{\stackrel{~}{\omega }_{f,𝐤}t}.`$ (242)
Now a few comments are in order. First, the coefficient of the positive frequency asymptotic solution for each short wavelength mode
$$\phi _𝐤^{out}(t)=\frac{1}{\sqrt{2\mathrm{}\omega _{f,𝐤}}}e^{i\omega _{f,𝐤}t}$$
(243)
leads to the rate for the initial vacuum to remain in the final vacuum
$$|0_𝐤,+\mathrm{}|0_𝐤,\mathrm{}|^2=\frac{\mathrm{sinh}^2\left[\frac{\pi \tau }{2}\left(\omega _{i,𝐤}+\omega _{f,𝐤}\right)\right]}{\mathrm{sinh}(\pi \tau \omega _{i,𝐤})\mathrm{sinh}(\pi \tau \omega _{f,𝐤})}.$$
(244)
On the other hand, the coefficient of the negative frequency solution $`\phi _𝐤^{out}(t)`$ leads to the particle production rate
$$1|0_𝐤,+\mathrm{}|0_𝐤,\mathrm{}|^2=\frac{\mathrm{sinh}^2\left[\frac{\pi \tau }{2}\left(\omega _{i,𝐤}\omega _{f,𝐤}\right)\right]}{\mathrm{sinh}(\pi \tau \omega _{i,𝐤})\mathrm{sinh}(\pi \tau \omega _{f,𝐤})}.$$
(245)
Second, when $`\stackrel{~}{\omega }_{f𝐤}\tau 1`$, the coefficient of the decaying mode in Eq. (242) can become infinite at
$$\stackrel{~}{\omega }_{f,𝐤}\tau =n,(n=1,2,3,\mathrm{}),$$
(246)
because the gamma function has simple poles at these negative integers. For a rapid quench $`\tau 1`$, there does not exist any integers that satisfy Eq. (246). Hence this kind of resonance can happen only for a non-negligible $`\tau `$, i.e., for a very slow quench, which will be treated elsewhere. Contrary to the resonance at (246), the apparent singularity at $`\stackrel{~}{\omega }_{f,𝐤}\tau =0`$ is removed by considering both terms in Eq. (242):
$`\phi _{f_U𝐤}(t)`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}}{\displaystyle \frac{\mathrm{\Gamma }(1i\omega _{i,𝐤}\tau )}{(i\frac{\tau }{2}\omega _{i,𝐤})\mathrm{\Gamma }^2(i\frac{\tau }{2}\omega _{i,𝐤})}}\left[\mathrm{\Gamma }(\stackrel{~}{\omega }_{f,𝐤}\tau )e^{\stackrel{~}{\omega }_{f,𝐤}t}+\mathrm{\Gamma }(\stackrel{~}{\omega }_{f,𝐤}\tau )e^{\stackrel{~}{\omega }_{f,𝐤}t}\right]`$ (247)
$`=`$ $`{\displaystyle \frac{1}{\sqrt{2\mathrm{}\omega _{i,𝐤}}}}{\displaystyle \frac{\mathrm{\Gamma }(1i\omega _{i,𝐤}\tau )}{(i\frac{\tau }{2}\omega _{i,𝐤})\mathrm{\Gamma }^2(i\frac{\tau }{2}\omega _{i,𝐤})}}\left[2{\displaystyle \frac{\mathrm{sinh}(\stackrel{~}{\omega }_{f,𝐤}t)}{\stackrel{~}{\omega }_{f,𝐤}\tau }}\right].`$ (248)
After the completion of quench $`(t\tau )`$, the unstable modes (242) dominate the correlation function (229) over the stable modes (240). In particular, the first term of Eq. (242) grows exponentially, so by using the asymptotic form in Appendix D, one has approximately
$$\phi _{f_U,𝐤}^{}(t)\phi _{f_U,𝐤}(t)=\frac{1}{8\mathrm{}}\left\{\left(\frac{\omega _{i,𝐤}^2+\stackrel{~}{\omega }_{f,𝐤}^2}{\omega _{i,𝐤}\stackrel{~}{\omega }_{f,𝐤}}\right)\frac{\pi \omega _{i,𝐤}\tau }{\mathrm{sinh}(\pi \omega _{i,𝐤})}\left[\frac{\left(1+\frac{\tau }{2}\stackrel{~}{\omega }_{f,𝐤}\right)^2+\frac{\tau ^2}{4}\omega _{i,𝐤}}{1+\tau \stackrel{~}{\omega }_{f,𝐤}}\right]^2e^{\frac{\tau ^2}{2}(\zeta (2)1)(\omega _{i,𝐤}^2+\stackrel{~}{\omega }_{f,𝐤}^2)}\right\}e^{2\stackrel{~}{\omega }_{f,𝐤}\stackrel{~}{t}},$$
(249)
where $`\zeta (n)`$ is the Riemann zeta function and
$$\stackrel{~}{t}=t\frac{\tau ^3}{8}\left(\zeta (3)1\right)\left(\omega _{i,𝐤}^2+\stackrel{~}{\omega }_{f,𝐤}^2\right).$$
(250)
Note that $`\omega _{i,𝐤}^2+\stackrel{~}{\omega }_{f,𝐤}^2=m_i^2+m_f^2`$, so $`\stackrel{~}{t}`$ lags by a constant in proportion to the cubic power of the quench period $`\tau `$. This time-lag is determined not only by the quench period but also by the initial and final coupling constants $`m_i`$ and $`m_f`$. After applying the steepest decent method to the correlation function (see Appendix C), we finally obtain
$$G_{f_U,\mathrm{T}}(r,t)G_{f_U,\mathrm{T}}(0,t)\frac{\mathrm{sin}\left(\sqrt{\frac{m_f}{2\stackrel{~}{t}}}r\right)}{\sqrt{\frac{m_f}{2\stackrel{~}{t}}r}}\mathrm{exp}\left[\frac{m_fr^2}{8\stackrel{~}{t}}\right],$$
(251)
where
$$G_{f_U,\mathrm{T}}(0,t)=\frac{\mathrm{}}{16\pi ^2}\sqrt{\frac{\pi }{2e}}\left(\frac{m_f}{2\stackrel{~}{t}}\right)^{3/2}e^{2m_f\stackrel{~}{t}}F_{\mathrm{III}}\left(k_0=\sqrt{\frac{m_f}{2\stackrel{~}{t}}}\right).$$
(252)
Here
$$F_{\mathrm{III}}(k)=\left(\frac{\omega _{i,𝐤}^2+\stackrel{~}{\omega }_{f,𝐤}^2}{\omega _{i,𝐤}\stackrel{~}{\omega }_{f,𝐤}}\right)\frac{\pi \omega _{i,𝐤}\tau }{\mathrm{sinh}(\pi \omega _{i,𝐤})}\left[\frac{\left(1+\frac{\tau }{2}\stackrel{~}{\omega }_{f,𝐤}\right)^2+\frac{\tau ^2}{4}\omega _{i,𝐤}}{1+\tau \stackrel{~}{\omega }_{f,𝐤}}\right]^2e^{\frac{\tau ^2}{2}(\zeta (2)1)(\omega _{i,𝐤}^2+\stackrel{~}{\omega }_{f,𝐤}^2)}\mathrm{coth}\left(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2}\right).$$
(253)
Remarkably, the scaling relation of the domain size has the same form as Eq. (205) from the instantaneous quench
$$\xi _D(t)=\sqrt{\frac{8\stackrel{~}{t}}{m_f}}.$$
(254)
However, there is a definite time-lag due to the finite quench as claimed in Ref. . The scaling relation (254) is robust, because only the asymptotic form of the exact solutions (216) is used.
## VII Back-Reaction in $`U(1)`$-Theory for Phase Transitions
The free scalar field model in Sec. VI does not describe a real system for the second order phase transition because the spinodal instability continues indefinitely. This model describes more appropriately an intermediate process of phase transition toward the spinodal line. In this section we consider $`|\mathrm{\Phi }|^4`$-theory, in which there is a natural exit for the spinodal decomposition. The spinodal decomposition ends at the spinodal line, the regime in which the instability from the quench competes with the back-reaction from the $`|\mathrm{\Phi }|^4`$-potential term.
The $`U(1)`$-model for the second order phase transition is described by the Lagrangian density
$$(t)=\dot{\mathrm{\Phi }}^{}\dot{\mathrm{\Phi }}\mathrm{\Phi }^{}\mathrm{\Phi }\frac{\lambda }{4!}\left(\mathrm{\Phi }^{}\mathrm{\Phi }+\frac{12m^2(t)}{\lambda }\right)^2,$$
(255)
where $`m^2(t)`$ is assumed to take either (186) or (206). Hence, the coupling parameter $`m^2(t)`$ of the quadratic term starts with the positive constant value $`m_i^2`$ at earlier times before the quench, keeps decreasing across zero at the moment of quench $`t_0`$, and finally reaches the negative constant value $`m_f^2`$ at later times after the quench. So the $`\mathrm{\Phi }=0`$ remains the true minimum of the potential until $`t_0`$. However, $`\mathrm{\Phi }=0`$ is no longer the true minimum after $`t_0`$, but becomes a local maximum. The true minimum occurs at $`|\mathrm{\Phi }|^2=12m^2(t)/\lambda `$, and the system undergoes a second order phase transition. When the system is in the initial thermal equilibrium at earlier times, it is invariant under $`\mathrm{\Phi }\mathrm{\Phi }e^{i\theta }`$ and has a global $`U(1)`$-symmetry. As the system undergoes the second order phase transition, the symmetry is broken and there occurs topological defects.
¿From the Lagrangian density (255) follows the Hamiltonian density
$$(𝐱,t)=\mathrm{\Pi }^{}\mathrm{\Pi }+\mathrm{\Phi }^{}\mathrm{\Phi }^{}+m^2(t)\mathrm{\Phi }^{}\mathrm{\Phi }+\frac{\lambda }{4!}(\mathrm{\Phi }^{}\mathrm{\Phi })^2,$$
(256)
where we neglected a time-dependent $`c`$-number term. Since it has been shown in Sec. V that the coherent and coherent-thermal states in the LvN approach give the identical result as the mean-field and Hartree-Fock methods, we first divide the field into a classical background and quantum fluctuations and then quantize the latter according to the LvN approach. The field may be divided into the homogeneous classical background field and quantum fluctuations
$$\mathrm{\Phi }(𝐱,t)=\varphi _c(t)+\mathrm{\Phi }_f(𝐱,t)$$
(257)
such that the classical background is a coherent state of the quantum field and the quantum fluctuations have symmetric states such as the vacuum, number and thermal states:
$$\widehat{\mathrm{\Phi }}=\varphi _c(t),\widehat{\mathrm{\Phi }}_f=0.$$
(258)
The Hamiltonian density is then the sum of the classical background, quantum fluctuations and interactions
$$(𝐱,t)=_c(t)+_f(𝐱,t)+_{int}(𝐱,t)+\delta _{int}(𝐱,t),$$
(259)
where
$`_c(t)`$ $`=`$ $`\pi _c^2+m^2(t)\varphi _c^2+{\displaystyle \frac{\lambda }{4!}}\varphi _c^4,`$ (260)
$`_f(𝐱,t)`$ $`=`$ $`\mathrm{\Pi }_f^{}\mathrm{\Pi }_f+\mathrm{\Phi }_f^{}\mathrm{\Phi }_f+m^2(t)\mathrm{\Phi }_f^{}\mathrm{\Phi }_f+{\displaystyle \frac{\lambda }{4!}}(\mathrm{\Phi }_f^{}\mathrm{\Phi }_f)^2,`$ (261)
$`_{int}(𝐱,t)`$ $`=`$ $`{\displaystyle \frac{\lambda }{3!}}\varphi _c^2\left(\mathrm{\Phi }_f^{}\mathrm{\Phi }_f\right),`$ (262)
$`\delta _{int}(𝐱,t)`$ $`=`$ $`\pi _c\left(\mathrm{\Pi }_f^{}+\mathrm{\Pi }_f\right)+\left(\varphi _c\right)\left(\mathrm{\Phi }_f^{}+\mathrm{\Phi }_f\right)+m^2(t)\varphi _c\left(\mathrm{\Phi }_f^{}+\mathrm{\Phi }_f\right)`$ (264)
$`+{\displaystyle \frac{\lambda }{4!}}\{\varphi _c^2(\mathrm{\Phi }_f^2+\mathrm{\Phi }_f^2)+2\varphi _c^3(\mathrm{\Phi }_f^{}+\mathrm{\Phi }_f))+2\varphi _c\left(\mathrm{\Phi }_f^{}\mathrm{\Phi }_f\right)(\mathrm{\Phi }_f^{}+\mathrm{\Phi }_f)\},`$
with $`\pi _c=\dot{\varphi }_c,\pi _f=\dot{\varphi }_f`$. As we are interested in the symmetric state of quantum fluctuations $`\mathrm{\Phi }_f`$, so the expectation value of $`\delta _{int}`$ vanishes and will not be considered any more, since only the terms involving $`(\mathrm{\Phi }_f^{}\mathrm{\Phi }_f)`$ do not vanish when one takes the expectation value with respect to the symmetric state.
Now the field $`\mathrm{\Phi }_f`$ and $`\mathrm{\Phi }_f^{}`$ and their conjugate momenta are decomposed into Fourier-modes according to Eqs. (159) and (161). The Fourier transform of the quartic term leads to
$$d^3𝐱\left(\mathrm{\Phi }_f^{}\mathrm{\Phi }_f\right)^2=\underset{j=1}{\overset{4}{}}\frac{d^3𝐤_j}{(2\pi )^3}\left(\varphi _{𝐤_1}^{}\varphi _{𝐤_2}\right)\left(\varphi _{𝐤_3}^{}\varphi _{𝐤_4}\right)\delta \left(𝐤_1+𝐤_3𝐤_2𝐤_4\right).$$
(265)
The symmetric state further restricts the integral in Eq. (265) to either $`(𝐤_1=𝐤_2,𝐤_3=𝐤_4)`$ or $`(𝐤_1=𝐤_4,𝐤_2=𝐤_3)`$, so Eq. (265) leads to
$$d^3𝐱\left(\widehat{\mathrm{\Phi }}_f^{}\widehat{\mathrm{\Phi }}_f\right)^2=2\left[\frac{d^3𝐤}{(2\pi )^3}\widehat{\varphi }_𝐤^{}\widehat{\varphi }_𝐤\right]^2=2\widehat{\mathrm{\Phi }}_f^{}\widehat{\mathrm{\Phi }}_f^2.$$
(266)
Keeping the symmetric state in mind, we find the the Fourier transform of $`_f+_{int}`$
$`_f+_{int}`$ $`=`$ $`{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\pi _𝐤^{}(t)\pi _𝐤(t)+\left(𝐤^2+m^2(t)+\frac{\lambda }{3!}\varphi _c^2(t)\right)\varphi _𝐤^{}(t)\varphi _𝐤(t)\right]}`$ (268)
$`+{\displaystyle \frac{\lambda }{23!}}\left[{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\varphi _𝐤^{}(t)\varphi _𝐤(t)}\right]^2.`$
Therefore the Hamiltonian for fluctuations is mode-decomposed into the sum of infinite number of coupled anharmonic oscillators.
We apply the method for quantum anharmonic oscillators in Sec. V to the Hamiltonian (268). According to the LvN approach, we first introduce two pairs of the annihilation and creation operators
$`\widehat{A}_𝐤(t)=i\left(\phi _𝐤^{}(t)\widehat{\pi }_𝐤^{}\dot{\phi }_𝐤^{}(t)\widehat{\varphi }_𝐤\right),`$ (269)
$`\widehat{A}_𝐤^{}(t)=i\left(\phi _𝐤(t)\widehat{\pi }_𝐤\dot{\phi }_𝐤(t)\widehat{\varphi }_𝐤^{}\right),`$ (270)
and
$`\widehat{A}_𝐤^{}(t)=i\left(\phi _𝐤^{}(t)\widehat{\pi }_𝐤\dot{\phi }_𝐤^{}(t)\widehat{\varphi }_𝐤^{}\right),`$ (271)
$`\widehat{A}_𝐤^{}(t)=i\left(\phi _𝐤(t)\widehat{\pi }_𝐤^{}\dot{\phi }_𝐤(t)\widehat{\varphi }_𝐤\right).`$ (272)
We then require the creation and annihilation operators (270) to satisfy the LvN equation to obtain the equation of motion
$$\ddot{\phi }_𝐤(t)+\left[𝐤^2+m^2(t)+\frac{\lambda }{3!}\varphi _c^2(t)+\frac{\lambda }{3!}\frac{d^3𝐤_1}{(2\pi )^3}\mathrm{}^2\phi _{𝐤_1}^{}(t)\phi _{𝐤_1}(t)\right]\phi _𝐤(t)=0,$$
(273)
where the expectation value is taken with respect to the vacuum state, a symmetric state, as mentioned. The LvN equations for the operators (272) are the complex conjugate of Eq. (273).
The Fock space for each mode can be found similarly according to Sec. V. We consider two symmetric states: the vacuum and thermal states. The vacuum state that is annihilated by all $`\widehat{A}_𝐤`$ and $`\widehat{A}_𝐤^{}`$
$$\widehat{A}_𝐤(t)|0,t=0,\widehat{A}_𝐤^{}(t)|0,t=0,$$
(274)
leads to the expectation values
$`\widehat{\mathrm{\Phi }}_f^{}\widehat{\mathrm{\Phi }}_f_\mathrm{V}={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\right]},`$ (275)
$`\widehat{\mathrm{\Pi }}_f^{}\widehat{\mathrm{\Pi }}_f_\mathrm{V}={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\mathrm{}^2\dot{\phi }_𝐤^{}(t)\dot{\phi }_𝐤(t)\right]},`$ (276)
The initial thermal state defined by the density operator for each mode
$`\widehat{\rho }_\mathrm{T}(t)={\displaystyle \underset{𝐤}{}}\widehat{\rho }_𝐤(t)={\displaystyle \underset{𝐤}{}}\{{\displaystyle \frac{1}{Z_𝐤}}\mathrm{exp}[\beta \mathrm{}\omega _{i,𝐤}(\widehat{A}_𝐤^{}(t)\widehat{A}_𝐤(t)+{\displaystyle \frac{1}{2}})]`$ (277)
$`\times {\displaystyle \frac{1}{Z_𝐤^{}}}\mathrm{exp}[\beta \mathrm{}\omega _{i,𝐤}(\widehat{A}_𝐤^{}(t)\widehat{A}_𝐤^{}(t)+{\displaystyle \frac{1}{2}})]\}`$ (278)
leads to the expectation values
$`\widehat{\mathrm{\Phi }}_f^{}\widehat{\mathrm{\Phi }}_f_\mathrm{T}`$ $`=`$ $`\mathrm{Tr}\left[\widehat{\rho }_\mathrm{T}(t)\widehat{\mathrm{\Phi }}_f^{}\widehat{\mathrm{\Phi }}_f\right]={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})\right]},`$ (279)
$`\widehat{\mathrm{\Pi }}_f^{}\widehat{\mathrm{\Pi }}_f_\mathrm{T}`$ $`=`$ $`\mathrm{Tr}\left[\widehat{\rho }_\mathrm{T}(t)\widehat{\mathrm{\Phi }}_f^{}\widehat{\mathrm{\Phi }}_f\right]={\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\mathrm{}^2\dot{\phi }_𝐤^{}(t)\dot{\phi }_𝐤(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})\right]}.`$ (280)
Then the effective Hamiltonian for the classical background with all the contributions from fluctuations in the initial thermal state is given by
$`H_C(t)`$ $``$ $`\widehat{}_c+\widehat{}_{int}_\mathrm{T}`$ (281)
$`=`$ $`\pi _c^2+\left[m^2(t)+{\displaystyle \frac{\lambda }{3!}}{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})}\right]\varphi _c^2+{\displaystyle \frac{\lambda }{4!}}\varphi _c^4,`$ (282)
and that for fluctuations by
$`H_\mathrm{F}`$ $``$ $`\widehat{}_f+\widehat{}_{int}_\mathrm{T}`$ (283)
$`=`$ $`{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\left[\mathrm{}^2\dot{\phi }_𝐤^{}(t)\dot{\phi }_𝐤(t)+\left(𝐤^2+m^2(t)+\frac{\lambda }{3!}\varphi _c^2(t)\right)\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\right]\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})}`$ (285)
$`+{\displaystyle \frac{\lambda }{23!}}\left[{\displaystyle \frac{d^3𝐤}{(2\pi )^3}\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})}\right]^2.`$
The effective Hamiltonian from the initial vacuum state is the zero-temperature limit $`(\beta \mathrm{})`$, i.e., $`\mathrm{coth}(\beta \mathrm{}\omega _{i,𝐤}/2)1`$.
On the other hand, in the effective action method of Sec. V.B, the $`\phi _𝐤`$ is a parameter that will be determined by the Hamilton equations. By writing $`\phi _𝐤`$ in the polar form
$$\phi _𝐤=\frac{\zeta _𝐤(t)}{\sqrt{\mathrm{}}}e^{i\theta _𝐤},$$
(286)
and by introducing $`p_{\zeta _𝐤}=\dot{\zeta }_𝐤`$, the effective Hamiltonian for the $`𝐤`$-mode (285) can be rewritten as
$$H_{\mathrm{T},𝐤}(t)=\mathrm{}\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})\left[p_{\zeta _𝐤}^2+\omega _𝐤^2(t)\zeta _𝐤^2+\frac{1}{8\zeta _𝐤^2}\right]=\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})H_{\mathrm{V},𝐤}(t).$$
(287)
where
$$\omega _𝐤^2(t)=𝐤^2+m^2(t)+\frac{\lambda }{3!}\varphi _c^2(t)+\frac{\lambda \mathrm{}}{3!}\frac{d^3𝐤}{(2\pi )^3}\zeta _𝐤^2(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2}).$$
(288)
The Hamilton equations are
$`{\displaystyle \frac{d\zeta _𝐤}{dt}}`$ $`=`$ $`{\displaystyle \frac{}{p_{\zeta _𝐤}}}\left[{\displaystyle \frac{H_{\mathrm{T},𝐤}(t)}{\mathrm{}\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})}}\right]=p_{\zeta _𝐤},`$ (289)
$`{\displaystyle \frac{dp_{\zeta _𝐤}}{dt}}`$ $`=`$ $`{\displaystyle \frac{}{\zeta _𝐤}}\left[{\displaystyle \frac{H_{\mathrm{T},𝐤}(t)}{\mathrm{}\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})}}\right]=\omega _𝐤^2(t)\zeta _𝐤+{\displaystyle \frac{1}{4\zeta _𝐤^3}}.`$ (290)
These Hamilton equations are identical to Eq. (273), because the effective action method is equivalent to the LvN approach as shown in Sec. V.B. The Hamilton equations of motion for the classical background are given by
$$\ddot{\varphi }_c(t)+\left[m^2(t)+\frac{\lambda }{3!}\frac{d^3𝐤}{(2\pi )^3}\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})\right]\varphi _c(t)+\frac{\lambda }{23!}\varphi _c^3=0.$$
(291)
The different quantum states given in Sec. V can be chosen to describe the different processes for the phase transition. We assume that the system starts from either the thermal or coherent-thermal states before the quench and evolves according to the functional Schrödinger equation during and after the quench. In this sense the process is completely determined once the initial quantum state is prescribed. In the first case of the symmetric thermal state, each mode has the zero expectation value for the position and momentum, but its dynamics is still governed by Eq. (273) with $`\varphi _c=0`$. In the second case of the coherent-thermal state, the classical field $`\varphi _c`$, which is a coherent state of the homogeneous field, plays the role of an order parameter and is influenced by thermal fluctuations of $`\varphi _𝐤`$. As the first case is a limit of the second one, we shall first focus on the coherent-thermal and then obtain the thermal state result by taking the limit $`\varphi _c(t)=\dot{\varphi }_c(t)=0`$ in the end.
It is very difficult to solve analytically the equations of motion (273) and (291). At best we have to rely on the adiabatic solutions that can be found in some important physical regimes. We assume that far before the quench $`(t\mathrm{})`$ the system starts from a thermal fluctuations around the order parameter with
$$\varphi _c(\mathrm{})0,\dot{\varphi }_c(\mathrm{})0.$$
(292)
As $`\varphi _c`$ is a classical field, the uncertainty principle does not prohibit us from even taking the limit $`\varphi _c(\mathrm{})=\dot{\varphi }_c(\mathrm{})=0`$, which leads to the symmetric thermal state. The initial data for $`\varphi _c`$ take very small values with high probability by some distribution function. For the sake of simplicity we consider the instantaneous quench first. Before the quench time, $`m^2(t)`$ takes the initial constant value $`m_i^2`$ and $`\varphi _c`$ remains close to zero, so $`\omega _𝐤`$ changes very little during the evolution. The solutions to Eq. (273) are approximately given by
$$\phi _{i,𝐤}(t)=\frac{1}{\sqrt{2\mathrm{}\mathrm{\Omega }_{i,𝐤}}}e^{i\mathrm{\Omega }_{i,𝐤}t},$$
(293)
and Eq. (288) yields the gap equation
$$\mathrm{\Omega }_{i,𝐤}^2=m_i^2+k^2+\frac{\lambda \mathrm{}}{3!}\frac{d^3𝐤}{(2\pi )^3}\frac{1}{2\mathrm{\Omega }_{i,𝐤}}\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2}).$$
(294)
The infinite quantity that appears in Eq. (294) will be absorbed by the bare coupling parameters $`m_i^2`$ and $`\lambda `$ to result in the renormalized ones and a finite equation for $`\mathrm{\Omega }_{i,𝐤}^2\text{[44, 12]}`$. The solutions (293) hold till the quench time $`t=0`$.
But after the quench time, $`m^2(t)`$ changes to $`m_f^2`$, and therefore, the long wavelength modes become unstable and the spinodal instability begins. The long wavelength solutions are given by
$$\phi _{f_U,𝐤}(t)=\frac{1}{\sqrt{2\mathrm{}\mathrm{\Omega }_{i,𝐤}}}\left[i\frac{\mathrm{\Omega }_{i,𝐤}}{\stackrel{~}{\mathrm{\Omega }}_{f,𝐤}(t)}\mathrm{sinh}(_0^t\stackrel{~}{\mathrm{\Omega }}_{f,𝐤}(t))+\mathrm{cosh}(_0^t\stackrel{~}{\mathrm{\Omega }}_{f,𝐤}(t))\right],$$
(295)
and the short wavelength solutions by
$$\phi _{f_S,𝐤}(t)=\frac{1}{\sqrt{2\mathrm{}\mathrm{\Omega }_{i,𝐤}}}\left[i\frac{\mathrm{\Omega }_{i,𝐤}}{\mathrm{\Omega }_{f,𝐤}(t)}\mathrm{sin}(_0^t\mathrm{\Omega }_{f,𝐤}(t))+\mathrm{cos}(_0^t\mathrm{\Omega }_{f,𝐤}(t))\right],$$
(296)
where
$`\stackrel{~}{\mathrm{\Omega }}_{f,𝐤}^2(t)`$ $`=`$ $`m_f^2k^2{\displaystyle \frac{\lambda }{3!}}\varphi _c^2`$ (299)
$`{\displaystyle \frac{\lambda \mathrm{}}{3!}}{\displaystyle _0^{k_\mathrm{\Lambda }}}{\displaystyle \frac{d^3𝐤}{(2\pi )^3}}{\displaystyle \frac{\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})}{2\mathrm{\Omega }_{i,𝐤}}}\left[\left\{\left({\displaystyle \frac{\mathrm{\Omega }_{i,𝐤}}{\stackrel{~}{\mathrm{\Omega }}_{f,𝐤}(t)}}\right)^2+1\right\}^2\mathrm{sinh}({\displaystyle _0^t}\stackrel{~}{\mathrm{\Omega }}_{f,𝐤}(t))+1\right]`$
$`{\displaystyle \frac{\lambda \mathrm{}}{3!}}{\displaystyle _{k_\mathrm{\Lambda }}^{\mathrm{}}}{\displaystyle \frac{d^3𝐤}{(2\pi )^3}}{\displaystyle \frac{\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})}{2\mathrm{\Omega }_{i,𝐤}}}\left[\left\{\left({\displaystyle \frac{\mathrm{\Omega }_{i,𝐤}}{\mathrm{\Omega }_{f,𝐤}(t)}}\right)^21\right\}^2\mathrm{sin}({\displaystyle _0^t}\mathrm{\Omega }_{f,𝐤}(t))+1\right],`$
and $`\mathrm{\Omega }_{f,𝐤}^2(t)=\stackrel{~}{\mathrm{\Omega }}_{f,𝐤}^2(t)`$. In contrast with the free scalar model in Sec. VI, not only the duration of the instability but also the band of unstable modes decrease in time due to the back-reaction $`\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)`$ of the self-interaction and $`\varphi _c^2(t)`$ of the classical background. The classical background obeying Eq. (291), though it remained around the initial vacuum under thermal fluctuations before the quench time, begins to roll down from the false vacuum to the true vacuum at first largely due to the exponentially growing unstable modes and then to the self-interaction. The competition between $`m_f^2`$ and $`\widehat{\mathrm{\Phi }}_f^{}\widehat{\mathrm{\Phi }}_f`$ together $`\varphi _c^2`$ determines the spinodal line, beyond which the instability stops and the fluctuations begin to oscillate around the true vacuum. Therefore the self-interacting phase transition model has a natural exit to the spinodal instability .
We now turn to the case of symmetric thermal state. Once the initial condition is prescribed such that $`\varphi _c(\mathrm{})=\dot{\varphi }_c(\mathrm{})=0`$, the classical field $`\varphi _c(t)`$ remains in the false vacuum even during the quench and has the trivial solution $`\varphi _c(t)=0`$ for all times. Hence Eq. (291) is identically satisfied, and the dynamics of the second order phase transition is entirely governed by Eq. (273) for quantum fluctuations, which extends the free scalar field model considered on Sec. VI. Though each mode of quantum fluctuations has the zero expectation value, the Wigner function becomes sharply peaked around its classical trajectory as the quench proceeds . Even without the classical field $`\varphi _c`$ the quantum contribution from the self-interaction in Eq. (273) still prevents the unstable modes from growing indefinitely and provides a natural exit to the spinodal instability. This implies that the quantum dynamics of the second order phase transition is classically correlated and exhibits most of the essential points described by the order parameter under thermal fluctuations.
Finally we comment on the formation process of topological defects. The topological defects formed from the second order phase transition in this paper are domain walls. The correlation of domain walls can be determined by the two-point thermal correlation function
$$G_\mathrm{T}(𝐲,𝐱,t)=\widehat{\mathrm{\Phi }}^{}(𝐲,t)\widehat{\mathrm{\Phi }}(𝐱,t)_\mathrm{T}=\frac{d^3k}{(2\pi )^3}\left[\mathrm{}^2\phi _𝐤^{}(t)\phi _𝐤(t)\mathrm{coth}(\frac{\beta \mathrm{}\omega _{i,𝐤}}{2})\right]e^{i𝐤(𝐱𝐲)},$$
(300)
where $`\omega _{i,𝐤}=\omega _𝐤(t=\mathrm{})`$. After the quench, the two-point correlation function is dominated by the unstable modes $`\phi _{f_U}`$. The scaling behavior of two-point correlation function in Sec. VI holds still before reaching the spinodal line. That is, the size of domains grows according to the power law $`t^{1/4}`$ during the quench and the Cahn-Allen relation $`t^{1/2}`$ after the completion of quench. However, the domains can not grow indefinitely due to the back-reaction of the self-interaction and the classical background. As each unstable mode reaches further the spinodal line, its solution stops exponential growing and oscillates. The behavior of the two-point thermal correlation function changes from that in Sec. VI, and one would expect a different scaling relation for the domain size, probably with a small power than the Cahn-Allen relation.
## VIII Conclusion
In this paper we have elaborated the recently introduced Liouville-von Neumann (LvN) approach to describe properly the time-dependent nonequilibrium systems. The systems interacting directly with environments or undergoing phase transitions are such nonequilibrium systems. These systems are characterized by time-dependent coupling parameters and their true nonequilibrium evolution deviates significantly from the equilibrium one when their coupling parameters differ greatly from their initial values. In this case the systems evolve completely out of equilibrium. For that purpose there have been developed many different methods such as the closed time-path integral method, sometimes in conjunction with the large $`N`$-expansion, mean-field, Hartree-Fock method. The LvN approach developed in this paper is a canonical method that unifies the functional Schrödinger equation for the quantum evolution of pure states and the LvN equation for the quantum description of mixed states of either equilibrium or nonequilibrium. Because the LvN approach shares all the useful techniques with quantum mechanics and quantum many-particle systems, it turns out to be a powerful method for describing time-dependent harmonic oscillators and anharmonic oscillators, and provides a rigorous and systematic method for describing time-dependent phase transitions.
By applying the LvN approach to time-dependent harmonic oscillators, we have found exactly the nonequilibrium quantum evolution evolving from various initial states such as the vacuum, number, coherent and thermal states. In this case the LvN approach is based on two operators, the so-called the annihilation and creation operators, that satisfy the quantum LvN equation, so it is straightforward to construct the Fock space of number states and the density operator according to the standard technique of quantum mechanics. We have thus obtained the density operator in terms of the classical solution, and by using the exact wave functions for number states, have been able to find the explicit form of the density matrix. In particular, the density matrix provides us with a criterion on nonequilibrium vs. equilibrium evolution. Moreover, the LvN approach has been applied to the time-dependent inverted harmonic oscillators, which can be regarded as models for second order phase transitions. For time-dependent anharmonic oscillators, we have found approximately the nonequilibrium evolution of the symmetric Gaussian, coherent and thermal states at the lowest order of the coupling constant of the quartic term. It has been shown that the LvN approach is equivalent to the effective action method and to the mean field or Hartree-Fock method.
Finally we have applied the LvN approach to the systems undergoing the symmetry breaking second order phase transition. In particular, due to the quench the coupling parameters change the sign during the evolution. As field models we have studied a free massive scalar field with an instantaneous and a finite smooth quench. By applying the LvN approach to this symmetry breaking system we have found the two-point vacuum and thermal correlation functions. It has proved that the spinodal instability leads to the $`t^{1/4}`$-scaling relation for domain sizes during the quench and the classical Cahn-Allen relation after the completion of quench. The Cahn-Allen scaling relation confirms the result for the instantaneous quench model in Refs. . One prominent feature of the finite smooth quench model is the time-lag occurring at the cubic power of the quench period in the Cahn-Allen scaling relation after the completion of quench. The inclusion of a self-interacting term shuts off the spinodal instability after crossing the spinodal line and gives rise to a natural exit for the spinodal decomposition. Not treated in detail in this paper is the very slow quench effect, which may show a transient resonance of decaying solution of long wavelength modes and will be addressed in a future research.
###### Acknowledgements.
We would like to thank S. G. Kim for many useful discussions. This work was supported by the KOSEF under Grant No. 1999-2-112-003-5.
## A Harmonic Oscillator Wave Functions
The vacuum state of the time-dependent oscillator is annihilated by $`\widehat{a}(t)`$:
$$\widehat{a}(t)|0,t=0.$$
(A1)
In the coordinate representation
$$\mathrm{\Psi }_0(q,t)=q|0,t,$$
(A2)
Eq. (A1) becomes
$$i\left[u^{}\frac{\mathrm{}}{i}\frac{}{q}m\dot{u}^{}q\right]\mathrm{\Psi }_0(q,t)=0.$$
(A3)
We thus obtain the normalized wave function for the vacuum state
$$\mathrm{\Psi }_0(q,t)=\left(\frac{1}{2\pi \mathrm{}^2u^{}u}\right)^{1/4}\mathrm{exp}\left[\frac{i}{2}\frac{m}{\mathrm{}}\frac{\dot{u}^{}}{u^{}}q^2\right].$$
(A4)
The wave function for the $`n`$th number state is obtained by applying the creation operator $`\widehat{a}^{}(t)`$ $`n`$-times
$$\mathrm{\Psi }_n(q,t)=\frac{1}{\sqrt{n!}}\left(\widehat{a}^{}(t)\right)^n\mathrm{\Psi }_0(q,t).$$
(A5)
By making use of the relation
$$\left(\widehat{a}^{}(t)\right)^n\mathrm{\Psi }(q,t)=\left(\mathrm{}u\right)^ne^{\frac{i}{2}\frac{m}{\mathrm{}}\frac{\dot{u}}{u}q^2}\left(\frac{}{q}\right)^n\left(e^{\frac{i}{2}\frac{m}{\mathrm{}}\frac{\dot{u}}{u}q^2}\mathrm{\Psi }(q,t)\right),$$
(A6)
and the definition of the Hermite polynomial
$$H_n(x)=(1)^ne^{x^2}\left(\frac{d}{dx}\right)^ne^{x^2},$$
(A7)
we obtain the wave function
$$\mathrm{\Psi }_n(q,t)=\left(\frac{1}{2\pi \mathrm{}^2u^{}u}\right)^{1/4}\frac{1}{\sqrt{2^nn!}}\left(\frac{u}{u^{}}\right)^nH_n(x)\mathrm{exp}\left[\frac{i}{2}\frac{m}{\mathrm{}}\frac{\dot{u}^{}}{u^{}}q^2\right],$$
(A8)
where
$$x=\frac{q}{\sqrt{2\mathrm{}^2u^{}u}}.$$
(A9)
Here we have also used the wronskian (14).
## B Density Matrix
In the coordinate representation, the density operator defined by
$$\widehat{\rho }_\mathrm{T}(t)=\frac{1}{Z_N}e^{\beta \mathrm{}\omega _0(\widehat{N}(t)+\frac{1}{2})},$$
(B1)
where $`Z_N`$ is the partition function, becomes
$$\rho _\mathrm{T}(q^{},q,t)=\frac{1}{Z_N}\underset{n=0}{\overset{\mathrm{}}{}}\mathrm{\Psi }_n(q^{},t)\mathrm{\Psi }_n^{}(q,t)e^{\beta \mathrm{}\omega _0(n+\frac{1}{2})}.$$
(B2)
By substituting (A8) into (B2), we obtain
$$\rho _\mathrm{T}(q^{},q,t)=\frac{1}{Z_N}\left(\frac{1}{2\pi \mathrm{}^2u^{}u}\right)^{1/2}\underset{n=0}{\overset{\mathrm{}}{}}\frac{1}{2^nn!}H_n(x^{})H_n(x)e^{\beta \mathrm{}\omega _0(n+\frac{1}{2})}e^{\frac{i}{2}\frac{m}{\mathrm{}}\frac{\dot{u}^{}}{u^{}}q^2\frac{i}{2}\frac{m}{\mathrm{}}\frac{\dot{u}}{u}q^2}.$$
(B3)
Following Kubo’s method , we rewrite the product of Hermite polynomials as
$$H_n(x^{})H_n(x)=\frac{1}{\pi }e^{x^2+x^2}_{\mathrm{}}^{\mathrm{}}_{\mathrm{}}^{\mathrm{}}𝑑z_1𝑑z_2(2iz_1)^n(2iz_2)^ne^{z_1^22ix^{}z_1z_2^22ixz_2}$$
(B4)
and sum over $`n`$ to obtain
$`\rho _\mathrm{T}(q^{},q,t)`$ $`=`$ $`{\displaystyle \frac{1}{Z_N}}\left({\displaystyle \frac{1}{2\pi \mathrm{}^2u^{}u}}\right)^{1/2}{\displaystyle \frac{e^{\frac{\beta \mathrm{}\omega _0}{2}}}{\pi }}\mathrm{exp}\left[(x^2+x^2)+{\displaystyle \frac{i}{2}}{\displaystyle \frac{m}{\mathrm{}}}\left\{{\displaystyle \frac{\dot{u}^{}}{u^{}}}q^2{\displaystyle \frac{\dot{u}}{u}}q^2\right\}\right]`$ (B6)
$`\times {\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}dz_1dz_2\mathrm{exp}[z_1^22ix^{}z_1z_2^22ixz_22z_1z_2e^{\beta \mathrm{}\omega _0}].`$
After doing the integral and using the identity
$$x^2=\frac{i}{2}\frac{m}{\mathrm{}}\left(\frac{\dot{u}}{u}\frac{\dot{u}^{}}{u^{}}\right)q^2,$$
(B7)
we finally obtain
$`\rho _\mathrm{T}(q^{},q,t)`$ $`=`$ $`{\displaystyle \frac{1}{Z_N}}\left[{\displaystyle \frac{1}{4\pi \mathrm{}^2u^{}u\mathrm{sinh}(\beta \mathrm{}\omega _0)}}\right]^{1/2}\mathrm{exp}\left[{\displaystyle \frac{i}{2}}\mathrm{}m{\displaystyle \frac{d}{dt}}\mathrm{ln}(u^{}u)(x^2x^2)\right]`$ (B9)
$`\times \mathrm{exp}\left[{\displaystyle \frac{1}{4}}\left\{(x^{}+x)^2\mathrm{tanh}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}})+(x^{}x)^2\mathrm{coth}({\displaystyle \frac{\beta \mathrm{}\omega _0}{2}})\right\}\right].`$
In the special case of the time-independent oscillator we recover the classical result by Kubo by substituting the complex solution
$$u(t)=\frac{e^{i\omega _0t}}{\sqrt{2\mathrm{}m\omega _0}}$$
(B10)
into Eq. (B9).
## C Steepest Decent Method
The integral appearing in the two-point correlation function has the form
$$I=𝑑xxe^{\gamma x^2}\mathrm{sin}(yx)F(x)$$
(C1)
where $`F(x)`$ is a slowly varying function. Note that $`xe^{\gamma x^2}`$ is a highly peaked function that varies rapidly. We let
$$e^{g(x)}xe^{\gamma x^2}$$
(C2)
where
$$g(x)=\gamma x^2+\mathrm{ln}(x).$$
(C3)
We expand $`g(x)`$ in a Taylor series and truncate it up to the quadratic term around the maximum point $`x_0=\frac{1}{\sqrt{2\gamma }}`$
$$g(x)g(x_0)2\gamma (xx_0)^2,$$
(C4)
and rewrite the integrand as
$`xe^{\gamma x^2}\mathrm{sin}(yx)F(x)`$ $``$ $`x_0e^{\gamma x_0^2}F(x_0)e^{2\gamma (xx_0)^2}{\displaystyle \frac{e^{iyx}e^{iyx}}{2i}}`$ (C5)
$`=`$ $`x_0e^{\gamma x_0^2}F(x_0)\{\mathrm{exp}[2\gamma (xx_0i{\displaystyle \frac{y}{4\gamma }})^2]{\displaystyle \frac{e^{iyx_0}}{2i}}`$ (C7)
$`+\mathrm{exp}[2\gamma (xx_0+i{\displaystyle \frac{y}{4\gamma }})^2]{\displaystyle \frac{(1)e^{iyx_0}}{2i}}\}e^{\frac{y^2}{8\gamma }}.`$
The Gaussian integrals contribute equally, so we obtain
$$I\left(\frac{\pi }{4\gamma }\right)^{1/2}x_0e^{\gamma x_0^2}F(x_0)\mathrm{sin}(yx_0)e^{\frac{y^2}{8\gamma }}.$$
(C8)
## D Asymptotic Form for Unstable Modes in a Finite Quench
To find the contribution to the two-point functions from the unstable growing modes for a finite quench $`(\stackrel{~}{\omega }_{f,𝐤}\tau <1)`$, we need to evaluate $`\mathrm{\Gamma }(\stackrel{~}{\omega }_{f,𝐤})`$ and $`\mathrm{\Gamma }(\frac{\tau }{2}(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤}))`$. These gamma functions are rewritten as
$`\mathrm{\Gamma }(\stackrel{~}{\omega }_{f,𝐤}\tau )`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(1+\stackrel{~}{\omega }_{f,𝐤}\tau )}{\stackrel{~}{\omega }_{f,𝐤}\tau }},`$ (D1)
$`\mathrm{\Gamma }({\displaystyle \frac{\tau }{2}}(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤}))`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Gamma }(1+\frac{\tau }{2}(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤}))}{\frac{\tau }{2}(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤})}}.`$ (D2)
We further make use of the expansion formula
$$\mathrm{ln}\mathrm{\Gamma }(1+z)=\mathrm{ln}(1+z)+(1\gamma )z+\underset{n=2}{\overset{\mathrm{}}{}}(1)^n\left[\zeta (n)1\right]\frac{z^n}{n},(|z|<2),$$
(D3)
where
$$\gamma =\underset{m\mathrm{}}{lim}\left[\underset{k=1}{\overset{m}{}}\frac{1}{k}\mathrm{ln}(m)\right]=0.5772156649\mathrm{}$$
(D4)
is the Euler’s constant and $`\zeta (n)`$ is the Riemann Zeta function
$$\zeta (n)=\underset{k=1}{\overset{\mathrm{}}{}}\frac{1}{k^n}.$$
(D5)
Here
$$z=\stackrel{~}{\omega }_{f,𝐤}\tau ,z=\frac{\tau }{2}(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤}\tau ).$$
(D6)
We now expand the gamma functions up to the cubic power of $`\tau `$:
$`\mathrm{\Gamma }(1+\stackrel{~}{\omega }_{f,𝐤}\tau )`$ $`=`$ $`{\displaystyle \frac{1}{1+\stackrel{~}{\omega }_{f,𝐤}\tau }}\mathrm{exp}[(1\gamma )(\stackrel{~}{\omega }_{f,𝐤}\tau )+{\displaystyle \frac{1}{2}}[\zeta (2)1](\stackrel{~}{\omega }_{f,𝐤}\tau )^2`$ (D7)
$``$ $`{\displaystyle \frac{1}{3}}[\zeta (3)1](\stackrel{~}{\omega }_{f,𝐤}\tau )^3+𝒪(\tau ^4)]`$ (D8)
and
$`\left|\mathrm{\Gamma }(1+{\displaystyle \frac{\tau }{2}}(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤}))\right|^2`$ $`=`$ $`{\displaystyle \frac{1}{|1+\frac{\tau }{2}(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤})|^2}}\mathrm{exp}[(1\gamma )(\stackrel{~}{\omega }_{f,𝐤}\tau )`$ (D9)
$`+`$ $`{\displaystyle \frac{1}{2}}[\zeta (2)1]({\displaystyle \frac{\tau }{2}})^2\left\{(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤})^2+(\stackrel{~}{\omega }_{f,𝐤}+i\omega _{i,𝐤})^2\right\}`$ (D10)
$``$ $`{\displaystyle \frac{1}{3}}[\zeta (3)1]({\displaystyle \frac{\tau }{2}})^3\{(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤})^3+(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤})^3\}+𝒪(\tau ^4)].`$ (D11)
Therefore it follows that
$`\left|{\displaystyle \frac{\mathrm{\Gamma }(1+\stackrel{~}{\omega }_{f,𝐤}\tau )}{\mathrm{\Gamma }^2(1+\frac{\tau }{2}(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤}))}}\right|^2=\left[{\displaystyle \frac{\left(1+\frac{\tau }{2}\stackrel{~}{\omega }_{f,𝐤}\right)^2+\frac{\tau ^2}{4}\omega _{i,𝐤}^2}{1+\stackrel{~}{\omega }_{f,𝐤}\tau }}\right]^2`$ (D12)
$`\times \mathrm{exp}\left[{\displaystyle \frac{1}{2}}[\zeta (2)1]\tau ^2\left(\stackrel{~}{\omega }_{f,𝐤}^2+\omega _{i,𝐤}^2\right){\displaystyle \frac{1}{4}}[\zeta (3)1]\tau ^3\stackrel{~}{\omega }_{f,𝐤}\left(\stackrel{~}{\omega }_{f,𝐤}^2+\omega _{i,𝐤}^2\right)+𝒪(\tau ^4)\right].`$ (D13)
We finally get
$`\left|{\displaystyle \frac{\mathrm{\Gamma }(\stackrel{~}{\omega }_{f,𝐤}\tau )}{\mathrm{\Gamma }^2(\frac{\tau }{2}(\stackrel{~}{\omega }_{f,𝐤}i\omega _{i,𝐤}))}}\right|^2=\left[{\displaystyle \frac{\frac{\tau ^2}{4}\left(\stackrel{~}{\omega }_{f,𝐤}^2+\omega _{i,𝐤}^2\right)}{\stackrel{~}{\omega }_{f,𝐤}\tau }}\right]^2\left[{\displaystyle \frac{\left(1+\frac{\tau }{2}\stackrel{~}{\omega }_{f,𝐤}\right)^2+\frac{\tau ^2}{4}\omega _{i,𝐤}^2}{1+\stackrel{~}{\omega }_{f,𝐤}\tau }}\right]^2`$ (D14)
$`\times \mathrm{exp}\left[{\displaystyle \frac{1}{2}}[\zeta (2)1]\tau ^2\left(\stackrel{~}{\omega }_{f,𝐤}^2+\omega _{i,𝐤}^2\right){\displaystyle \frac{1}{4}}[\zeta (3)1]\tau ^3\stackrel{~}{\omega }_{f,𝐤}\left(\stackrel{~}{\omega }_{f,𝐤}^2+\omega _{i,𝐤}^2\right)+𝒪(\tau ^4)\right].`$ (D15)
The last relation necessary for the two-point correlation function is
$$\left|\mathrm{\Gamma }(1i\omega _{i,𝐤}\tau )\right|^2=\frac{\pi \omega _{i,𝐤}\tau }{\mathrm{sinh}(\pi \omega _{i,𝐤}\tau )}.$$
(D16)
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# DPNU-00-21hep-lat/0005015 Gauge anomaly cancellations in SU(2)L× U(1)Y Electroweak theory on the lattice
## 1 Introduction
The gauge interaction of the Weyl fermions now can be described in the framework of lattice gauge theory. The clue to this development is the construction of gauge covariant and local Dirac operators which solve the Ginsparg-Wilson relation . The Ginsparg-Wilson relation implies an exact chiral symmetry for the Dirac fermion and a gauge-field-dependent chiral projection to the Weyl degrees of freedom .
The functional measure for the Weyl fermion field is defined based on this chiral projection. It leads to a mathematically reasonable definition of the chiral determinant, which generically has the structure as an overlap of two vacua . It has been shown by Lüscher in that for anomaly-free abelian chiral gauge theories, the functional measure for the Weyl fermion fields can be constructed so that the gauge invariance is maintained exactly on the lattice. This issue of the gauge-invariant construction of the functional measure in non-abelian chiral theories has been related to the cohomological classification of a certain topological field which is defined on the four-dimensional lattice plus two continuum dimensions . It has been shown that in all orders in the lattice spacing $`a`$, the topological field has trivial cohomology for anomaly free theories. This problem has also been examined by Suzuki from the point of view of the Wess-Zumino consistency condition and the BRST cohomology in four-dimensions . The gauge anomaly cancellation has been proved for general gauge groups in all powers of gauge potential.<sup>1</sup><sup>1</sup>1 The topological aspect of the non-abelian anomaly for Weyl fermions defined based on the overlap formalism / the Ginsparg-Wilson relation has been examined by D.H. Adams in close relation to the argument of L. Alvarez-Gaumé and P. Ginsparg in the continuum theory . The global SU(2) anomaly has been examined by H. Neuberger and O. Bär and I. Campos in detail. A lattice implementation of the $`\eta `$-invariant and its relation to the effective action for chiral Ginsparg-Wilson fermions has been examined by T. Aoyama and Y.K. in . Non-compact formulation of abelian chiral gauge theories has been considered recently by Neuberger .
The above result for the abelian chiral gauge theories implies that U(1)<sub>Y</sub> hyper-charge chiral gauge theory now can be constructed on the lattice. In this paper, we consider a first step towards the extension of this work to the case of the SU(2)<sub>L</sub>$`\times `$ U(1)<sub>Y</sub> electroweak theory. We examine the exact cancellation of gauge anomalies in the SU(2)<sub>L</sub>$`\times `$ U(1)<sub>Y</sub> electroweak theory, through the cohomological classification of the 4+2-dimensional topological field proposed by Lüscher . Here we will discuss the cancellation of the local anomaly in infinite volume lattice only and leave the issue related to possible global obstructions to the non-perturbative construction of the theory for future study.
The SU(2)<sub>L</sub>$`\times `$U(1)<sub>Y</sub> electroweak theory contains the following fermions as the first generation:
$$\left(\begin{array}{c}\nu _L(x)\\ e_L(x)\end{array}\right)_{Y=\frac{1}{2}},e_R(x)_{Y=1},\left(\begin{array}{c}u_{Li}(x)\\ d_{Li}(x)\end{array}\right)_{Y=\frac{1}{6}},\begin{array}{c}u_{Ri}(x)_{Y=+\frac{2}{3}}\\ d_{Ri}(x)_{Y=\frac{1}{3}}\end{array},$$
(1.1)
where $`i`$ is the color index $`(i=1,2,3)`$. The left-handed leptons and quarks are SU(2)<sub>L</sub> doublets. The right-handed fermions are SU(2)<sub>L</sub> singlet. Taking into account of the color degrees of freedom, there are four doublets. The hyper-charge $`Y`$, which is related to electromagnetic charge $`Q`$ by the Gell-Mann-Nishijima relation,
$$Q=I_3+Y,$$
(1.2)
are assigned as shown above.
There are two types of gauge anomalies in the electroweak theory. The first one is the gauge anomaly associated with the abelian U(1)<sub>Y</sub> gauge group. The second one is the gauge anomaly of the mixed type among SU(2)<sub>L</sub> and U(1)<sub>Y</sub> gauge groups. In the continuum theory, these gauge anomalies are generated from the following diagrams:
Then the conditions for the gauge anomaly cancellation in the electroweak theory are given in terms of the hypercharges as
$$\underset{L}{}Y^3\underset{R}{}Y^3=0,$$
(1.3)
and
$$\underset{\text{doublet(L)}}{}Y=0.$$
(1.4)
We can see that the assignment of the hyper-charges shown above indeed satisfies these conditions and one more condition as
$$\underset{\text{singlet(R)}}{}Y=0.$$
(1.5)
In order to show the exact cancellations of the (local) gauge anomalies at finite lattice spacing, we consider the cohomological classification of the 4+2-dimensional topological field for SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> electroweak theory. Our approach is then to determine the dependence on the admissible abelian gauge field of U(1)<sub>Y</sub> through topological argument, with SU(2)<sub>L</sub> gauge field fixed as background (cf. ). Although it does not determine the explicit dependence on SU(2)<sub>L</sub> gauge field, it turns out to be sufficient to show the exact cancellations of the gauge anomalies: we can show the cancellation of the gauge anomaly of the mixed type SU(2)<sub>L</sub><sup>2</sup> $`\times `$ U(1)<sub>Y</sub> at finite lattice spacing, as well as U(1)<sub>Y</sub><sup>3</sup>, using the pseudo reality of SU(2)<sub>L</sub> and the anomaly cancellation conditions in the electroweak theory given in terms of the hyper-charges of U(1)<sub>Y</sub>.
This paper is organized as follows. In section 2, we introduce the 4+2 dimensional topological field for the electroweak theory and discuss its specific features for SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> gauge groups. In section 3, we formulate the Poincaré lemma in 4+2 dimensions and determine the dependence on the admissible U(1)<sub>Y</sub> field at finite lattice spacing. In section 4, we show the exact cancellations of gauge anomalies of the mixed type SU(2)<sub>L</sub><sup>2</sup> $`\times `$ U(1)<sub>Y</sub> as well as U(1)<sub>Y</sub><sup>3</sup>, using the pseudo reality of SU(2)<sub>L</sub> and the anomaly cancellation conditions in the electroweak theory given in terms of the hyper-charges of U(1)<sub>Y</sub>. In section 5, we give some discussions.
## 2 4+2 dimensional topological field for Electroweak theory on the lattice
### 2.1 Weyl fermions on the lattice
Let us consider lattice Dirac fermion which is described by a gauge-covariant and local lattice Dirac operator which satisfies the Ginsparg-Wilson relation.
$$D\gamma _5+\gamma _5D=aD\gamma _5D.$$
(2.1)
The action of the Dirac fermion is written as
$$S=a^4\underset{x}{}\overline{\psi }(x)D\psi (x).$$
(2.2)
In the case of Neuberger’s Dirac operator
$$D=\frac{1}{a}\left(1+\gamma _5\frac{H}{\sqrt{H^2}}\right),$$
(2.3)
where $`H`$ is defined by the hermitian Wilson-Dirac operator
$$H=\gamma _5\left(\underset{\mu }{}\left\{\frac{1}{2}\gamma _\mu \left(_\mu _\mu ^{}\right)+\frac{a}{2}_\mu _\mu ^{}\right\}\frac{m_0}{a}\right),$$
(2.4)
locality of the action has been proved rigorously for gauge fields with bounded field strength .
$$1U_{\mu \nu }(x)ϵ,16(2+\sqrt{2})ϵ>|1m_0|^2.$$
(2.5)
This proof has been extended to the case where $`H`$ is defined by the transfer matrix of the five-dimensional Wilson fermion .
The action is invariant under the transformation which can be regarded as the chiral transformation on the lattice:
$$\delta \psi (x)=\gamma _5\left(1aD\right)\psi (x),\delta \overline{\psi }(x)=\overline{\psi }(x)\gamma _5.$$
(2.6)
By virtue of this exact chiral symmetry, we can define left-handed Weyl fermion on the lattice by the following projections:
$$\widehat{P}_{}\psi _L(x)=\psi _L(x),\overline{\psi }_L(x)P_+=\overline{\psi }_L(x).$$
(2.7)
$`\widehat{P}_{}`$ is the chiral projector defined as
$$\widehat{P}_{}=\frac{1\widehat{\gamma }_5}{2},\widehat{\gamma }_5=\gamma _5(1aD).$$
(2.8)
$`P_+`$ is the usual chiral projector defined with $`\gamma _5`$. The right-handed Weyl fermions can be defined in the similar manner.
The functional measure for the Weyl fermion can be defined as follows: we first introduce chiral bases $`\left\{v_j(x)\right\}`$ and $`\left\{\overline{v}_k(x)\right\}`$ as
$$\widehat{P}_{}v_j(x)=v_j(x),\overline{v}_k(x)P_+=\overline{v}_k(x),$$
(2.9)
and expand the Weyl fermion fields in terms of the chiral bases with the coefficients which generate the Grassmann algebra,
$$\psi (x)=\underset{j}{}v_j(x)c_j,\overline{\psi }(x)=\underset{k}{}\overline{c}_k\overline{v}_k(x).$$
(2.10)
Then the functional measure of the Weyl fermion can be defined as
$$\underset{x}{}d\psi _L(x)d\overline{\psi }_L(x)=\underset{j}{}dc_j\underset{k}{}d\overline{c}_k.$$
(2.11)
Given the definition for the functional measure of the Weyl fermion, the chiral determinant is evaluated as
$`Z_W`$ $`=`$ $`{\displaystyle \underset{x}{}d\psi _L(x)d\overline{\psi }_L(x)\mathrm{exp}\left(a^4\underset{x}{}\overline{\psi }_L(x)D\psi _L(x)\right)}`$ (2.12)
$`=`$ $`detM_{kj},`$ (2.13)
where
$$M_{kj}=a^4\underset{x}{}\overline{v}_k(x)Dv_j(x)=\left(\overline{v}_kDv_j\right).$$
(2.14)
### 2.2 4+2 dimensional topological field
The question is how to construct the functional measure for the Weyl fermions so that gauge invariance is maintained at finite lattice spacing. As shown by Lüscher , this question can be formulated as the cohomological problem of a certain topological field which is defined on the four-dimensional lattice plus two continuum dimensions.
The 4+2 dimensional topological field is introduced as follows. We consider lattice gauge fields
$$U_\mu (z)G,z=(x_\mu ,t,s),\mu =1,2,3,4$$
(2.15)
which depend on two additional real coordinates $`t`$ and $`s`$. We also introduce gauge potentials $`A_t(z)`$ and $`A_s(z)`$ along these directions and define the associated field tensor by
$$F_{ts}(z)=_tA_s(z)_sA_t(z)+i[A_t(z),A_s(z)].$$
(2.16)
The covariant derivative in these directions is defined as
$$D_r^AU_\mu (z)=_rU_\mu (z)+i\left[A_r(z)U_\mu (z)U_\mu (z)A_r(z+a\widehat{\mu })\right],r=t,s$$
(2.17)
which transforms in the same way as $`U_\mu (z)`$ under gauge transformations in 4+2 dimensions. Then we consider the following 4+2 dimensional field which is gauge invariant and local:
$`q(z)=i\text{tr}\{[{\displaystyle \frac{1}{4}}\widehat{\gamma }_5[D_t^A\widehat{P}_{},D_s^A\widehat{P}_{}]+{\displaystyle \frac{1}{4}}[D_t^A\widehat{P}_{},D_s^A\widehat{P}_{}]\widehat{\gamma }_5`$
$`+{\displaystyle \frac{i}{2}}R(F_{ts})\widehat{\gamma }_5](x,x)\}.`$
The trace is taken over the Dirac and flavor indices only.
By noting
$$a^4\underset{x}{}q(x)=i\mathrm{Tr}\left\{\widehat{P}_{}[_t\widehat{P}_{},_s\widehat{P}_{}]\frac{i}{2}_t\left[R(A_s)\widehat{\gamma }_5\right]+\frac{i}{2}_s\left[R(A_t)\widehat{\gamma }_5\right]\right\}$$
(2.19)
and making use of the identity
$$\mathrm{Tr}\left\{\delta _1\widehat{P}_{}\delta _2\widehat{P}_{}\delta _3\widehat{P}_{}\right\}=0,$$
(2.20)
we can show that this 4+2 dimensional field satisfies
$$a^4\underset{x}{}𝑑t𝑑s\delta q(z)=0$$
(2.21)
for all local variations of the link variables $`U_\mu (z)`$ and the potential $`A_r(z)`$, i.e. it is a topological field.
It has been shown by Lüscher that if this topological field is in the trivial cohomology class, i.e. it is equal to the divergence of a gauge-invariant local current,
$$q(z)=_\mu ^{}k_\mu (z)+_tk_s(z)_sk_t(z),$$
(2.22)
then using $`k_r(z)`$ it is possible to construct a gauge-covariant local current $`j_\mu ^a(x)`$ which satisfies the integrability condition in differential form and the anomalous conservation law.
### 2.3 4+2 dimensional topological field for SU(2)<sub>L</sub>$`\times `$U(1)<sub>Y</sub> electroweak theory
In order to construct the 4+2 dimensional topological field for SU(2)<sub>L</sub>$`\times `$U(1)<sub>Y</sub> electroweak theory, we consider lattice SU(2)<sub>L</sub> and U(1)<sub>Y</sub> gauge fields,
$$U_\mu ^{(1)}(z)\mathrm{U}(1),U_\mu ^{(2)}(z)\mathrm{SU}(2),$$
(2.23)
which satisfy the admissibility conditions with sufficiently small constants $`ϵ^{(2)}`$ and $`ϵ^{(1)}`$:
$$1U_{\mu \nu }^{(2)}(x)<ϵ^{(2)},1U_{\mu \nu }^{(1)}(x)<ϵ^{(1)}.$$
(2.24)
When $`ϵ^{(1)}<1/6Y\times \pi /3`$, the admissible abelian lattice gauge fields can be expressed in terms of vector potentials
$$U_\mu ^{(1)}(x)=\mathrm{exp}\left(iA_\mu (x)\right)$$
(2.25)
which has the following properties:
$$F_{\mu \nu }(x)\frac{1}{i}\mathrm{ln}U_{\mu \nu }(x)=_\mu A_\nu (x)_\nu A_\mu (x),|A_\mu (x)|\pi (1+8x).$$
(2.26)
This representation of the link variable is unique up to the gauge transformation with the parameter $`\omega (x)`$ which takes values in integer multiple of $`2\pi `$.
$$\stackrel{~}{A}_\mu (x)=A_\mu (x)+_\mu \omega (x).$$
(2.27)
We also introduce gauge potentials along the two additional dimensions $`A_r(z),r=t,s`$ for U(1)<sub>Y</sub> and $`B_r(z),r=t,s`$ for SU(2)<sub>L</sub> and denote the 4+2 dimensional gauge fields as follows:
$$A_\mu (z)=(A_k(z),A_t(z),A_s(z)),U_\mu (z)=(U_k^{(2)}(z),iB_t(z),iB_s(z)),$$
(2.28)
where $`\mu =1,\mathrm{},6`$ and we use the Latin index $`i=1,2,3,4`$ for four-dimensional lattice here after. Then the 4+2 dimensional topological field $`q(z)`$ for SU(2)<sub>L</sub>$`\times `$U(1)<sub>Y</sub> electroweak theory can be regarded as a gauge-invariant local functional of the 4+2 dimensional gauge field variables $`A_\mu (z)`$ and $`U_\mu (z)`$:
$$q(z)=q(z;A_\mu (z),U_\mu (z))$$
(2.29)
It follows from the charge conjugation property of the lattice Dirac operator that $`q(z)`$ changes sign under complex-conjugation of the representations of the gauge fields
$$q(z;A_\mu (z),U_\mu (z))=q(z;A_\mu (z),U_\mu ^{}(z))$$
(2.30)
Since SU(2)<sub>L</sub> is pseudo real, there exists a unitary transformation $`S`$ such that
$$SU_\mu ^{}S^1=U_\mu .$$
(2.31)
Then we obtain
$$q(z;A_\mu (z),U_\mu (z))=q(z;A_\mu (z),U_\mu (z)).$$
(2.32)
We note also that $`q(z)`$ vanishes identically when the U(1)<sub>Y</sub> gauge fields are switched off:
$$q(z;0,U_\mu (z))=0.$$
(2.33)
## 3 Cohomological classification of the topological field in 4+2 dimensions
### 3.1 Analysis of the 4+2 dimensional topological field
In this section, we will formulate the Poincaré lemma in 4+2 dimensions and examine the dependence of $`q(x,t,s)`$ on the admissible U(1)<sub>Y</sub> field through topological argument. In the course of the analysis, SU(2)<sub>L</sub> gauge field is fixed as background.
Since $`q(x,t,s)`$ smoothly depends on 4+2 dimensional U(1) gauge potential $`A_\mu (x,t,s)`$ and its differentials, the variation of the topological field can be expressed as
$$\delta q(x,t,s)=\underset{m,n=0,1}{}\underset{y}{}\frac{q(x,t,s)}{[_s^m_t^nA_\mu (y,t,s)]}_s^m_t^n\delta A_\mu (y,t,s).$$
(3.1)
By definition, $`q(x,t,s)`$ contains at most the first-order differentials of the vector potential $`A_\mu (x,t,s)`$ in the continuous coordinates. Therefore, we may restrict the sum over $`m,n`$ to $`0,1`$.
If we define the differential operator in 4+2 dimensions in the above expression as
$$L_\mu (x,y,t,s)=\underset{m,n=0,1}{}\frac{q(x,t,s)}{[_s^m_t^nA_\mu (y,t,s)]}_s^m_t^n,$$
(3.2)
then the topological property and the gauge invariance of the 4+2 dimensional field lead to the following conditions for the differential operator $`L_\mu `$.
$$𝑑t𝑑sa^4\underset{x}{}\underset{y}{}L_\mu (x,y,t,s)\delta A_\mu (y,t,s)=0,$$
(3.3)
and
$$L_\mu (x,y,t,s)_\mu =0.$$
(3.4)
The topological field $`q(x,t,s)`$ itself can be expressed with this operator as
$$q(x,t,s)=\alpha (x,t,s)+\underset{y}{}_0^1𝑑uL_\mu (x,y,t,s)|_{AuA}A_\mu (y,t,s),$$
(3.5)
where $`\alpha (x,t,s)`$ is the part which does not depend on the abelian gauge field $`A_\mu (x,t,s)`$. Then the problem reduces to examine the cohomological properties of the differential operator $`L_\mu `$. In order to examine such an operator and determine the form of the topological field $`q(x,t,s)`$, we will next formulate the Poincaré lemma which is applicable to the differential operators in 4+2 dimension. This is the extension of the Poincaré lemma on the lattice given in along the line of the analysis in the continuum theory of .
### 3.2 Poincaré lemma in 4+2 dimensions
We first introduce a Grassmann algebra with basis element
$$dx_1,dx_2,dx_3,dx_4,dx_5=dt,dx_6=ds$$
(3.6)
and denote them by $`dx_\mu (\mu =1,\mathrm{}6)`$ collectively. A $`k`$-form ($`0k6`$) in 4+2 dimensions is then defined as:
$$f(z)=\frac{1}{k!}f_{\mu _1\mathrm{}\mu _k}(z)dx_{\mu _1}\mathrm{}dx_{\mu _k}\mathrm{\Omega }_k,z=(x_i,t,s).$$
(3.7)
We assume that $`f_{\mu _1\mathrm{}\mu _k}(z)`$ is a smooth function in the continuous coordinates $`t,s`$ and is locally supported in the lattice coordinate $`x_i`$.
The exterior differential operator, which is a map $`\mathrm{\Omega }_k\mathrm{\Omega }_{k+1}`$, is defined by $`d=_{i=1}^4dx_i_i+_{r=t,s}dr_r`$ where $`_i`$ is forward difference operator on the four-dimensional lattice. It acts on the $`k`$-forms according to
$$df(z)=\frac{1}{k!}_\mu f_{\mu _1\mathrm{}\mu _k}(z)dx_\mu dx_{\mu _1}\mathrm{}dx_{\mu _k}.$$
(3.8)
The associated divergence operator $`d^{}:\mathrm{\Omega }_k\mathrm{\Omega }_{k1}`$ is defined as
$$d^{}f(z)=\frac{1}{(k1)!}_\mu ^{}f_{\mu \mu _2\mathrm{}\mu _k}(z)dx_{\mu _2}\mathrm{}dx_{\mu _k},$$
(3.9)
where $`_i^{}`$ is backward difference operator.
Now we consider a class of differential operators $`L`$ which is a map $`L:\mathrm{\Omega }_l\mathrm{\Omega }_k`$ such that
$$Lf(x,t,s)=\frac{1}{k!l!}dx_{\mu _1}\mathrm{}dx_{\mu _k}\underset{y,n,m}{}L_{\mu _1\mathrm{}\mu _k;\nu _1\mathrm{}\nu _l}^{n,m}(x;y,t,s)_t^n_s^mf_{\nu _1\mathrm{}\nu _l}(y,t,s),$$
(3.10)
where $`L_{\mu _1\mathrm{}\mu _k;\nu _1\mathrm{}\nu _l}^{n,m}(x;y,t,s)`$ is a local function in $`x(y),t,s`$ which is exponentially decaying in $`x(y)`$ with respect to the reference point $`y(x)`$. We assume $`m,n=0,1`$ in the following discussions. When $`L`$ is not a differential operator, i.e. $`n=m=0`$ and is proportional to the Kronecker delta $`\delta _{x,y}`$, we refer such an operator as zero degree.
$$L^0f(x,t,s)=\frac{1}{k!l!}dx_{\mu _1}\mathrm{}dx_{\mu _k}L_{\mu _1\mathrm{}\mu _k;\nu _1\mathrm{}\nu _l}^0(x,t,s)f_{\nu _1\mathrm{}\nu _l}(x,t,s).$$
(3.11)
We can show the following Poincaré lemma for these differential operators in 4+2 dimensions.
Lemma (Poincaré lemma) If $`L:\mathrm{\Omega }_l\mathrm{\Omega }_k`$ is any differential operator satisfying
$$dL(x,y,t,s)=0,$$
(3.12)
then there exists a differential operator $`M(x,y,t,s)`$ and an operator of zero degree $`L^0`$ such that
$$L(x,y,t,s)=\delta _{k,6}\delta _{x,y}L^0(x,t,s)+dM(x,y,t,s),$$
(3.13)
Note that the product of $`d`$ and $`M`$ is a product of operators.
The proof of the lemma is given as follows. We first note the fact that for any differential operator $`L`$ in concern and any fixed continuous dimension $`r`$, there is a unique decomposition of the operator into two operators so that
$$L=S+_rR,$$
(3.14)
where $`S`$ is zero degree with respect to $`_r`$. This is because $`L`$ is at most a first-order differential operator in terms of $`_r`$ assuming the form
$$L=L_0+L_1_r$$
(3.15)
with $`L_0`$ and $`L_1`$ zero degree with respect to $`_r`$ and it can be rewritten uniquely into
$$L=\left(L_0[_rL_1]\right)+_rL_1,$$
(3.16)
where the bracket $`[_rL_1]`$ stands for the fact that $`_r`$ in it acts only on $`L_1`$. When $`L`$ is gauge invariant, by acting it on a constant, we infer that $`L_0`$ is also gauge invariant. Then it also follows that $`L_1`$ is gauge invariant. Therefore, in the decomposition of Eq. (3.14), both $`S`$ and $`R`$ are gauge invariant.
Then we can decompose any differential operator $`L:\mathrm{\Omega }_l\mathrm{\Omega }_k`$ into the sequence:
$$L=(L_5+_sR_5)ds+Z_6,$$
(3.17)
and
$$L_5=(L_4+_tR_4)dt+Z_5,$$
(3.18)
where $`L_5`$ is degree zero with respect to $`_s`$ and $`L_4`$ is degree zero with respect to both $`_s`$ and $`_t`$.
From the condition Eq. (3.12),
$$dL=\stackrel{~}{d}\left(L_5+_sR_5\right)ds+\left(\stackrel{~}{d}+ds_s\right)Z_6=0,$$
(3.19)
where $`\stackrel{~}{d}=_{i=1}^4dx_i_i+dt_t`$. Then the coefficient of $`ds`$ must vanish
$$\stackrel{~}{d}L_5+_s\left(\stackrel{~}{d}R_5+()^kZ_6\right)=0.$$
(3.20)
Since this can be regarded as the unique decomposition of zero operator with respect to $`_s`$, we have
$`\stackrel{~}{d}L_5=0,`$ (3.21)
$`\stackrel{~}{d}R_5+()^kZ_6=0.`$ (3.22)
In a similar manner, from the condition Eq. (3.21), we obtain
$`\overline{d}L_4=0,`$ (3.23)
$`\overline{d}R_4+()^{k1}Z_5=0,`$ (3.24)
where $`\overline{d}=_{i=1}^4dx_i_i`$.
Combing these results, we have
$$L=L_4dtds+dK,$$
(3.25)
where
$$K=()^{k1}R_5+()^{k2}R_4ds.$$
(3.26)
$`L_4`$ is zero degree with respect to both $`_s`$ and $`_t`$ and satisfies $`\overline{d}L_4=0`$.
Now we recall the Poincaré lemma on four-dimensional lattice given in . From the condition $`\overline{d}L_4=0`$, we have
$$L_4(x;y,t,s)=\delta _{k2,4}\delta _{x,y}L_4^0(x,t,s)+\overline{d}T(x;y,t,s)$$
(3.27)
where
$$L_4^0(x,t,s)=\underset{y}{}L_4(y;x,t,s).$$
(3.28)
Then we can write $`L`$
$$L(x,y,t,s)=\delta _{k,6}\delta _{x,y}L^0(x,t,s)+dM(x,y,t,s),$$
(3.29)
where
$$M=Tdtds+K,L^0=L_4^0dtds.$$
(3.30)
It is clear from the above construction that $`L^0`$ and $`M`$ possesses the same locality and gauge-transformation properties as $`L`$. When $`L`$ is gauge invariant under SU(2)<sub>L</sub> and U(1)<sub>Y</sub> gauge transformations, $`L_0`$ and $`M`$ are also gauge invariant.
For the operators $`L`$ which satisfy $`d^{}L=0`$, we can show the following form of the Poincaré lemma: If $`L:\mathrm{\Omega }_l\mathrm{\Omega }_k`$ is any differential operator satisfying
$$d^{}L(x,y,t,s)=0,$$
(3.31)
then there exists a differential operator $`M(x,y,t)`$ and an operator of zero degree $`L^0`$ such that
$$L(x,y,t,s)=\delta _{k,0}\delta _{x,y}L^0(x,t,s)+d^{}M(x,y,t,s).$$
(3.32)
For the operators $`L`$ which satisfy $`Ld=0`$ and $`Ld^{}=0`$, the lemma can be derived in similar manners.
### 3.3 Structure of the 4+2 dimensional topological field
By using the Poincaré lemma in 4+2 dimensions, we next determine the dependence of $`q(z)`$ on the admissible abelian gauge field of U(1)<sub>L</sub>. We will show the following lemma:
Lemma The 4+2 dimensional topological field $`q(z)`$ for the SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> electroweak theory is written in the following form.
$`q(z)`$ $`=`$ $`\alpha (z)+\beta _{\mu \nu }(z\widehat{\mu }\widehat{\nu })F_{\mu \nu }(z\widehat{\mu }\widehat{\nu })`$ (3.33)
$`+`$ $`\gamma _{\mu \nu \rho \sigma }(z\widehat{\mu }\widehat{\nu }\widehat{\rho }\widehat{\sigma })F_{\rho \sigma }(z\widehat{\mu }\widehat{\nu }\widehat{\rho }\widehat{\sigma })F_{\mu \nu }(z\widehat{\mu }\widehat{\nu })`$
$`+`$ $`\delta ϵ_{\mu \nu \rho \sigma \lambda \tau }F_{\lambda \tau }(z\widehat{\mu }\widehat{\nu }\widehat{\rho }\widehat{\sigma }\widehat{\lambda }\widehat{\tau })\times `$
$`F_{\rho \sigma }(z\widehat{\mu }\widehat{\nu }\widehat{\rho }\widehat{\sigma })F_{\mu \nu }(z\widehat{\mu }\widehat{\nu })`$
$`+`$ $`_\mu ^{}k_\mu (z)`$
where $`F_{\mu \nu }=_\mu A_\nu _\nu A_\mu `$ is the 4+2 dimensional field strength of U(1)<sub>Y</sub> gauge potential. $`k_\mu (z)`$ is a local current which is gauge invariant under SU(2)<sub>L</sub> and U(1)<sub>Y</sub> gauge transformations. $`\delta `$ is a constant. $`\alpha (z)`$, $`\beta _{\mu \nu }(z)`$ and $`\gamma _{\mu \nu \rho \sigma }(z)`$ are certain gauge invariant local functions which may depend on the SU(2)<sub>L</sub> gauge field and satisfy
$$_\mu ^{}\beta _{\mu \nu }(z)=0,_\mu ^{}\gamma _{\mu \nu \rho \sigma }(z)=0.$$
(3.34)
Note also that $`\widehat{5}=\widehat{6}=0`$ and $`z=(x_i,t,s)`$.
The proof consists of three steps.
#### 3.3.1 Step one
We first consider the differential operator
$$L_\mu (x,y,t,s)=\underset{m,n=0,1}{}\frac{q(x,s,t)}{[_s^m_t^nA_\mu (y,s,t)]}_s^m_t^n.$$
(3.35)
We may regard this operator as a map $`L_\mu :\mathrm{\Omega }_1\mathrm{\Omega }_0`$. Then using the Poincaré lemma, we can write as
$$L_\mu (x,y,t,s)=\delta _{x,y}L_\mu ^0(x,t,s)+_\nu ^{}K_{\nu ;\mu },$$
(3.36)
where $`L^0`$ is zero degree and $`K_{\nu ;\lambda }`$ is a map $`K_{\nu ;\lambda }:\mathrm{\Omega }_1\mathrm{\Omega }_1`$. The topological property of the 4+2 dimensional field implies
$$0=𝑑t𝑑sa^4\underset{x}{}\delta q(x,t,s)=𝑑t𝑑sa^4\underset{x}{}L_\mu ^0(x,t,s)\delta A_\mu (x,t,s).$$
(3.37)
Since $`L_\mu ^0(x,t,s)`$ is zero degree, it must vanish identically,
$$L_\mu ^0(x,t,s)=0.$$
(3.38)
On the other hand, the gauge invariance of the 4+2 dimensional field implies
$$L_\mu (x,y,t,s)_\mu =0$$
(3.39)
or
$$_\nu ^{}K_{\nu ;\mu }(x,y,t,s)_\mu =0.$$
(3.40)
Then, using the Poincaré lemma, we obtain
$$K_{\nu ;\mu }(x,y,t,s)_\mu =_\lambda ^{}H_{\lambda \nu }(x,y,t,s),$$
(3.41)
where $`H_{\lambda \nu }:\mathrm{\Omega }_0\mathrm{\Omega }_2`$. Using the Poincaré lemma again, it can be cast into the form
$$H_{\lambda \nu }(x,y,t,s)=\delta _{x,y}H_{\lambda \nu }^0(x,t,s)+R_{\lambda \nu ;\rho }(x,y,t,s)_\rho ,$$
(3.42)
where $`R_{\lambda \nu ;\rho }:\mathrm{\Omega }_1\mathrm{\Omega }_2`$. We can eliminate $`R_{\lambda \nu ;\rho }`$ term in the above expression by redefining $`K_{\nu ;\mu }`$ so that
$$K_{\nu ;\mu }K_{\nu ;\mu }+_\lambda ^{}R_{\lambda \nu ;\mu }(x,y,t,s),$$
(3.43)
which does not affect the relation $`L_\mu =_\nu ^{}K_{\nu ;\mu }`$. Thus we obtain
$$H_{\lambda \nu }(x,y,t,s)=\delta _{x,y}H_{\lambda \nu }^0(x,t,s).$$
(3.44)
We substitute this result into Eq. (3.41) and make it act on a constant, we obtain
$$[_\lambda ^{}H_{\lambda \nu }^0(x,t,s)]=0,$$
(3.45)
where the square bracket means that $`_\lambda ^{}`$ acts on the function $`H_{\lambda \nu }^0`$ rather than a product of the differential operators. Then it follows that<sup>2</sup><sup>2</sup>2 $`_\lambda ^{}\left[H_{\lambda \nu }^0(x,t,s)f(x,t,s)\right]`$ $`=\left[_\lambda ^{}H_{\lambda \nu }^0(x,t,s)\right]f(x,t,s)+H_{\lambda \nu }^0(x\widehat{\lambda },t,s)\left[_\lambda ^{}f(x,t,s)\right]`$ $`=\left[_\lambda ^{}H_{\lambda \nu }^0(x,t,s)\right]f(x,t,s)+H_{\lambda \nu }^0(x\widehat{\lambda },t,s)\left[_\lambda f(x\widehat{\lambda },t,s)\right].`$
$$_\lambda ^{}\delta _{x,y}H_{\lambda \nu }^0(x,t,s)=\delta _{x\widehat{\lambda },y}H_{\lambda \nu }^0(x\widehat{\lambda },t,s)_\lambda .$$
(3.46)
Then using the Poincaré lemma we obtain
$$K_{\nu ;\mu }(x,y,t,s)=\delta _{x\widehat{\mu },y}H_{\mu \nu }^0(x\widehat{\mu },t,s)+\omega _{\nu ;\mu \rho }(x,y,t,s)_\rho .$$
(3.47)
We now substitute this result into Eq. (3.36) and obtain
$$L_\mu (x,y,t,s)=\delta _{x\widehat{\nu }\widehat{\mu },y}H_{\mu \nu }^0(x\widehat{\nu }\widehat{\mu },t,s)_\nu +_\nu ^{}\omega _{\nu ;\mu \rho }(x,y,t,s)_\rho .$$
(3.48)
Then the 4+2 dimensional topological field is given as follows.
$`q(x,t,s)`$ $`=`$ $`\alpha (x,t,s)+{\displaystyle _0^1}𝑑u{\displaystyle \underset{y}{}}L_\mu (x,y,t,s)|_{AuA}A_\mu (y,t,s)`$ (3.49)
$`=`$ $`\alpha (x,t,s)+{\displaystyle _0^1}𝑑u\left[H_{\mu \nu }^0(x\widehat{\nu }\widehat{\mu },t,s)\right]|_{AuA}_\nu A_\mu (x\widehat{\nu }\widehat{\mu },t,s)`$
$`+_\nu ^{}\left({\displaystyle \underset{y}{}}{\displaystyle _0^1}𝑑u\omega _{\nu ,\mu \rho }(x,y,t,s)|_{AuA}_\rho A_\mu (y,t,s)\right).`$
Setting
$`\varphi _{\mu \nu }(x,t,s)`$ $`=`$ $`{\displaystyle \frac{1}{2}}{\displaystyle _0^1}𝑑u\left[H_{\mu \nu }^0(x,t,s)\right]|_{AuA},`$ (3.50)
$`\theta _\nu (x,t,s)`$ $`=`$ $`{\displaystyle \underset{y}{}}{\displaystyle _0^1}𝑑u\omega _{\nu ,\mu \rho }(x,y,t,s)|_{AuA}_\rho A_\mu (y,t,s),`$ (3.51)
we obtain
$$q(z)=\alpha (z)+\varphi _{\mu \nu }(z\widehat{\nu }\widehat{\mu })F_{\mu \nu }(z\widehat{\nu }\widehat{\mu })+_\nu ^{}\theta _\nu (z).$$
(3.52)
From the topological property of the 4+2 topological field,
$`0`$ $`=`$ $`{\displaystyle 𝑑t𝑑sa^4\underset{x}{}\delta q(x,t,s)}`$ (3.53)
$`=`$ $`{\displaystyle 𝑑s𝑑ta^4\underset{x}{}\frac{1}{2}\varphi _{\mu \nu }(z\widehat{\nu }\widehat{\mu })_\nu \delta A_\mu (z\widehat{\nu }\widehat{\mu })}.`$
By the integration by parts, we obtain
$$_\nu ^{}\varphi _{\mu \nu }(z)=0.$$
(3.54)
#### 3.3.2 Step two
Next we examine the differential operator which is obtained from the variation of $`\varphi _{\mu \nu }(x,t,s)`$.
$$\delta \varphi _{\mu \nu }(z)=\underset{y}{}\underset{m,n=0,1}{}\frac{\varphi _{\mu \nu }(z)}{[_s^m_t^nA_\lambda (y,t,s)]}_s^m_t^n\delta A_\lambda (y,t,s).$$
(3.55)
We denote the differential operator in the above expression as $`L_{\mu \nu \lambda }(x,y,t,s)`$.
$$L_{\mu \nu \lambda }(x,y,t,s)=\underset{m,n=0,1}{}\frac{\varphi _{\mu \nu }(z)}{[_s^m_t^nA_\lambda (y,t,s)]}_s^m_t^n.$$
(3.56)
$`\varphi _{\mu \nu }(x,t,s)`$ itself can be expressed with $`L_{\mu \nu \lambda }(x,y,t,s)`$ as
$$\varphi _{\mu \nu }(z)=\beta _{\mu \nu }(z)+_0^1𝑑u\underset{y}{}L_{\mu \nu ,\lambda }(x,y,t,s)|_{AuA}A_\lambda (y,t,s).$$
(3.57)
It follows from the property Eq. (3.54) that
$$_\mu ^{}\delta \varphi _{\mu \nu }(z)=\delta \left[_\mu ^{}\varphi _{\mu \nu }(z)\right]=0,$$
(3.58)
which implies
$$_\mu ^{}L_{\mu \nu ,\lambda }(x,y,t,s)=0.$$
(3.59)
and in turn implies
$$_\mu ^{}\beta _{\mu \nu }(z)=0.$$
(3.60)
From the gauge invariance of $`\varphi _{\mu \nu }(z)`$, we have
$$L_{\mu \nu ,\lambda }(x,y,t,s)_\lambda =0.$$
(3.61)
Then following the similar argument as the first step, we obtain
$`\varphi _{\mu \nu }(z)`$ $`=`$ $`\beta _{\mu \nu }(z)`$ (3.62)
$`+\eta _{\mu \nu \lambda \rho }(z\widehat{\lambda }\widehat{\rho })F_{\lambda \rho }(z\widehat{\lambda }\widehat{\rho })`$
$`+_\rho ^{}\theta _{\mu \nu \rho }(z),`$
and
$$_\rho ^{}\eta _{\mu \nu \lambda \rho }(z)=0,_\mu ^{}\beta _{\mu \nu }(z)=0.$$
(3.63)
#### 3.3.3 Step three
Finally, we examine the differential operator which is obtained from the variation of $`\eta _{\mu \nu \lambda \rho }`$. In the course, we encounter the operator of zero degree $`H_{\mu \nu \lambda \rho \sigma \tau }^0:\mathrm{\Omega }_0\mathrm{\Omega }_6`$, which satisfies the condition
$$_\tau ^{}H_{\mu \nu \lambda \rho \sigma \tau }^0(x,t,s)=0.$$
(3.64)
$`H^0`$ is the six form and it may be written with the totally antisymmetric tensor $`ϵ_{\mu \nu \lambda \rho \sigma \tau }`$ as
$$H_{\mu \nu \lambda \rho \sigma \tau }^0(x,t,s)=\delta (x,t,s)ϵ_{\mu \nu \lambda \rho \sigma \tau }.$$
(3.65)
Then the above condition implies that $`\delta `$ does not depend on $`(x,t,s)`$ and is a constant.
Following the similar argument as the first and second steps, we obtain
$`\eta _{\mu \nu \lambda \rho }(z)`$ $`=`$ $`\gamma _{\mu \nu \lambda \rho }(z)`$ (3.66)
$`+\delta ϵ_{\mu \nu \lambda \rho \sigma \tau }F_{\sigma \tau }(z\widehat{\sigma }\widehat{\tau })`$
$`+_\sigma ^{}\theta _{\mu \nu \lambda \rho \sigma }(z),`$
and
$$_\mu ^{}\gamma _{\mu \nu \lambda \rho }(z)=0,$$
(3.67)
where $`ϵ_{\mu \nu \lambda \rho \sigma \tau }`$ is the totally antisymmetric tensor and $`\delta `$ is a constant.
Combining these three results and using the Bianchi identity,
$$_{[\rho }^{}F_{\mu \nu ]}(x\widehat{\mu }\widehat{\nu })=_{[\rho }F_{\mu \nu ]}(x\widehat{\mu }\widehat{\nu }\widehat{\rho })=0,$$
(3.68)
we finally obtain Eq. (3.33) and complete the proof of the lemma.
In Eq. (3.33), there is a difference in the shifts of the lattice indices from the result obtained in . This difference, however, can be shown to be a total divergence.<sup>3</sup><sup>3</sup>3 In order to show this, we can use the following identity in Eqs. (3.52), (3.62), (3.66). $`f_{\mu \nu }(z)g_{\mu \nu }(z)f_{\mu \nu }(z\widehat{\mu }\widehat{\nu })g_{\mu \nu }(z\widehat{\mu }\widehat{\nu })`$ $`=`$ $`f_{\mu \nu }(z)g_{\mu \nu }(z)f_{\mu \nu }(z\widehat{\mu })g_{\mu \nu }(z\widehat{\mu })`$ $`+f_{\mu \nu }(z\widehat{\mu })g_{\mu \nu }(z\widehat{\mu })f_{\mu \nu }(z\widehat{\mu }\widehat{\nu })g_{\mu \nu }(z\widehat{\mu }\widehat{\nu })`$ $`=`$ $`_i^{}\left(f_{i\nu }(z)g_{i\nu }(z)+f_{\mu i}(z\widehat{\mu })g_{\mu i}(z\widehat{\mu })\right).`$
## 4 Exact cancellations of gauge anomalies in <br>SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> Electroweak theory
The result Eq. (3.33) concerning the dependence on the the admissible U(1)<sub>Y</sub> gauge field of the 4+2 dimensional topological field may be written symbolically in following form.
$$q(z;A_\mu ,U_\mu )=\alpha (z;U_\mu )+\beta (z;U_\mu )F+\gamma (z;U_\mu )F^2+\delta F^3+d^{}k(z;U_\mu ,A_\mu ).$$
(4.1)
As we discussed in section 2, the 4+2 dimensional topological field for SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> electroweak theory has the properties
$$q(z;A_\mu (z),U_\mu (z))=q(z;A_\mu (z),U_\mu (z))$$
(4.2)
and
$$q(z;0,U_\mu (z))=0.$$
(4.3)
From these two conditions, we infer
$$\alpha (z;U_\mu )=0.$$
(4.4)
and
$$\gamma (z;U_\mu )F^2=d^{}\frac{1}{2}\left(k(z;U_\mu ,A_\mu )+k(z;U_\mu ,A_\mu )\right).$$
(4.5)
Therefore, the 4+2 topological field turns out to have the following structure
$$q(z;A_\mu ,U_\mu )=\beta (z;U_\mu )F+\delta F^3+d^{}\frac{1}{2}\left(k(z;U_\mu ,A_\mu )k(z;U_\mu ,A_\mu )\right).$$
(4.6)
We note that $`\beta (z;U_\mu )`$ is a gauge-invariant local functional of SU(2)<sub>L</sub> gauge field, while $`\delta `$ is a constant
Now we recall the anomaly cancellation conditions for the electroweak theory which is given in terms of U(1)<sub>Y</sub> hyper-charges. Because of the cubic condition
$$\underset{L}{}Y^3\underset{R}{}Y^3=0,$$
(4.7)
the term $`\delta F^3`$ vanishes identically, if all the contributions from the fermions are summed up. On the other hand, because of the linear conditions
$$\underset{\text{doublet(L)}}{}Y=0.$$
(4.8)
and
$$\underset{\text{singlet(R)}}{}Y=0.$$
(4.9)
the term $`\beta (z;U_\mu )F`$, which represents the gauge anomaly of the mixed type, also vanishes identically, in each sectors of doublets and singlets. Thus we can see that the 4+2 dimensional topological field for the electroweak theory is indeed in the trivial cohomology class.
$$q(z;A_\mu ,U_\mu )=_\mu ^{}\frac{1}{2}\left(k_\mu (z;U_\mu ,A_\mu )k_\mu (z;U_\mu ,A_\mu )\right).$$
(4.10)
## 5 Summary and discussion
We have shown the exact cancellations of gauge anomalies of the SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> electroweak theory on the lattice, which is formulated based on the lattice Dirac operator satisfying the Ginsparg-Wilson relation. Our approach is to consider the cohomological classification of the 4+2-dimensional topological field proposed by Lüscher for SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> electroweak theory. Using the Poincaré lemma in 4+2 dimensions, we have determined the dependence on the admissible abelian gauge field of U(1)<sub>Y</sub> through topological argument, with SU(2)<sub>L</sub> gauge field fixed as background (cf. ). This turned out to be sufficient to show the exact cancellations of the gauge anomalies: using the pseudo reality of SU(2)<sub>L</sub> and the anomaly cancellation conditions in the electroweak theory given in terms of the hyper-charges of U(1)<sub>Y</sub>, we have shown the exact cancellation of the gauge anomaly of the mixed type SU(2)<sub>L</sub><sup>2</sup> $`\times `$ U(1)<sub>Y</sub> at finite lattice spacing, as well as U(1)<sub>Y</sub><sup>3</sup>.
As to the question of the cohomological classification of the 4+2 dimensional topological field for the electroweak theory, we may also invoke the elegant method based on the non-commutative differential calculus and BRST cohomology in order to explore the structure of the 4+2 dimensional topological field .
Towards the lattice construction of the SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> electroweak theory, the next step would be to show the integrability condition given in , which assures the existence of the functional measure of fermions with desired properties. For this purpose, we need to examine the possible global anomalies in SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> electroweak theory. Global SU(2) anomaly has been examined by Neuberger and by Bär and Campos in detail . It is also desirable to establish the existence of the model in a finite volume, as in the case of the abelian chiral gauge theories .
In order to extend our result to the whole SU(3)<sub>C</sub> $`\times `$ SU(2)<sub>L</sub> $`\times `$ U(1)<sub>Y</sub> standard model, we need to attack directly the non-abelian nature of the gauge anomalies of the mixed type of SU(3)<sub>C</sub> $`\times `$ SU(2)<sub>L</sub>, although there is no corresponding anomaly of this type in the continuum theory. This is also true, of course, when we consider more general non-ablelian chiral gauge theories. The recent work by Suzuki based on the BRST cohomology in four-dimensions could shed lights on this issue.
## Acknowledgments
The authors are grateful to A.I. Sanda, K. Yamawaki, S. Kitakado, S. Uehara, K. Morita and M. Harada for discussions and encouragements. Y.K. would like to thank S. Aoki, T. Onogi, Y. Taniguchi, T. Izubuchi, K. Nagai, J. Noaki, K. Nagao, N. Ukita, and H. So for enlightening discussions. The intensive discussions at the Summer Institute 99 at Yamanashi, Japan, were very suggestive and useful to complete this work. Y.K. would like to thank M. Lüscher, E. Seiler and P. Weisz for their kind hospitality at the Ringberg workshop. The authors would like to thank M. Lüscher and H. Suzuki for valuable comments. Y.K. is supported in part by Grant-in-Aid for Scientific Research of Ministry of Education (#10740116).
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# 1 Introduction
## 1 Introduction
Recent measurements of charm production in deep inelastic scattering (DIS) at HERA have shown that up to 25% of the total cross-section at small-x contains charm in the final state. This is within the expectations of perturbative QCD based on conventional parton distributions. We are now in a position to utilize this process to study details of the production mechanism of heavy quarks in general, and to extract useful information on the charm and gluon structure of the proton in particular.
Conventional perturbative QCD (PQCD) theory is formulated in terms of zero-mass quark-partons. For processes depending on one hard scale $`Q,`$ the well-known factorization theorem then provides a straightforward procedure for order-by-order perturbative calculations, as well as an associated intuitive parton picture interpretation of the perturbation series. Heavy quark production presents a challenge in PQCD because the heavy quark mass, $`mH`$ $`(H=c,b,t),`$ provides an additional hard scale which complicates the situation – it requires a different organization of the perturbative series depending on the relative magnitudes of $`mH`$ and $`Q`$.
The two standard methods for PQCD calculation of heavy quark processes represent two diametrically opposite ways of reducing the two-scale problem to a one-scale problem. (i) In the conventional *parton model approach* used in many global QCD analyses of parton distributions and Monte Carlo programs, the zero-mass parton approximation is applied to a heavy quark calculation as soon as the typical energy scale<sup>d</sup><sup>d</sup>dWe use $`Q`$ as the generic name for a typical kinematic physical scale. It could be $`Q`$, $`W`$, or $`pT`$, depending on the process. of the physical process $`Q`$ is above the mass threshold $`mH`$. This leaves $`Q`$ as the only apparent hard scale in the problem. (ii) In the *heavy quark approach* which played a dominant role in “NLO calculations” of the production of heavy quarks , the quark $`H`$ is always treated as a “heavy” particle and never as a parton. The mass parameter $`mH`$ is explicitly kept along with $`Q`$ as if they are of the same order, irrespective of their real relative magnitudes.
The co-existence of these two opposite approaches represents an uneasy dichotomy in the current literature. On physical grounds, the zero-mass parton picture of heavy quarks should be applicable at energy scales very much larger than the relevant quark mass, $`mH`$ $``$ $`Q`$, whereas the heavy quark approach (often referred to as the fixed-flavor-number (FFN) scheme) should be more appropriate at energy scales comparable to the quark mass $`mH`$ $``$ $`Q`$. The actual experimental regime often lies *in between* these two extreme regions, where the validity of either approach can be called into question. There is, however, a natural way to incorporate both approaches in a unified framework in PQCD which provides a smooth transition between the two. This has been formulated in a series of papers over the years, Refs. <sup>e</sup><sup>e</sup>eSee Ref. for a brief review., which has now been adopted, in different guises, by most recent literature on heavy quark production in PQCD.
To see the basic ideas behind this unified picture, let us focus explicitly on the production of charm ($`H=c`$) in deep inelastic scattering. All considerations apply to a generic heavy quark. Consider the PQCD calculation of the $`F2(x,Q)`$ structure function which receives substantial contribution from charm production as mentioned earlier. The underlying physical ideas are illustrated graphically in Fig. 1 where the charm contribution to this structure function, denoted $`F2^c(x,Q)`$, is plotted as a function of $`Q`$ at some fixed value of $`x`$.
Near threshold $`Qmc`$, it is natural to consider the charm quark as a heavy particle, and to adopt the 3 active-parton-flavor scheme of calculation (the “heavy quark approach”). As $`Q`$ becomes large compared to $`mc`$, the fixed-order calculation in this approach becomes unreliable since the perturbative expansion contains terms of the form $`\alpha s^n\mathrm{log}^n\left(mc^2/Q^2\right)`$ at any order $`n`$, which ruin the convergence of the series—these terms are *not infra-red safe* as $`mc`$ $`0`$ or $`Q\mathrm{}`$. Thus the uncertainty of the 3-flavor calculation grows as $`Q/mc`$ becomes large. This is illustrated in Fig. 1 as an error band marked by horizontal hashes which is narrow near threshold but becomes ever wider as $`Q/mc`$ increases. On the other hand, starting from the high energy end ($`Qmc`$), the most natural calculational scheme to adopt is the conventional 4-flavor scheme with active charm partons. (In this approach, the infra-red unsafe large logarithms mentioned earlier are “resummed” and absorbed into the finite charm parton distributions.) However, as we go down the energy scale toward the charm production threshold region, the 4-flavor calculation becomes unreliable because the approximation $`mc=0`$ deteriorates as $`Qmc`$. The uncertainty band of such a calculation is outlined in Fig. 1 by the vertical hashes – it is narrow at high energies, but becomes increasingly wider as one approaches the threshold region.
The intuitive ideas embodied in Fig 1 illustrate that: (i) these two conventional approaches are individually unsatisfactory over the full energy range, but are mutually complementary; and (ii) the most reliable PQCD prediction for the physical $`F2(x,Q)`$, at a given order of calculation, can be obtained by utilizing the most appropriate scheme at that energy scale $`Q`$, resulting in a composite scheme, as represented by the cross-hashed region in Fig. 1. The use of a composite scheme consisting of different numbers of flavors in different energy ranges, rather than a fixed number of flavors, is familiar in the conventional zero-mass parton picture. The new formalism espoused in Refs. provides a quantum field theoretical basis for this intuitive picture in the presence of non-zero quark mass.
This generalization brings about several important distinguishing features and insights. First, after the hard cross-section is rendered infra-red safe by factorizing out the $`\mathrm{ln}(mc)`$ terms, the remaining finite dependence on $`mc`$ can be kept in the hard cross-sections to maintain better accuracy in the intermediate energy region.<sup>f</sup><sup>f</sup>fBy contrast, in the standard literature the $`mc0`$ limit is tied to the proof of factorization in the first place. This association is not needed, cf. One can then show that the unified formalism reproduces the two conventional approaches in well-defined ways in their respective regions of applicability -. Secondly, it has become clear recently that there is much inherent flexibility in the choice of the transition energy scale (cf. Fig. 1) as well as the detailed matching condition between the 3- and 4-flavor calculations. This makes it possible to have a variety of different implementations of the new formalism , with different emphases. Properly understood, this feature can be exploited to lend strength to the formalism; but, during the development of the theory, the subtleties have given rise to much misunderstanding and confusion in existing literature.
In this paper, we study charm production in DIS using the general mass formalism, and compare the results with recent HERA data. The main goals of this work are:
* (i) Through this specific example, we give a concise and careful presentation of the general formalism, including the relatively simple operational procedure for calculating its various components at any order of $`\alpha _s`$. We hope this will help fill the gap between the original, relatively sketchy, ACOT paper and the recent, more technical, all-order proof of the formalism by Collins .
* (ii) We carry out the numerical calculations to make concrete the intuitive ideas illustrated in Fig. 1; and to demonstrate the validity of the physical principles underlying the composite scheme.
* (iii) We show that the flexibility of the general formalism mentioned above can result in efficient PQCD calculations of inclusive quantities at relatively low order in $`\alpha _s`$ compared to FFN calculations – because the relevant physics has been effectively “resummed” by the appropriate scheme adopted for the given energy scale.
Item (i) is basic; it forms the foundation for the other two parts. However, since this part is about the clarification of the existing theoretical formalism rather than new work, and since not all readers are concerned with theoretical precision, we have elected to place it in the Appendix. Hopefully, this will increase the accessibility of the main body of the paper. In regards to the underlying physics, we believe the discussion in the appendix should make a useful contribution to clear the confusion and misunderstanding among the various approaches that have been proposed in the recent literature following .
Based on the terminologies discussed in this Appendix, in Sec. 2 we describe the complete order $`\alpha _s`$ calculation carried out in this paper in relation to previous work on this subject. The new calculation extends the validity of the original ACOT results to NLO in the high energy regime – on the same level as the conventional zero-mass total inclusive structure functions. In addition, the new perspective, as illustrated in Fig. 1 above and discussed quantitatively in the paper proper, allows a re-assessment of the physical predictions of the order $`\alpha _s`$ calculation near the threshold region, making it a viable alternative to the order $`\alpha _s^2`$ FFN calculation. In Sec. 3 we present the numerical results on inclusive charm production, and demonstrate that the validity and the efficiency of the general formalism, as described in (ii) and (iii) above, are indeed seen at this order. We show that very good agreement with recent HERA data on inclusive $`F2^c`$ is obtained in practice. For this study, we have developed a new implementation of the generalized $`\overline{\mathrm{MS}}`$ formalism using Monte Carlo methods. This implementation allows the computation of differential distributions, with kinematic cuts, such as $`d\sigma ^D^{}/dpt,d\sigma ^D^{}/dQ,`$ and $`d\sigma ^D^{}/d\eta `$ which we present in Sec. 4. These results, are in qualitative agreement with available data from HERA. However, they are not as good as those of the order $`\alpha s^2`$ calculation in the 3-flavor scheme. This is to be expected for differential distributions with experimental cuts, since the (resummed) low-order calculation contains more severe approximations to the kinematics of the final state partons. For future quantitative studies, the general formalism needs to be expanded to incorporate higher-order results (adaptable from existing FFN calculations). This point is discussed in the concluding section, along with other observations.
## 2 Total Inclusive Structure Functions in the general formalism
We consider the inclusive DIS structure functions, such as $`F_2`$, focusing on the contribution of a massive quark. For definiteness, we assume that the only relevant quark with non-zero mass is the charm quark. The generic leptoproduction process is depicted in Fig. 2:
$$\mathrm{}_1+A\mathrm{}_2+X,$$
(1)
where $`A`$ is a hadron, $`\mathrm{}_{1,2}`$ are leptons, and $`X`$ represents the summed-over final state hadronic particles. Note that $`X`$ may or may not contain a visible heavy-flavor hadron.
After the calculable leptonic part of the cross-section has been factored out, we work with the hadronic process induced by the virtual vector boson $`\gamma ^{}`$ of momentum $`q`$ and polarization $`\lambda `$:
$$\gamma ^{}(q,\lambda )+A(P)X(P_X).$$
(2)
Although our considerations apply to DIS processes induced by $`W`$ and $`Z`$ as well, we shall explicitly refer to the neutral current interaction with the exchange of a virtual photon $`\gamma ^{}`$ in order to be concrete. The cross-section is expressed in terms of the hadronic tensor
$$W_{\lambda \sigma }(q,P)=\frac{1}{4\pi }\underset{X(P_X),spin}{\overline{}}P|e_\lambda ^{}J^{}|P_X(2\pi )^4\delta ^{(4)}\left(P+qP_X\right)P_X|e_\sigma J|P.$$
(3)
where $`\overline{}`$ denotes a sum over all final hadronic states. In most cases, it suffices to consider the diagonal elements of the tensor $`F^\lambda W_{\lambda \lambda }(q,P).`$
The factorization theorem in the presence of non-zero quark masses – assumed in and established to all orders in PQCD – states that the inclusive cross-section can be written as a convolution:
$$F_A^\lambda (Q^2,x,..)=\underset{a}{}f_A^a(x,\mu )\widehat{\omega }_{a,\lambda }(x,Q/\mu ,Q/m_c,\alpha _s\left(\mu \right))+𝒪(\mathrm{\Lambda }/Q)^p$$
(4)
where $`f_A^a`$ is the distribution of parton $`a`$ inside the hadron $`A`$, $`\widehat{\omega }_{a,\lambda }`$ is the perturbatively calculable hard cross-section for $`\gamma ^{}+aX`$, $`p`$ is some positive number, $`\mu `$ denotes collectively the renormalization and factorization scales, and a convolution in the $`x`$ variable is implied. The helicity structure functions $`F^\lambda `$ are simply related to the familiar $`F_{1,2,3}`$ .
The exact way that the physical structure function factorizes into the long-distance ($`f_A^a`$) and the short-distance ($`\widehat{\omega }_{a,\lambda }`$) pieces on the right-hand side of Eq. (4) depends on the scheme used to define the parton distributions. The physical structure function $`F_A^\lambda `$ should be independent of any calculational scheme; therefore, the definition of the hard cross-sections is determined by the subtraction procedure used to define the parton distributions. As discussed in the Introduction, the general formalism consists of the 3-flavor scheme at low energy scales, the 4-flavor scheme at high energy scales, and a suitably chosen transition region where matching conditions between the two schemes are applied. The precise description of these elements of the formalism is given in the Appendix. Operational definitions of quantities needed in subsequent discussions are also discussed in some detail there.
### 2.1 Previous calculations
To put the current calculation in context, we first summarize the existing calculations of leptoproduction of charm using the precise definitions given in the Appendix.
* NLO 3-flavor (3$`\alpha _s^2`$) calculation : Most dedicated calculations of heavy flavor production in recent years have been carried out in this scheme. The LO process is $`𝒪(\alpha _s^1)`$ heavy-flavor creation (HC), $`\gamma ^{}gc\overline{c}`$. The NLO processes consist of the $`𝒪(\alpha _s^2)`$ virtual corrections to $`\gamma ^{}gc\overline{c}`$ as well as the real HC $`\gamma ^{}lc\overline{c}l`$ process, where $`l`$ denotes any light parton. Cf. Fig. 3. This calculation becomes questionable when $`Qm_c`$i.e., it ceases to be “NLO” in accuracy, as indicated in Fig. 1, because the perturbative expansion is actually in $`\alpha _sln(Q^2/m_c^2)`$ for large $`Q`$.
* Zero-mass 4-flavor (ZM) (4$`\alpha _s^1;`$ $`m_c=0`$) calculation: This is the formalism used in most conventional QCD parton model calculations and popular Monte Carlo programs. It represents an approximation to the general mass (GM) 4-flavor scheme by setting $`m_c=0`$ in the hard cross-sections $`\widehat{\omega }_a(x,\frac{Q}{\mu },\frac{m_c}{Q},\mu )\stackrel{m_c/Q0}{}\widehat{\omega }_a^{m_c=0}(x,\frac{Q}{\mu },\mu ).`$ The LO contribution consists of the $`𝒪(\alpha _s^0)`$ $`\gamma ^{}cc`$ heavy-flavor excitation (HE) process. The NLO contribution consists of the $`𝒪(\alpha _s^1)`$ virtual corrections to $`\gamma ^{}cc`$ plus the real HC $`\gamma ^{}gc\overline{c}`$ and, $`\gamma ^{}cgc`$ HE processes. Cf. Fig. 4. This calculation is unreliable in the threshold region, as mentioned in the introduction.
* LO generalized $`\overline{\mathrm{MS}}`$ calculation (ACOT) : This represents the simplest implementation of the generalized formalism . It emphasizes the overlapping physics underlying the 3-flavor and 4-flavor calculations in the region not far above threshold. The 4-flavor calculation consists of HE $`\gamma ^{}cc`$ (Fig. 4a) plus $`m_c0`$ HC $`\gamma ^{}gc\overline{c}`$ (Fig. 4d), with the requisite subtraction term which removes the $`m_c`$ mass-logarithm. Mathematically, the result on $`F_2(x,Q,\mu )`$ can be shown to match that of the $`𝒪(\alpha _s^1)`$ 3-flavor scheme calculation (HC $`\gamma ^{}gc\overline{c}`$) as the (unphysical) factorization scale $`\mu `$ approaches $`m_c`$ from above . Physically, the predicted behavior of $`F_2^c`$ for $`Qm_c`$ will depend on the choice of $`\mu `$ as a function of the physical variables, cf. next section. At a high energy scale, $`\mu Qm_c`$, this calculation only approximates the 4-flavor NLO results: it contains the most important $`𝒪(\alpha _s^1)`$ term, HC $`\gamma ^{}gc\overline{c}`$ (because of the large gluon distribution and the need for matching), but does not include the smaller $`𝒪(\alpha _s^1)`$ terms represented by Fig. 4b,c. This calculation has been further studied by Kretzer and Schienbein and Krämer et al .
* Variations on the “variable flavor number” theme: In recent years, the general approach proposed in has been adopted by other groups, starting from different historical perspectives, . In contrast to the fixed flavor-number (FFN) approach, these are usually referred to as being in the variable flavor-number (VFN) scheme. This terminology has caused some confusion, since although the common theme is that of , as shown in Fig. 1, the various implementations differ considerably. Whether a particular implementation is self-consistent within the general framework of PQCD, or whether two implementations are compatible within the accuracy of the perturbative expansion, are often unclear or controversial (e.g. ) because of the complexity of the multi-scale problem and because of possible misunderstandings. It is beyond the scope of this paper to critically review these approaches. By pursuing the goals described in the introduction, we hope to provide a clearer picture of the general formalism of , hence a better basis to help address some of the controversial issues in the future.
### 2.2 The full order $`\alpha _s`$ generalized $`\overline{\mathrm{MS}}`$ calculation
The calculation of charm contribution to the inclusive structure functions reported in this paper completes the order $`\alpha _s`$ calculation in the general formalism, initiated in , including all the hard processes of Fig. 4. The additional terms, although relatively small numerically at current energies, are required to make the 4-flavor part of this calculation truly NLO at high energies, so that it becomes equivalent to the conventional zero-mass NLO theory used in most modern analyses of precision DIS data. In the remainder of this section we discuss the theoretical issues and uncertainties in the order $`\alpha _s`$ 4-flavor calculation at all energy scales, in order to address the issues highlighted at the end of the introduction, Sec. 1.
At order $`\alpha _s^1`$ the 3-flavor component of the composite scheme consists of only the (unsubtracted) $`𝒪(\alpha _s^1)`$ HC $`\gamma ^{}gc\overline{c}`$ process. The result is standard. Therefore, our main calculation concerns the *4-flavor scheme component*. At order $`\alpha _s^1`$ the right-hand side of Eq. (4) consists of three terms
$$\begin{array}{ccc}F_{A,\lambda }^{(4)}(Q^2,x,m_c,\mu )\hfill & =\hfill & f_A^c^0\widehat{\omega }_{c,\lambda }^c\hfill \\ & & +f_A^g^1\widehat{\omega }_{g,\lambda }^{c\overline{c}}\hfill \\ & & +f_A^c^1\widehat{\omega }_{c,\lambda }^{cX}\hfill \\ & & +\mathrm{light}\mathrm{parton}\mathrm{contributions},\hfill \end{array}$$
(5)
where the superscript (4) indicates that this is a 4-flavor calculation, and the *hard-scattering cross-sections* $`{}_{}{}^{i}\widehat{\omega }_{a,\lambda }^{X}`$ for the various subprocesses are calculated from the corresponding *partonic cross-sections* $`{}_{}{}^{i}\omega _{a,\lambda }^{X}`$ (without the hat) according to the procedures described in the Appendix. A description of each of the terms follows:
* The leading order $`\gamma ^{}cc`$ (Fig. 4a) partonic cross-section $`{}_{}{}^{0}\omega _{c,\lambda }^{c}`$ is infra-red safe, thus
$${}_{}{}^{0}\widehat{\omega }_{c,\lambda }^{c}=^0\omega _{c,\lambda }^c.$$
(6)
* The $`\gamma ^{}gc\overline{c}`$ (Fig. 4d) partonic cross-section contains a single power of $`\mathrm{ln}\mu ^2/m_c^2`$ which can be factorized into the charm distribution function by the subtraction
$${}_{}{}^{1}\widehat{\omega }_{g,\lambda }^{c\overline{c}}=^1\omega _{g,\lambda }^{c\overline{c}}^1\stackrel{~}{f}_g^c^0\omega _{c,\lambda }^c,$$
(7)
where the cancelling logarithm with mass-singularity resides in the $`𝒪(\alpha _s)`$ perturbative parton distribution function
$${}_{}{}^{1}\stackrel{~}{f}_{g}^{c}=(\alpha _s/2\pi )P_{gq}(x)\mathrm{ln}(\mu ^2/m_c^2)$$
(8)
* The virtual correction to $`\gamma ^{}cc`$ (Fig. 4b) plus the real $`\gamma ^{}cgc`$ (Fig. 4c) partonic process also contain $`\mathrm{ln}\left(\mu ^2/m_c^2\right)`$ terms which are factorized into the charm distribution function by the subtraction
$${}_{}{}^{1}\widehat{\omega }_{c,\lambda }^{cX}=^1\omega _{c,\lambda }^{cX}^1\stackrel{~}{f}_c^c^0\omega _{c,\lambda }^c,$$
(9)
where the logarithm appears in the $`𝒪(\alpha _s^1)`$ perturbative parton distribution function<sup>g</sup><sup>g</sup>gWe have calculated these terms keeping a finite charm quark mass. As discussed in Ref. , it would also have been consistent to calculate diagrams with an initial state charm quark using a zero charm quark mass. The errors near threshold in both methods of calculation are comparable and of order $`\alpha _s^2`$.
$${}_{}{}^{1}\stackrel{~}{f}_{c}^{c}=\frac{\alpha _s}{2\pi }\frac{4}{3}\left[\left(\frac{1+x^2}{1x}\right)\left(\mathrm{ln}\frac{\mu ^2}{m_c^2}12\mathrm{ln}(1x)\right)\right]_+.$$
(10)
Substituting Eqs. (6)-(9) into Eq. (5), the right-hand side can be re-organized as
$$\begin{array}{ccc}F_{A,\lambda }^{(4)}(Q^2,x,m_c,\mu )\hfill & =\hfill & f_A^g^1\omega _{g,\lambda }^{c\overline{c}}\hfill \\ & & +(f_A^cf_A^g^1\stackrel{~}{f}_g^cf_A^c^1\stackrel{~}{f}_c^c)^0\omega _{c,\lambda }^c\hfill \\ & & +f_A^c^1\omega _{c,\lambda }^{cX}\hfill \\ & & +\mathrm{light}\mathrm{quark}\mathrm{terms},\hfill \end{array}$$
(11)
where the $`\mathrm{ln}\left(\mu /m_c\right)`$ terms in the $`\omega _{a,\lambda }`$ factors are kept intact, and the needed subtraction terms are explicitly grouped with the leading 2$``$1 term with the same kinematics. This is the form we use for the actual numerical calculations, which we implement using a Monte Carlo approach.
It is useful to compare this calculation with the NLO 3-flavor calculation which can be written, in the same notation, as follows
$$\begin{array}{ccc}F_{A,\lambda }^{(3)}(Q^2,x,m_c,\mu )\hfill & =\hfill & f_A^g^1\omega _{g,\lambda }^{c\overline{c}}\hfill \\ & & +f_A^g^2\omega _{g,\lambda }^{gc\overline{c}}\hfill \\ & & +f_A^q^2\omega _{q,\lambda }^{qc\overline{c}}\hfill \\ & & +\mathrm{light}\mathrm{quark}\mathrm{terms}\hfill \end{array}.$$
(12)
The term common to the two schemes is $`𝒪\left(\alpha _s^1\right)`$ $`\gamma ^{}gc\overline{c}`$, which appears as the first line in both Eq. 11 and Eq. 12. For this comparison, one can consider the rest of the terms in these equations as complementary “corrections” to the first term. In particular, the last two terms in the 3-flavor formula Eq. 12 are genuine $`𝒪\left(\alpha _s^2\right)`$ corrections to the common term in the threshold region; hence are commonly referred to as NLO. But at high $`Q^2m_c^2`$, these terms contain large logarithm factors $`\mathrm{ln}Q^2/m_c^2`$ which vitiates the perturbation expansion. They are *not* NLO in this region.
In the 4-flavor calculation as organized in the form Eq. 11, the last two lines are also effectively $`𝒪\left(\alpha _s^2\right)`$. This is because the distribution $`f_A^c`$ is effectively $`𝒪\left(\alpha _s^1\right)`$ near threshold, and there is a built-in cancellation between the leading order $`\gamma ^{}cc`$ partonic cross-section and the first subtraction term in Eq. 11, as discussed in detail in Ref.. Although these $`𝒪\left(\alpha _s^2\right)`$ “correction” terms to the $`𝒪\left(\alpha _s^1\right)`$ common term (first line) do not contain the full NLO $`𝒪\left(\alpha _s^2\right)`$ corrections at threshold, they do contain all the $`𝒪\left(\alpha _s^2\right)`$ contributions which are enhanced by $`\mathrm{ln}\left(Q^2/m_c^2\right)`$ and quickly dominate as $`Q^2`$ increases. In fact, these large logarithmic terms have been resummed to all orders in perturbation theory via the DGLAP-evolved charm parton distribution contributions in the second and third lines. Therefore, in the large $`Q^2`$ region the 4-flavor result represents a true NLO calculation.
In principle, if carried out to all orders in $`\alpha _s`$, the 3-flavor and 4-flavor calculations in the form of $`F_{A,\lambda }^{(i)}(Q^2,x,m_c,\mu )`$, $`i=3,4`$, would give exactly the same prediction for all values of the arguments. The dependence on the scheme $`(i)`$ and the scale $`(\mu )`$ arises from the truncation of the perturbation series. Therefore, in order to produce a physical prediction $`F_{A,\lambda }^{phys}(Q^2,x,m_c)`$ and to relate the predictions in the two schemes, a number of additional steps must be taken. In the presence of a non-zero heavy quark mass, some of these steps are non-obvious; hence they can be the source of confusion. An explicit discussion of these elements of the calculation will make clear the flexibility as well as the uncertainties inherent in the formalism. This we do in the Appendix, as part of the more precise description of the formalism. Here, we mention only two features which are particularly relevant for the subsequent discussions of physical predictions.
First, within each scheme ($`i=3`$ or $`4`$), one needs to specify $`\mu `$ as a function of the physical variables in order to make a *physical prediction*, i.e.
$$F_{A,\lambda }^{phys}(Q^2,x,m_c)=F_{A,\lambda }^{(i)}(Q^2,x,m_c,\mu (x,Q,m_c))$$
Although there is considerable freedom in choosing $`\mu (x,Q,m_c)`$, two conditions should be met so that the prediction can be reliable: (i) $`\mu `$ must be of the order of $`Q`$ or $`m_c`$ so that PQCD applies, and (ii) $`F_{A,\lambda }^{(i)}(Q^2,x,m_c,\mu )`$ must be relatively stable with respect to variations of $`\mu `$ for the ($`x,Q`$)-range of interest. This is the well-known *scale-dependence* of any PQCD calculation. For the problem at hand, a common choice for $`\mu (x,Q,m_c)`$ is $`\sqrt{Q^2+m_c^2}`$: it represents the typical virtuality of the internal parton lines in the important subprocesses. The presence of the uncertainty associated with the choice of $`\mu (x,Q,m_c)`$ in each scheme is illustrated in Fig. 1 by the respective bands.<sup>h</sup><sup>h</sup>hIn the original ACOT paper, a different scale function $`\mu (x,Q,m_c)`$ was chosen: it was designed to enforce the condition $`F^{(4)}(Q^2,x,m_c,\mu (x,Q,m_c))F_{LO}^{(3)}`$ $`(Q^2,x,m_c,m_c)`$ as $`Qm_c.`$ In retrospect, this choice is artificial and unnecessary. The general formalism automatically ensures that, $`F^{(4)}(Q^2,x,m_c,\mu )F_{LO}^{(3)}`$ $`(Q^2,x,m_c,m_c)`$ as $`\mu m_c`$ for given $`(Q^2,x)`$ to the order which we are working; but it does not place any restriction on the behavior of $`F^{phys}(Q^2,x,m_c)`$ as $`Qm_c`$ in any given scheme. As shown in Sec. 3, the more natural choice of scale $`\mu (x,Q,m_c)=\sqrt{Q^2+m_c^2}`$ leads to much improved physical predictions.
Secondly, as shown in Fig. 1, one needs to identify an appropriate scale at which the 3-flavor and the 4-flavor scheme predictions are both reasonable and mutually comparable, so that the transition from one to the other in the composite scheme can be made smoothly. In the next section we show that, for the inclusive charm production cross-section, these conditions can be met over a rather large range of $`Q`$, extending down to near the threshold region. One possibility then is to choose the transition point (cf. Fig. 1) at a low value, close to $`m_c`$, so that in effect the 4-flavor calculation by itself covers the full range of physical interest.
In existing literature, it is already known that the 3-flavor order $`\alpha _s^2`$ (NLO) calculation can be extended to most of the currently accessible energy scales without manifest ill-effects of the large logarithms; and it agrees with data rather well. Its efficacy at very large $`Q`$ is not fully tested (cf. and results of next section). Our calculation will demonstrate the robustness of the complementary, much simpler (order $`\alpha _s`$) 4-flavor calculation. It is worthwhile pointing out that, in the 4-flavor scheme, the order $`\alpha _s`$ calculation is, in fact, also *NLO* – since the LO $`\gamma ^{}cc`$ term is of order $`\alpha _s^0`$, as is the case in the standard QCD theory of inclusive structure functions. This important point is discussed in detail in the Appendix, Sec. A.5.
## 3 Inclusive Charm Structure Function
In the previous section, the theoretical formulation of charm production in DIS is presented within the context of the totally inclusive structure functions $`F_\lambda (x,Q)`$. Although terms involving at least one charm quark in the final state are the main subject of discussion, they have been added to “light quark contributions” to form the totally inclusive structure functions, cf. Eqs. 11 and 12. The reason to do this (rather than simply talk about an “inclusive charm structure function”, say $`F_2^c)`$ is: charm quarks are not physically observable; there is no unique or obvious definition of a “$`F_2^c`$” either experimentally or theoretically. Any experimental definition necessarily depends on the procedure or prescription for tagging final state charm (analogous to the jet-algorithm for defining jets). Any theoretical definition will be scheme-dependent, and will be subjected to questions such as infra-red safety (IRS), e.g. free from $`\mathrm{ln}(m_c)`$ terms. The proper matching of the experimental and theoretical definitions is also necessary for a meaningful comparison of theory with experiment.
In this section, we shall be concerned with the comparison of our calculation of inclusive charm production structure function with previous calculations, such as that of ACOT and that of the NLO three-flavor scheme , as well as with existing experimental results. Since the energy range where these comparisons will be made is moderate, the question of infra-red safety is not a critical one – as evidenced by the general phenomenological success of the NLO 3-flavor calculation (which is not IRS in the sense used in this paper). For this purpose, it suffices to define the theoretical 4-flavor $`F_2^c`$ as the quantity consisting of terms with at least one charm quark in the final state in Eq. 11, i.e. it is the right-hand side less the light-quark contribution. At the order we are calculating, this $`F_2^c`$ is actually well-defined and infra-red safe, – there are no large logarithm terms.<sup>i</sup><sup>i</sup>iHowever, at the next order ($`\alpha _s^2`$) and beyond, $`F_2^c`$ is, strictly speaking, not totally infra-red safe by itself (in the sense that it does contain some un-cancelled $`\mathrm{log}(\frac{Q}{m_c})`$ factors); only the sum with the light-quark contributions (in Eqs.5,11) is free from such potentially large logarithms. Thus, the 4-flavor formula without the light-parton term should not be used far beyond the physical range $`Q\mu >m_c`$ where it is well-defined. This issue is discussed in Ref. , where a cutoff parameter is used to remove the “soft” parts of the charm contribution from $`F_2^c`$ and to combine them with the light-parton terms. In addition, since the charm mass is just about at the scale for PQCD to become valid and, as discussed at the end of the previous section, the transition scale can in practice be chosen close to $`m_c`$, *the generalized $`\overline{\mathrm{MS}}`$ calculation reduces in practice to a 4-flavor (general mass) calculation*.
#### The 4-flavor NLO calculation:
The parton distribution functions $`f_A^a(x,\mu )`$ needed for this calculation must be defined in the same renormalization scheme as the hard cross-sections used. In the 4-flavor general mass calculation, most CTEQ parton distributions have been defined in the ACOT scheme as described in the Appendix, with the matching scales chosen at the heavy quark masses. In particular, the CTEQ5HQ set is generated with non-zero heavy quark masses in the hard cross-sections, hence it is the relevant one for our application. This set is evolved using NLO evolution, which is appropriate to our calculation which is NLO at high energies. For the numerical results presented below, we use an updated version of CTEQ5HQ . In order to compare with results from the LO and NLO 3-flavor scheme, we need the corresponding parton distributions in the 3-flavor scheme. The CTEQ5F3 distributions satisfy this need, since they are obtained from global analysis of the same data sets as the CTEQ5HQ set, but with parton distributions and hard cross-sections in the 3-flavor scheme.
The calculation of the hard cross-section at $`𝒪(\alpha _s)`$ is based on Eqs. 7 and 9. In the expression for the overall structure functions, the combination $`f_A^g^1\omega _{g,\lambda }^{c\overline{c}}`$ $`+(f_A^cf_A^g^1\stackrel{~}{f}_g^c`$) $`^0\omega _{c,\lambda }^c,`$ due to the subprocesses $`\gamma ^{}cc`$ and $`\gamma ^{}gc\overline{c},`$ comprise the original ACOT calculation . With non-zero $`m_c,`$ they are all finite. The implementation of these terms in the new Monte Carlo calculation is straightforward. We have verified that the new MC program reproduces the original ACOT results in detail. Illustrations of the relative contributions of the various terms can be found in Ref. . Also, the smooth transition of the results of the general formalism from the 3-flavor one near threshold to the 4-flavor one at high energies is described in that paper.
The additional contributions $`f_A^c^1\omega _{c,\lambda }^{cX}f_A^c^1\stackrel{~}{f}_c^c^0\omega _{c,\lambda }^c,`$ due to the $`\gamma ^{}cgc`$ subprocess and virtual correction to the Born $`\gamma ^{}cc`$ term, are new for this calculation. They contain soft divergences which must be cancelled. In this Monte Carlo implementation, we use the phase-space splicing method to achieve the proper cancellation of the soft divergences between the real and virtual parts. Details of this calculation are contained in . For double-checking, we have independently implemented an analytic calculation based on the formulas of Hoffmann and Moore . The two calculations agree quite well with each other over the full phase space, with the exception of small values of $`Q/m_c`$. The difference could be attributed to a different treatment of the massive charm quark kinematics adopted by Ref. in deriving their formulas. This effect goes away when $`m_c`$ is small compared to $`Q,`$ as expected.<sup>j</sup><sup>j</sup>jThe authors of Ref. have also identified some differences between their recent calculation with Ref.. We thank S. Kretzer for providing us with some details of this comparison. Since the difference in question is numerically not significant, none of the results presented below are affected by this problem.
#### Results and Comparison with order $`\alpha _s^2`$ 3-flavor calculation:
We now present some typical numerical results on the theoretical $`F_2^c(x,Q)`$ obtained in our order $`\alpha _s`$ 4-flavor calculation compared to those of order $`\alpha _s^2`$ 3-flavor scheme. The 3-flavor results were obtained from the parametrization of Ref. .<sup>k</sup><sup>k</sup>kWe thank Brian Harris for furnishing us with his interface to this program. Each calculation corresponds to a different way of organizing the perturbation series, hence has its natural region of applicability, as discussed in the previous sections. To do a meaningful comparison, it is important to take into account the estimated uncertainties of each calculation. Following the example of Fig. 1, we plot a band for each of the predictions, obtained with two common choices of the renormalization and factorization scales: $`\mu =c\sqrt{Q^2+m_c^2}`$ with $`c=0.5`$ and $`2.0`$. (This choice is more natural than that of the original ACOT paper, c.f. footnote h.) The 3-flavor calculation requires parton distributions defined in the same scheme. We use the CTEQ5F3 set. The results vary in appearance, depending on the kinematic variables. Two representative plots relevant for the HERA measurements are shown in Fig. 5 where $`F_2^c(x,Q)`$ is plotted against $`Q`$ for two values of $`x`$: (a) $`x=0.01;`$ and (b) $`x=0.0001.`$ These constitute real examples of the cartoonistic Fig. 1, which is designed to emphasize the underlying ideas.
These plots show that the overlapping region of the two schemes is quite wide for $`x=10^4`$; and the overall behavior conforms with expectations. For $`x=10^2`$, the overlapping region is more limited and it is confined to low $`Q`$ values. The 3-flavor results fall below the 4-flavor ones at large values of $`Q`$ where the latter should be more reliable. Compared to current data, both are within the experimental range; the 4-flavor results are closer to the data points, cf. next subsection, although the uncertainties of the 4-flavor calculation are fairly large for $`x=10^4`$.
Since the order $`\alpha _s`$ 4-flavor results are comparable to the order $`\alpha _s^2`$ ones at energy scales close to the threshold in both cases (and for other values of $`x`$), it is reasonable to choose the transition scale at a relatively low value, as mentioned earlier in the paper. The band representing the 3-flavor calculation does become wider at large Q for $`x=10^2`$ (where the absolute values also are lower than the 4-flavor calculation and data); but for $`x=10^4`$, it remains quite narrow. Thus, the theoretically infra-red unsafe logarithms, $`\mathrm{ln}^{1,2}(\mu /m_c)`$, do not seem to cause serious problems, at least for very low $`x`$.
#### Results Compared to recent Zeus data:
The general agreement between existing data on charm production with order $`\alpha _s^2`$ 3-flavor calculations, using the scale choice $`\mu =\sqrt{Q^2+m_c^2}`$, is well known. It is of interest to compare the same data on “inclusive charm production structure function” $`F_2^c`$ with our order $`\alpha _s`$ 4-flavor calculation. Fig. 6 compares the results of our calculation, using the same scale choice, with data from the recent ZEUS data. The agreement is clearly excellent. Data also agree with the order $`\alpha _s^2`$ 3-flavor calculation as shown in . This higher order calculation is obviously much more elaborate than the order $`\alpha _s`$ 4-flavor calculation presented here. We see that, for inclusive cross-sections, the resummation of the $`\alpha _s^nln^n(\mu /m_c)`$ terms into the charm distribution function $`f^c(x,\mu )`$ in the 4-flavor scheme offers a more efficient way to organize the perturbative series, resulting in an effective NLO calculation already at order $`\alpha _s`$, cf. Sec. A.5.
## 4 Semi-inclusive Cross Sections with Tagged Charm Hadrons
### 4.1 General considerations
Now we consider semi-inclusive cross sections, with a charm hadron tagged in the final state. Naively, to compute this cross section, one simply convolves the cross sections for parton final states, Eqs. 11 and 12, with a suitable fragmentation function of partons into the final charm hadron. However, the factorization of the final state particles through fragmentation functions is only rigorously defined in the limit $`Q^2m_c^2`$. Thus, the treatment of tagged charm particles in the final state can only be systematically applied at high energies, using the 4-flavor scheme. However, it is a common practice to introduce fragmentation functions into charm hadrons even in the 3-flavor scheme, and for energies not far above threshold. This approach should be considered a convenient phenomenological model of hadronization, perhaps adequate for current experimental accuracy, rather than rigorous theory. We follow this practice in our calculation, employing fragmentation functions over the full range of $`Q^2`$, while maintaining the correct factorization-scheme implementation at high $`Q^2`$.
The fragmentation functions $`d_a^H(x,\mu )`$ obey the standard (mass-independent) QCD evolution equations, and are determined from suitable initial functions at some given scale $`\mu =Q_0`$, of the order of $`m_c`$. Following Mele and Nason , for a given final-state charm hadron $`H,`$ we write
$$d_a^H(x,Q_0)=d_a^c(x,Q_0)D_c^H(x,Q_0)$$
(13)
where the partonic charm fragmentation functions $`\left\{d_a^c;a=l,c\right\}`$ are perturbatively calculable, and $`D_c^H(x,Q_0)`$ is considered non-perturbative and is to be obtained by comparison with experiment.
For the perturbatively calculable fragmentation functions, Ref. gives, to order $`\alpha _s`$:
$`d_c^c(x,Q_0)`$ $`=`$ $`\delta (1x)+{\displaystyle \frac{\alpha _s(Q_0)C_F}{2\pi }}\left[{\displaystyle \frac{1+x^2}{1x}}\left(\mathrm{ln}{\displaystyle \frac{Q_0^2}{m_c^2}}2\mathrm{ln}(1x)1\right)\right]_+`$
$`d_g^c(x,Q_0)`$ $`=`$ $`{\displaystyle \frac{\alpha _s(Q_0)T_F}{2\pi }}(x^2+(1x)^2)\mathrm{ln}{\displaystyle \frac{Q_0^2}{m_c^2}}`$ (14)
$`d_{q,\overline{q},\overline{c}}^c(x,Q_0)`$ $`=`$ $`0`$
where $`T_F=1/2`$ and $`C_F=4/3`$. In keeping with the choice of the matching scale in our overall calculation, we choose $`Q_0=m_c`$ for convenience in this paper.
For the non-perturbative charm quark into charmed mesons fragmentation function $`D_c^H(z),`$ we used the conventional Peterson form ,
$$D_c^H(z)=\frac{A}{z[11/zϵ/(1z)]^2},$$
(15)
For the charm meson $`D^{()},`$ which will be our focus because of available experimental data, we take $`ϵ=0.02,`$ cf. , and a value for $`A`$ such that the branching fraction $`B(cD^{})=0.22`$ . We note that, although the perturbative fragmentation functions, Eq. 14, contain singular (generalized) functions, the overall parton-to-charm-meson fragmentation functions $`d_a^H(x,\mu ),`$ Eq. 13, are well behaved after convolution with the above non-perturbative fragmentation function $`D_c^{D^{()}}(z)`$.
In principle, after evolving to high enough $`Q^2`$ so that $`\alpha _s\mathrm{log}(Q^2/m_c^2)`$ is of order one, all of the fragmentation functions $`d_c^H,d_g^H,d_{q,\overline{q},\overline{c}}^H`$ eventually become of the same size. In practice, however, at HERA energies we find $`d_c^Hd_g^Hd_{q,\overline{q},\overline{c}}^H`$. For currently required accuracy, it suffices to keep only the charm-to-hadron contributions, proportional to $`d_c^H`$. In this approximation, our calculation of the cross section with a tagged charm hadron can be written in the 4-flavor scheme as
$$\begin{array}{ccc}F_{A,\lambda }^H(Q^2,x,..)\hfill & =\hfill & f_A^g^1\omega _{g,\lambda }^{c\overline{c}}d_c^H\hfill \\ & & +(f_A^cf_A^g^1\stackrel{~}{f}_g^cf_A^c^1\stackrel{~}{f}_c^c)^0\omega _{c,\lambda }^cd_c^H\hfill \\ & & +f_A^c^1\omega _{c,\lambda }^{cX}d_c^H.\hfill \end{array}$$
(16)
To ensure that this calculation is adequate, we have also calculated the contribution from one of the more important remaining subprocesses: gluon fragmentation in an order $`\alpha _s`$ light parton hard scattering, i.e. $`\gamma ^{}qgq;gH`$. It is given by: $`f_A^q^1\widehat{\omega }_{q,\lambda }^{qg}d_g^H`$, where $`d_g^H`$ is the gluon fragmentation function computed from Eqs. 13 and 14. We have verified that its contribution remains small throughout the current energy range. It becomes more noticeable only in the large $`Q`$ limit. However, in this limit, the gluon fragmentation function term is not infra-red safe by itself. To insure consistency at high energies, one needs to include a full set of infra-red safe higher order subprocesses along with it. The full calculation is more appropriately considered as part of the next order project.
At the same level of accuracy, it is also reasonable to ignore the evolution of $`d_a^H(x,\mu (x,Q))`$, since the effect of QCD evolution is not significant over the currently accessible HERA $`Q`$ range. One can use the un-evolved $`d_c^H(x,\mu =Q_0)`$ in place of the fully evolved $`d_c^H(x,\mu (x,Q,m_c))`$ with much gain in efficiency of calculation and little sacrifice in accuracy. The error incurred is of the same order as that incurred by neglecting the subleading fragmentation functions, $`d_g^H`$ and $`d_{q,\overline{q},\overline{c}}^H`$ ; and the comments on accounting for $`\mathrm{ln}(m_c)`$ factors made there also apply here. To be sure about this, we have performed the calculation both with and without evolving $`d_c^H.`$ The difference is indeed small. Therefore, in all our subsequent plots we shall include only the direct charm-to-hadron contributions of Eq. (16) with the un-evolved fragmentation function $`d_c^H(x,Q_0)`$.
### 4.2 Differential distributions
We employ the Monte Carlo method to carry out the numerical phase-space integration; the new program package is implemented in the C++ programming language. Therefore, we can generate differential distributions involving final-state charm mesons, incorporating kinematic cuts appropriate for specific experimental measurements, in addition to fully inclusive cross-sections.
In working with on-mass-shell heavy flavor quarks and hadrons in the parton language, there is an ambiguity in defining the momentum fraction variables $`x(z)`$ for the parton distribution (fragmentation) function. This problem arises in all schemes; and it goes away at high energies ($`Qm_c`$), where the parton picture becomes accurate. Following the modern practice in proofs of factorization, we define the momentum fraction variables as ratios of the relevant light-cone momentum components, e.g.$`p_D^+=zp_c^+`$ for fragmentation of a charm quark into a $`D`$ meson. Other authors, e.g. , use the prescription $`\stackrel{}{p}_D=z\stackrel{}{p}_c`$ and adjust the energy variable to enforce the mass-shell condition. At moderate energies, any noticeable differences in results due to the choice of this prescription signals that the calculation using fragmentation functions is outside the region of applicability of the parton formalism. We have verified that the results presented below are insensitive to the choice between the two prescriptions.<sup>l</sup><sup>l</sup>lThere are some differential distributions, especially those associated with the unphysical partons (such as momentum fraction carried by the charm quark, sometimes seen in the literature) which are more sensitive to the choice of definition of the momentum fraction variable.
The QCD formula, Eq. 16, contains three scale choices in principle: the renormalization scale, the factorization scale and the fragmentation scale. For simplicity, we choose the same energy scale $`\mu (x,Q,m_c)`$ for all three. As in the case of the inclusive $`F_2^c(x,Q)`$, for results shown below, we choose the simple function $`\mu =c(Q^2+m_c^2)^{1/2}`$, which characterizes the typical virtuality of the process. The constant $`c`$ is of order 1; and is varied over same range when we try to estimate the scale-dependence of the physical predictions. The magnitude of the charm cross-section is sensitive to the value of $`m_c`$. For results presented here, we use $`m_c^{\overline{\mathrm{MS}}}=1.3GeV`$, the value used in the CTEQ5HQ parton distribution analysis (which is in the middle of the range given by the PDG review).
Fig. 7 shows plots of four differential distributions for $`D^{}`$ production at HERA, calculated using the NLO ($`\alpha _s`$) generalized $`\overline{\mathrm{MS}}`$ 4-flavor formalism described above. The kinematic variables and their ranges correspond to those of the 1996-97 ZEUS data: $`1<Q^2<600`$GeV$`{}_{}{}^{2};\mathrm{\hspace{0.33em}0.02}<y<0.7;\mathrm{\hspace{0.33em}1.5}<p_T^D^{}<15`$GeV ; $`\left|\eta _D^{}\right|<1.5.`$ Each distribution contains two curves obtained with two values of the constant $`c=0.5,1`$ in the definition of the scale parameter described above. These predictions, using the CTEQ5HQ parton distributions, are compared to the ZEUS data . We observe a rather large scale dependence in these results. This is not surprising, since the compensation among the various subprocesses which underlie the scale-independence of the physics predictions (up to some order of perturbation theory), strictly speaking, only apply to the inclusive cross-section. The experimental kinematical cuts implemented in these exclusive calculations to some extent undermine the mutual cancellation between diagrams which are necessary for relatively scale-independent predictions. For example, the order $`\alpha _s^0`$ $`\gamma ^{}cc`$ HE term (which resums the logarithms arising from the near-collinear configurations of an infinite tower of higher-order diagrams) implements the full contribution in collinear kinematics, a clear over-simplification.
Keeping this fact in mind, and with current relatively large experimental errors, the results of Fig. 7 can be considered rather encouraging: the $`Q^2`$ and $`p_T`$ distributions show very good general agreement; while the $`W`$ and $`\eta _D`$ distributions are “in the right ball park”, the shapes are too scale-dependent to allow for meaningful “predictions”. (A specific choice of scale, in between the two shown, will actually yield theory curves in reasonable agreement with data, within errors.) In order to make genuine predictions on differential distributions in the 4-flavor scheme, it is necessary to extend the calculation to order $`\alpha _s^2`$, which would be NNLO in the 4-flavor scheme. This can be done by transforming already available NLO results for 3-flavor calculations into the 4-flavor scheme. At the same order in $`\alpha _s`$, the 4-flavor scheme calculation is, of course, more involved than the (NLO) 3-flavor one because of the need for including the necessary subtraction terms in a NNLO calculation – such as those appearing in Eqs. 7, 9, and 11.
The calculation of these differential distributions in the $`\alpha _s^2`$ 3-flavor FFN scheme was carried out by . Generally good agreement between these calculations and the recent ZEUS data, using specific parton distributions, scale choices, etc. has been reported in Ref. . Although the dependence of the predictions to the scale choice was not discussed in this comparison, it is relatively mild, according to . This is to be expected, because the sensitivity to cuts is reduced with a better approximation of the final state particle configurations provided by the order $`\alpha _s^2`$ calculation. This fact implies greater predictive power for the differential distributions than the order $`\alpha _s`$ 4-flavor calculation.
The $`\eta _D^{}`$ distribution in both the 3-flavor and the 4-flavor calculations appear to differ in shape compared to the existing data points. This could be due to the inadequacy of applying the fragmentation function approach at less than asymptotic region, as discussed in the beginning of this section. In particular, if the $`D^{}`$ is not collinear to the parton, as assumed in this approach, the rapidity distribution will be affected. More extensive study of this effect is obviously needed.
## 5 Conclusions
In this work we have shown how the generalized $`\overline{\mathrm{MS}}`$ formalism can be used to calculate the production of a heavy quark over a wide range of energies, from threshold to high $`Q^2`$. For charm production, the formalism consists of a 3-flavor scheme calculation at low energies, a 4-flavor scheme calculation at high energies, a matching condition between the two schemes, and a transition scale chosen at which one switches between the two schemes. Specifically, we have extended the original ACOT calculation for charm production at HERA by adding those terms which are necessary to bring the 4-flavor part of the calculation to NLO accuracy at high energies. This brings our calculation to the same level of accuracy as the other theoretical inputs to the CTEQ and MRS global QCD analyses.
The generalized $`\overline{\mathrm{MS}}`$ formalism is ideally suited for inclusive calculations, such as $`F_2(x,Q)`$ in DIS. We have shown that for a physically-motivated choice of the renormalization/factorization scale, $`\mu `$, the transition scale can be chosen rather low, so that our 4-flavor scheme can be used in the HERA energy range with excellent agreement compared to data, and with considerable economy compared to the more elaborate NLO 3-flavor calculation. The calculation is much simpler and the resulting program runs faster. Furthermore, the NLO corrections are much smaller in the 4-flavor calculation. This suggests that the resummation involved in the charm quark distribution function picks up the most important higher-order corrections even at modest energies. It also indicates that the perturbation series in the 4-flavor calculation is very well-behaved.
Our calculation has been implemented as a Monte Carlo program, so that we have also calculated differential distributions for exclusive charm final states. By incorporating experimental cuts in the Monte Carlo we are able to ensure that our calculation is compared directly with the data, without the need for any theoretical extrapolation to all of phase space. The agreement with HERA data is reasonable. However, in this case the 3-flavor NLO calculation has somewhat of an edge, because there is no approximation on the kinematics as is necessary in the resummation used in the 4-flavor calculation.
The comparison between available data and the order $`\alpha _s`$ 4-flavor and order $`\alpha _s^2`$ 3-flavor calculations discussed in the last two sections demonstrates the complementary nature of the two schemes – both with regard to the kinematic regions and to the physical quantities they are suitable for. With more abundant and more precise experimental data on charm production, the composite generalized $`\overline{\mathrm{MS}}`$ formalism, which encompasses both, will be needed to make reliable comparisons. With this in mind we are now ready to extend the 4-flavor calculation to include the necessary $`𝒪(\alpha _s^2)`$ terms, along with the $`𝒪(\alpha _s^2)`$ matching conditions. Such a calculation would include all the advantages the current calculation has in addition to the advantages of the current 3-flavor NLO calculation.
Finally, the generalized $`\overline{\mathrm{MS}}`$ formalism provides a framework to extract information on the gluon and the charm distributions of the nucleon – to the same accuracy as the overall NLO global QCD analysis based on total DIS structure functions and other hard processes. Having established that our calculation does a reasonable job in describing the existing HERA data, we are now in a position to explore the question of whether the proton contains a non-perturbative charm component . Certainly, it must at some level, so the real question is how small is it? The handle on the charm distribution is unique to the generalized formalism, since the fixed 3-flavor scheme does not allow the charm parton as an independent degree of freedom.<sup>m</sup><sup>m</sup>mThe CTEQ5HQ parton distributions used in our calculations here contains charm partons, but does not use an independent non-perturbative charm component as input. This feature is not inherent to the formalism. We note that there have been recent phenomenological studies of “intrinsic charm” which take the theoretical cross-section to be the simple sum of the 3-flavor FFN scheme formulas and an intrinsic charm contribution by the heavy-quark excitation mechanism . This approach cannot be internally consistent, because the 3-flavor calculation assumes parton evolution with no charm quark distribution, while “intrinsic charm” explicitly requires one. The intimate interplay between the hard matrix elements and parton evolution is consistently incorporated only in the generalized $`\overline{\mathrm{MS}}`$ scheme. Although there has been a recent effort to examine this problem incorporating some of the ideas of the general scheme , a fully consistent study, preferably based on more extensive data, still awaits to be done.<sup>n</sup><sup>n</sup>nBecause all active partons are coupled, the inclusion of a non-perturbative charm component of the nucleon will affect all parton distributions. Hence, it will influence the global fitting of all available data sets, including the extensive DIS sets. If a reliable conclusion is to be drawn, it will not suffice to combine a subset of existing parton distributions (in a conventional scheme) with modified charm and gluon distributions (in the new scheme), as is done in using GRV and MRST distributions. Phenomenologically, such a combination can upset the original good global fit to existing data, particularly the precision DIS data. Among theoretical problems, the NLO evolution equation and the momentum sum rule will not be observed, due to the mixing of these hybrid parton distributions.
### Acknowledgments
We would like to thank Brian Harris for many helpful discussions and especially for providing us with the code for the NLO 3-flavor calculations, and for verification of the numerical results. We would also like to thank John Collins and Fredrick Olness for many discussions and help; Jim Whitmore for assistance about the ZEUS data; and Jack Smith for discussions on the extensive order $`\alpha _s^2`$ calculations carried out by the Stony-Brook/Leiden group and for clarification of the proper references to their work. This work was supported in part by the US National Science Foundation under grants PHY9802564 and PHY9722144.
## Appendix A Formalism for Inclusive DIS Cross Section
In this appendix, we give a concise and careful presentation of the general formalism, including the relatively simple operational procedure for calculating its various components. We hope this will help fill the gap between the original, relatively sketchy, ACOT paper and the recent, more technical, all-order proof of the formalism by Collins .
The basis for all discussions is the factorization theorem in the presence of non-zero quark masses , Eq. 4. To establish this theorem in PQCD (cf. Eq. 17 below), and to give precise meaning to the various factors, one must work with *partonic cross-sections* and *parton distributions inside partons*, rather than the corresponding physical quantities which appear in Eq. 4. In this concise summary, we proceed as follows: (i) spell out the operational procedure to establish the factorization formula and calculate the hard cross-sections in any scheme; (ii) describe the specifics of the 3-flavor and 4-flavor schemes respectively; (iii) discuss the matching conditions between the two; and (iv) consider the transition from one to the other in the composite scheme, which constitutes the general formalism of Collins . We finish with some remarks on the meaning of “LO” and “NLO” calculations in different schemes.
### A.1 Procedure to define the factorization scheme and calculate the hard cross-sections
Given the QCD Lagrangian, with non-zero masses for the heavy quarks, one arrives at the general factorization formula for partonic cross-sections and parton distributions as follows: <sup>o</sup><sup>o</sup>oThis procedure may sound familiar because it is based on the basic principles of PQCD. We include it here because the details, especially concerning the heavy quark mass dependence, are quite distinct from conventional practices. Thus, it is essential to explicitly spell out the steps involved.
* (i) Start with a set of relevant *partonic structure functions* $`\omega _a`$ similar to the left-hand side of Eq. (4) but with on-shell parton targets and calculate them in perturbation theory in a given renormalization scheme (i.e. with specific ultra-violet counter-terms) to a given order in $`\alpha _s`$. The result $`\omega _a(\frac{Q}{\mu },x,\frac{m_c}{\mu },\frac{1}{ϵ},\alpha _s\left(\mu \right))`$ will depend on the renormalization scale $`\mu `$ and will contain collinear singularities (represented by $`\frac{1}{ϵ}`$) as well as potentially large logarithm terms of the form $`(\alpha _s\mathrm{ln}(\frac{\mu }{m_c}))^n`$. (For simplicity we do not differentiate between the renormalization scale $`\mu _R`$ and the factorization scale $`\mu _f,`$ both of which are taken in practice to be of order $`Q`$.)
* (ii) Independently, calculate the set of process-independent *perturbative partonic distribution functions* $`\stackrel{~}{f}_a^b`$ in the same renormalization scheme, using either the (moment-space) operator-product expansion or, equivalently, the ($`x`$-space) bi-local operator definition of the distribution functions. Both ultra-violet and collinear singularities appear in this calculation. The ultra-violet singularities are removed by additional counter-terms which, along with the coupling constant renormalization counter-terms, define the factorization scheme. The result takes the form $`\stackrel{~}{f}_a^b(x,\frac{m_c}{\mu },\frac{1}{ϵ},\alpha _s\left(\mu \right))`$.
* (iii) Confirm that all collinear singularities in the form of $`\frac{1}{ϵ}`$ terms, appearing in $`\omega _a(\frac{Q}{\mu },x,\frac{m_c}{\mu },\frac{1}{ϵ},\alpha _s\left(\mu \right))`$ appear in the universal form given in the process-independent functions $`\stackrel{~}{f}_a^b(x,\frac{m_c}{\mu },\frac{1}{ϵ},\alpha _s\left(\mu \right))`$, so that they can be factorized out in the manner of Eq. (4),
$$\omega _a(\frac{Q}{\mu },x,\frac{m_c}{\mu },\frac{1}{ϵ},\alpha _s\left(\mu \right))=\underset{b}{}\stackrel{~}{f}_a^b(x,\frac{m_c}{\mu },\frac{1}{ϵ},\alpha _s\left(\mu \right))\widehat{\omega }_b(\frac{Q}{\mu },x,\frac{m_c}{\mu },\alpha _s\left(\mu \right)),$$
(17)
with $`\widehat{\omega }_a`$ being fully infra-red safe in the sense that it is free of all $`\frac{1}{ϵ}`$ dependence. In the 4-flavor scheme defined below, the functions $`\stackrel{~}{f}_a^b(x,\frac{m_c}{\mu },\frac{1}{ϵ},\alpha _s\left(\mu \right))`$ will also contain the same large logarithmic terms as $`\omega _a(\frac{Q}{\mu },x,\frac{m_c}{\mu },\frac{1}{ϵ},\alpha _s\left(\mu \right))`$, so that these too factorize in Eq. (17) with the result that $`\widehat{\omega }_a`$ is free of all $`(\alpha _s\mathrm{ln}(\frac{\mu }{m_c}))^n`$ terms, and it is well-behaved as $`m_c0.`$
* (iv) Systematically invert Eq. (17) to solve for the set of finite hard cross-sections $`\widehat{\omega }_a,`$ which are then used in Eq. (4) for calculating physical structure functions.
There are two points to note: (a) The inversion of Eq. (17) order-by-order in the perturbation series is equivalent to *subtracting* the singularities contained in $`\stackrel{~}{f}_a^b`$ from $`\omega _a;`$ (b) There is no need to set the quark mass(es) to zero anywhere in the above procedure.
In the following, we apply the above procedure to define the two simple renormalization schemes, involving *3 or 4 active quark flavors,* which underlies the general approach of Refs. ; and combine them to define the latter in the subsection under the heading of *the generalized $`\overline{\mathrm{MS}}`$ formalism*. These discussions are applicable to all orders in perturbation theory. Throughout these discussions, the three quarks $`\{u,d,s\},`$ with masses comparable to or less than $`\mathrm{\Lambda },`$ will be referred to as *light quarks*, and denoted collectively by $`q`$. The collection of light quarks plus the gluon $`g`$ will be referred to as *light partons*, and denoted by $`l`$. As mentioned earlier, although the formalism applies to all heavy quarks $`\{c,b,t\},`$ we shall use the case of charm as a generic representative, for concreteness and clarity – hence the 3- and 4-flavors. Because the real charm quark mass $`m_c`$ is not large compared to the on-set of the region of applicability of PQCD, the 4-flavor scheme plays a more prominent role in practical applications discussed in the main body of this paper. For a heavier quark, the two corresponding schemes and their proper matching, as discussed in the rest of this (theoretical) section, will be more relevant.
### A.2 Three-flavor Scheme
The 3-flavor scheme is precisely defined by choosing to work with only 3 active quark flavors, consisting of the light quarks, and using the subtraction procedure of Ref. . The prescription for subtracting ultra-violet divergences encountered in the calculation of the partonic structure functions depends on the particle that produces the divergence. Broadly speaking, divergences due to the light partons $`l,`$ are removed using $`\overline{\text{MS}}`$ counter terms, whereas those due to the charm quark $`c`$ are removed by BPH zero-momentum subtraction counter terms. The precise definition can be found in Ref. . This ultra-violet subtraction scheme has the nice feature that the charm quark explicitly decouples as its mass becomes large. In particular, the operators which make up the charm quark distribution function are suppressed by powers of order $`\mathrm{\Lambda }^2/m_c^2`$. Since these terms are power-suppressed in the “heavy quark” mass, they are usually excluded from the 3-flavor scheme parton picture.
In practice then the partonic calculations in this scheme are done by considering diagrams where the massive charm quark can only appear in the final state, and there are no charm quark distribution functions, cf. Fig. 3. The light parton distributions always evolve according to the 3-flavor DGLAP equation, for all values of the renormalization scale $`\mu `$—both below and above the heavy quark production threshold. The parton distribution functions defined in this scheme will be restricted to the light parton $`l=\{g,q,\overline{q}\}`$ sector, and they will be denoted by $`{}_{}{}^{3}f_{A}^{l}`$. In the perturbative calculation, $`{}_{}{}^{3}\stackrel{~}{f}_{l}^{l^{}}`$ contains $`ϵ^1`$ pole terms which are due to collinear singularities. The lowest order (LO, $`𝒪(\alpha _s^1)`$) partonic process in which the charm quark appears in this scheme is the $`\gamma ^{}gc\overline{c}`$ “heavy-flavor creation” (HC) process (also known as boson-gluon fusion), corresponding to the diagrams of Fig.(3a). The associated partonic structure function, denoted by $`\omega _g^{c\overline{c}},`$ is finite. The next-to-leading order (NLO) contribution includes the 1-loop virtual corrections to $`\gamma ^{}gc\overline{c}`$ (cf. Fig.(3b)), plus the real partonic HC processes $`\gamma ^{}lc\overline{c}l`$ (cf. Fig.(3c)). The collinear divergences which appear in the calculation of the $`𝒪(\alpha _s^2)`$ partonic structure functions $`{}_{}{}^{3}\omega _{g}^{c\overline{c}}`$ and $`{}_{}{}^{3}\omega _{l}^{c\overline{c}l}`$ arise from splitting of massless light partons in the collinear configuration, and take the form of $`ϵ^1`$ pole terms, precisely corresponding to those appearing in $`{}_{}{}^{3}\stackrel{~}{f}_{l}^{l^{}}`$ mentioned above. That is, the partonic structure functions have the factorized structure shown in Eq. (4), and the hard cross-section functions $`\widehat{\omega }_l`$ will be free from $`ϵ^1`$ collinear singularities. This 3-flavor scheme is the one used by Ref. to calculate charm production to NLO, i.e. $`𝒪(\alpha _s^2)`$.<sup>p</sup><sup>p</sup>pTo be consistent, the virtual correction to the process $`\gamma ^{}qqg`$, which contains a charm quark loop, must also be included at this order in the 3-flavor scheme calculation of the total inclusive structure functions. .
At high energies the hard cross-sections calculated in this scheme contain powers of $`\mathrm{ln}(Q^2/m_c^2),`$ as mentioned in the introduction. The perturbative expansion should be accurate at energy scales not too far above threshold, or $`Q^2m_c^2,`$ where $`\mathrm{ln}(Q^2/m_c^2)`$ is of order $`1`$. However, at high $`Q^2m_c`$ the perturbative expansion parameter is effectively $`\alpha _s\mathrm{ln}(Q^2/m_c^2),`$ and the large logarithm factor spoils the convergence of the perturbative series. In other words, the “hard cross-sections” $`\widehat{\omega }_a`$ defined in this scheme are finite, but *not infra-red safe* in the limit $`\frac{m_c}{Q}0`$.
### A.3 Four-flavor scheme with non-zero $`m_c`$
In order to better deal with these logarithms at high energies it is more useful to use the *4-flavor scheme*, in which the renormalization of $`\omega _a`$ and $`\stackrel{~}{f}_a^b`$ is carried out using dimensional regularization and the $`\overline{\text{MS}}`$ counter terms for all Feynman diagrams, *while keeping the full quark mass dependence* ($`m_c`$) of the Lagrangian.
Charm distribution functions calculated in this scheme, $`{}_{}{}^{4}\stackrel{~}{f}_{a}^{c}`$ ($`a=l,c`$), are not suppressed as in the 3-flavor scheme, but contain powers of $`\mathrm{ln}(m_c/\mu )`$, along with possible $`ϵ^1`$ poles. Because of the different subtraction procedures used in the two schemes, even the light parton distributions $`{}_{}{}^{4}\stackrel{~}{f}_{l}^{l^{}}`$ will differ from $`{}_{}{}^{3}\stackrel{~}{f}_{l}^{l^{}}`$ by a finite renormalization in general. (We will return to this point later.) Because renormalization constants in the $`\overline{\text{MS}}`$ subtraction procedure are independent of mass, the evolution kernels for the $`{}_{}{}^{4}\stackrel{~}{f}_{a}^{b}`$ parton distributions will be the same as the corresponding ones in the familiar zero-mass 4-flavor case. This is a significant convenience. The perturbative parton distribution functions $`{}_{}{}^{4}\stackrel{~}{f}_{a}^{b}`$ have been calculated to NLO in Ref. .
Since charm also has a parton interpretation in this scheme, the set of partonic processes are expanded to include those involving charm initial states. The LO partonic process that involves the charm quark in the 4-flavor scheme is the $`\gamma ^{}cc`$ “heavy-quark excitation” (HE) process (Fig.(4a)). NLO charm quark contributions in the 4-flavor scheme come from the 1-loop virtual corrections to HE $`\gamma ^{}cc`$ (Fig.(4b)), and from the real HE $`\gamma ^{}cgc`$ and HC $`\gamma ^{}gc\overline{c}`$ processes (Fig.(4c,d)). Partonic structure functions $`\omega _a`$ calculated beyond LO in this subtraction scheme contain both $`ϵ^1`$ poles (due to collinear singularities associated with light degrees of freedom) and powers of mass-logarithms, $`\mathrm{ln}(Q/m_c)`$, (due to collinear configurations associated with the heavy degree of freedom), just as in the 3-flavor scheme. The important difference compared to the latter case is that these potentially large logarithm terms also appear in the 4-flavor parton distributions $`{}_{}{}^{4}\stackrel{~}{f}_{a}^{b}`$. Consequently, *they are systematically factored out from* $`\omega _a`$ when we obtain the hard cross-sections $`\widehat{\omega }_a`$ by inverting the factorization formula Eq. (17). The charm distribution function represents the *resummed* contribution of all the large (infra-red unsafe) logarithm terms in $`\omega _a.`$ As a result, $`\widehat{\omega }_a`$ is free from both types of collinear “singularities” (in quotes since the logarithms become singular only in the zero-mass limit). In effect, all logarithm factors $`\mathrm{ln}(Q/m_c)`$ in $`\omega _a`$ are replaced by $`\mathrm{ln}(Q/\mu )`$ in $`\widehat{\omega }_a`$, (with accompanying finite subtractions), and the latter is *infra-red safe* in the $`\frac{m_c}{Q}0`$ limit.<sup>q</sup><sup>q</sup>qThe validity of these statements to order $`\alpha _s^2`$ can be inferred from the explicit calculations of Refs. . The proof to all orders of perturbation theory has been given in Ref. . Thus, the 4-flavor scheme has a well-defined high energy limit, and is expected to give a much more reliable description of the physics of charm production at large $`Q`$ than the 3-flavor scheme.
As formulated above, the hard cross-sections still contain finite charm-mass dependence, i.e. $`\widehat{\omega }_a=\widehat{\omega }_a(x,\frac{Q}{\mu },\frac{m_c}{Q},\mu )`$. Being infra-red safe, as $`m_c/Q0,`$ the limit $`\widehat{\omega }_a(x,\frac{Q}{\mu },\frac{m_c}{Q},\mu )\widehat{\omega }_a^{m_c=0}(x,\frac{Q}{\mu },\mu )`$ is well defined. In this limit, the 4-flavor scheme with non-zero charm mass reduces to the conventional zero-mass (ZM) 4-flavor parton scheme<sup>r</sup><sup>r</sup>rIn conventional zero-mass (ZM) 4-flavor theory, collinear singularities due to charm appear as $`ϵ^1`$ poles along with those from other flavors, and are regulated accordingly. When properly calculated, the massless limit of our ($`m_c0`$) Wilson coefficients, $`\widehat{\omega }_a^{m_c=0}(x,Q,\mu ),`$ should agree with the standard zero-mass results., as mentioned in the introduction. As emphasized in Ref. , however, the factorization of potentially dangerous $`\mathrm{ln}(m_c)`$ terms does not require taking the $`m_c0`$ limit in the infra-red safe coefficient functions. The conventional practice of always setting $`m=0`$ in the hard cross-section $`\widehat{\omega }_a(x,Q,\mu )`$ is a convenience, not a necessity; it results from the use of dimensional regularization of the zero-mass theory as a simple and efficient way to classify and to remove the collinear singularities. For a “heavy quark” with non-zero mass $`m_c,`$ this convenient method of achieving infra-red safety is not a natural one (as it is for light flavors), since $`m_c`$ itself already provides a natural cutoff. In other words, the theory has no real collinear “singularities” associated with the charm quark, and the universal (i.e. process-independent) and potentially large mass-logarithms can be factorized systematically as outlined above. In fact, by keeping the charm quark mass dependence, this scheme can be extended down to lower values of $`Q`$ with much more reliable results than in the zero-mass case. This is possible because of the well-defined relation between the 4-flavor calculation with non-zero $`m_c`$ and the 3-flavor (FFN) calculation; e.g. at order $`\alpha _s`$, Ref. showed that, for given $`x,Q`$
$${}_{}{}^{4}F_{2}^{c}(x,\frac{Q}{\mu },\frac{m_c}{\mu })\underset{lim\mu m_c}{}{}_{}{}^{3}F_{2}^{c}(x,\frac{Q}{\mu },\frac{m_c}{\mu })\text{ }𝒪(\alpha _s)$$
(18)
where the superscripts 3,4 refer to the 3- and 4-flavor scheme calculations respectively. To distinguish this more general 4-flavor scheme from the conventional zero-mass (ZM) 4-flavor scheme, we can refer it as the general-mass (GM) 4-flavor scheme.
The theoretical result (Eq. 18) does *not*, however, constrain the threshold behavior of the predicted physical structure function in the limit of $`limQ(\mathrm{or}W)m_c`$; to make a physical prediction, one needs to first choose $`\mu `$ as a function of the physical variables {$`x,Q,m_c`$}. This is related to the well-known scale dependence of PQCD prediction in general. We shall return to this problem at the end of the next subsection.
There is one additional advantage of the 4-flavor scheme. Since the charm quark distribution is explicitly included in the 4-flavor scheme, and since $`m_c`$ is not much larger than a typical non-perturbative scale such as the nucleon mass, one can allow for the existence of a possible nonperturbative (“intrinsic”) charm component inside a hadron at a low energy scale, say $`Q_0`$— as the boundary condition for evolution to higher scales, just like the other light flavors. This is a possibility not permitted in the 3-flavor scheme by assumption.
### A.4 The generalized $`\overline{\mathrm{MS}}`$ formalism with non-zero $`m_c`$
Both the 3-flavor and the 4-flavor schemes described above are valid schemes for defining the perturbative series of the inclusive cross section in principle. They are equivalent if both are carried out to all orders in the perturbation series. At a given finite order, they differ by a finite renormalization<sup>s</sup><sup>s</sup>sThe magnitude of the “finite” renormalization depends on the renormalization scale: e.g. $`\mathrm{ln}(m_c/\mu )`$ factors are finite, but can be numerically large if $`\mu m_c.`$ of the distribution functions, as well as the strong coupling $`\alpha _s`$. From the physics point of view, when calculated to the appropriate order (cf. below), the 3-flavor scheme provides a more natural and accurate description of the charm production mechanism near the threshold ($`Q^2m_c^2`$), whereas the 4-flavor scheme does the same in the high energy regime ($`Q^2m_c^2`$), as shown in Fig. 1.
The precise definitions given in the above subsections provide the means to implement the intuitive ideas discussed in the introduction. *A unified program* to calculate the inclusive structure functions, including charm, which maintains uniform accuracy over the full energy range, must be a composite scheme consisting of:
* (i) the 3-flavor scheme, applied from low energy scales, of the order of $`m_c`$, and extended up;
* (ii) the 4-flavor scheme, applied from high energy scales on down; and
* (iii) a set of matching conditions which define the perturbative relation between the two schemes applied at a specific matching scale $`\mu _m`$.
It is useful to explicitly discuss all the elements of this composite scheme which link the component 3-flavor and 4-flavor calculations discussed in previous subsections to physics predictions of the general formalism:
* Choice of scale: Within each scheme ($`i=3`$ or $`4`$), one needs to specify $`\mu `$ as a function of the physical variables in order to make a physical prediction, i.e.
$$F_{A,\lambda }^{phys}(Q^2,x,m_c)=F_{A,\lambda }^{(i)}(Q^2,x,m_c,\mu (x,Q,m_c))$$
Although there is considerable freedom in choosing $`\mu (x,Q,m_c)`$, two conditions should be met so that the prediction can be reliable: (i) $`\mu `$ must be of the order of $`Q`$ or $`m_c`$ so that PQCD applies, and (ii) $`F_{A,\lambda }^{(i)}(Q^2,x,m_c,\mu )`$ must be relatively stable with respect to variations of $`\mu `$ for the ($`x,Q`$) of interest. This is the well-known scale-dependence of any PQCD calculation. In Fig. 1, the presence of the uncertainty associated with the choice of $`\mu (x,Q,m_c)`$ in each scheme is represented by the respective bands.
* Matching conditions and choice of matching scale: For a given set of arguments, $`F_{A,\lambda }^{(3)}(Q^2,x,m_c,\mu )`$ and $`F_{A,\lambda }^{(4)}(Q^2,x,m_c,\mu )`$ are not independent. Being the same physical quantity calculated in two different schemes (cf. Sec. A.2 and A.3), they are related by a finite renormalization:
$$\begin{array}{c}{}_{}{}^{4}\alpha _{s}^{}(\mu )=^3\alpha _s(\mu )+\mathrm{\Delta }\alpha \\ {}_{}{}^{4}f_{}^{a}(x,\mu )=^3f^a(x,\mu )+\mathrm{\Delta }f^a(x,\mu )\end{array}\text{applied at }\mu =\mu _m$$
(19)
where $`\mathrm{\Delta }\alpha `$ and $`\mathrm{\Delta }f^a(x,\mu )`$ are fully calculable once the two schemes are defined. They have been calculated to order $`\alpha _s^2`$ . Specifically, the simpler results at order $`\alpha _s`$ are :
$${}_{}{}^{4}\alpha _{s}^{}(\mu )=^3\alpha _s(\mu )\left[1+\frac{{}_{}{}^{3}\alpha _{s}^{}(\mu )}{6\pi }\mathrm{ln}\frac{\mu ^2}{m_c^2}+𝒪(\alpha _s^2)\right]$$
(20)
$$\begin{array}{ccccccc}{}_{}{}^{4}f_{}^{q}(x,\mu )& =& {}_{}{}^{3}f_{}^{q}(x,\mu )& +& 0& +& 𝒪(\alpha _s^2)\\ {}_{}{}^{4}f_{}^{g}(x,\mu )& =& {}_{}{}^{3}f_{}^{g}(x,\mu )& & \frac{{}_{}{}^{3}\alpha _{s}^{}(\mu )}{6\pi }\mathrm{ln}\frac{\mu ^2}{m_c^2}^3f^g(x,\mu )& +& 𝒪(\alpha _s^2)\\ {}_{}{}^{4}f_{}^{c}(x,\mu )& =& 0& +& \frac{{}_{}{}^{3}\alpha _{s}^{}(\mu )}{4\pi }\mathrm{ln}\frac{\mu ^2}{m_c^2}\frac{dz}{z}(z^2+\left(1z\right)^2)^3f^g(\frac{x}{z},\mu )& +& 𝒪(\alpha _s^2)\end{array}$$
(21)
The scale at which these two schemes are matched will be called the matching point, and denoted by $`\mu _m.`$ Note that either scheme can still be used with $`\mu `$ above or below the matching point, it is just that the equations (19) are only enforced at $`\mu _m`$. In principle, $`\mu _m`$ can be chosen at any value – different choices lead to the same overall results, up to higher order corrections. As can be seen in Eqs. 20 and 21, in the generalized $`\overline{\mathrm{MS}}`$ scheme, $`\mathrm{\Delta }\alpha (\mu )`$ and $`\mathrm{\Delta }f^a(x,\mu )`$ are both of the form $`\alpha _s(\mu )\mathrm{ln}(\mu /m_c)C_1+𝒪(\alpha _s^2)`$. Thus, if one chooses $`\mu _m=m_c`$, both functions $`\alpha (\mu )`$ and $`f^a(x,\mu )`$ in the 3-flavor scheme are equal to their counterpart in the 4-flavor scheme to first order in $`\alpha _s`$ at the matching point. Most recent works adopt this choice ; we do the same in this paper. Although this choice is convenient, it is not required in the general formalism. The ideas behind the matching conditions, Eq. 19, are illustrated in Fig. 8 which show two possible matching points, the first being the special one $`\mu _m=m_c`$. This plot also shows that $`\mu _m`$ should not be chosen too far above $`m_c`$, lest the factor $`\alpha _s\mathrm{ln}(\mu /m_c)`$ in the discontinuity ceases to be perturbative.
* Choice of transition scale: With $`F_{A,\lambda }^{(i)}(Q^2,x,m_c,\mu (x,Q,m_c))`$, $`i=3,4`$, we have two sets of calculations, one for each scheme. For physics applications, we need to specify which of these to use, say
$$F_{A,\lambda }^{phys}(Q^2,x,m_c)=\{\begin{array}{ccc}F_{A,\lambda }^{(3)}(Q^2,x,m_c,\mu (x,Q,m_c))\hfill & & Q<\mu _t\hfill \\ F_{A,\lambda }^{(4)}(Q^2,x,m_c,\mu (x,Q,m_c))\hfill & & Q>\mu _t\hfill \end{array}$$
(22)
where we have introduced another scale $`\mu _t`$ – the transition point – where one switches from one scheme to the other, according to which one is more appropriate, as discussed in the introduction, cf. Fig.1 and the previous subsections of this appendix.
Conceptually, the transition point is distinct from the matching point, as should be clear from their defining equations, 19 and 22.<sup>t</sup><sup>t</sup>t This distinction was first made in Collins’ paper on the general formalism . The guiding principle for choosing the optimal $`\mu _t`$ is that it should be within the region where the $`F_{A,\lambda }^{(3)}`$ and $`F_{A,\lambda }^{(4)}`$ calculations are both valid, and that their differences within this region are small. For instance, in the idealized situation depicted in Fig. 1, $`\mu _t`$ is best chosen to be in the middle of the $`Q`$ range where both uncertainty bands are relatively narrow.
In practice, one can only estimate the range of uncertainties of the 3-flavor and 4-flavor calculations, say by examining the scale-dependence of the respective calculations and then make a judicious choice of $`\mu _t`$. In the case of inclusive charm production discussed in Sec. 3, Fig. 5 shows that the transition scale can be chosen at a relatively low value, close to $`m_c`$. For this choice, the composite scheme calculation reduces, in practice, to just the 4-flavor calculation.
### A.5 What do “LO” and “NLO” mean?
As already mentioned in the body of this paper, in a multi-scale problem such as heavy quark production, the designation of “LO” and “NLO” to a given calculation can be rather misleading in conventional fixed-order calculations, due to the presence of large logarithms which vitiates the naive counting of powers of $`\alpha _s`$. In the composite scheme, which has a wider range of applicability than FFN schemes, the meaning of “LO” and “NLO” can be better defined, provided the relative magnitudes of the large scales are properly kept in mind. We elaborate a little bit.
* In the 3-flavor scheme, “LO” consists of the $`𝒪(\alpha _s^1)`$ HC $`\gamma ^{}gc\overline{c}`$ process; whereas “NLO” involves $`𝒪(\alpha _s^2)`$ processes such as $`\gamma ^{}ggc\overline{c}`$. This formal designation makes physical sense only in the threshold region.
* In the 4-flavor scheme, the “LO” process is represented by the $`𝒪(\alpha _s^0)`$ HE $`\gamma ^{}cc`$ process; and the “NLO” ones consist of the $`𝒪(\alpha _s^1)`$ HC process as well as the $`𝒪(\alpha _s^1)`$ HE $`\gamma ^{}cgc`$ process. This formal designation coincides with the familiar one in the conventional treatment of DIS structure functions; it makes physical sense only when $`Q^2m_c^2`$.
The apparent mismatch of orders in $`\alpha _s`$ (e.g. order $`\alpha _s`$ being “LO” in the former, but “NLO” in the latter) can be understood within the composite scheme which takes into account the order of magnitudes of the other relevant quantities in the factorization formula at the appropriate energy ranges.
Specifically, the $`𝒪(\alpha _s^1)`$ and $`𝒪(\alpha _s^2)`$ 3-flavor calculations lose their “LO” and “NLO” meaning as $`Q^2`$ becomes very large compared to $`m_c^2,`$ since each power of $`\alpha _s`$ is accompanied (and neutralized) by a large logarithm $`\mathrm{ln}(Q^2/m_c^2)`$ factor in the hard cross-section. Consequently, all terms become effectively $`𝒪(\alpha _s^0)`$! In fact, the resummation of these large logarithms to all orders of $`\alpha _s`$ gives rise precisely to the charm distribution, resulting in the $`𝒪(\alpha _s^0)`$ charm excitation process $`\gamma ^{}cc`$ of the 4-flavor scheme. Conversely, in the 4-flavor scheme, although the charm quark distribution $`f_A^c(x,Q)`$ can be considered to be $`𝒪(\alpha _s^0)`$ at high $`Q^2`$ (where all flavors are on an equal footing), as we go down to a lower energy range $`Q^2m_c^2`$, one finds $`f_A^c(x,Q)\alpha _s^1\mathrm{ln}(Q^2/m_c^2)𝒪(\alpha _s^1)`$ compared to the dominant gluon and light quark distribution functions (assuming there is no large non-perturbative charm component). Therefore, to consistently match the 4-flavor calculation onto the LO ($`𝒪(\alpha _s^1)`$) 3-flavor calculation, one must include both the “LO” $`𝒪(\alpha _s^0)`$ $`\gamma ^{}cc`$ and the “NLO” $`𝒪(\alpha _s^1)`$ $`\gamma ^{}gc\overline{c}`$ contributions, along with the associated subtraction term in the 4-flavor calculation (Cf. Ref. ). This implies that our 4-flavor calculation, which is NLO at high energy scales, becomes effectively “LO” near the threshold, because both the formally $`𝒪(\alpha _s^0)`$ and $`𝒪(\alpha _s^1)`$ contributions become of the same order of magnitude, $`𝒪(\alpha _s^1)`$ – as in the LO 3-flavor scheme. (The calculations in the main part of this paper show that, with the natural choice of scale $`\mu =\sqrt{Q^2+m_c^2},`$ the numerical predictions of this calculation are actually fairly close to those of the NLO 3-flavor calculation, which is consistent with this observation because the size of the NLO correction is within the range of uncertainty of a LO result.)
This mixing of terms with different *apparent* powers of $`\alpha _s`$ is physically natural ( cf. Fig. 1 ) and logically consistent – it is a necessary feature of switching between different primary schemes, since any finite renormalization always entails a *resummation* (i.e. re-organization) of the perturbation series to all orders. In more concrete terms, the need for mixing terms of different apparent powers of $`\alpha _s`$ arises when:
* (i) the LO diagrams for different subprocesses start at different orders of $`\alpha _s;`$
* (ii) the associated parton densities are of different numerical orders of magnitude (such as between $`g,q,c`$);
* (iii) the order of magnitude of a parton distribution changes as it evolves with $`Q`$ (such as for $`c`$ in the region above the threshold); and
* (iv) the hard cross-section contains logarithms of ratios of energy scales which become large.<sup>u</sup><sup>u</sup>uFor these reasons, to require a naive uniform counting of powers of $`\alpha _s`$ over a wide range of $`Q,`$ when a composite scheme must be used, would miss the basic tenet of adapting the renormalization scheme to the appropriate number of active flavors as the physical scale varies.
The generalized $`\overline{\mathrm{MS}}`$ formalism, by keeping the physical $`m_c,`$ provides the appropriate scheme to describe the underlying physical processes in the different regions encountered in heavy quark production.
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# Hopf term induced by fermions
## Abstract
We derive an effective action for Dirac fermions coupled to O(3) non-linear $`\sigma `$-model (NL$`\sigma `$M) through the Yukawa-type interaction. The nonperturbative (global) quantum anomaly of this model results in a Hopf term for the effective NL$`\sigma `$M. We obtain this term using the “embedding” of the CP<sup>1</sup> model into the CP<sup>M</sup> generalization of the model which makes the quantum anomaly perturbative. This perturbative anomaly is calculated by means of a gradient expansion of a fermionic determinant and is given by the Chern-Simons term for an auxiliary gauge field.
It is well-known that nonperturbative anomalies in gauge theories can be reduced to perturbative ones by embedding the gauge group into a bigger one. In this paper we show that a similar method can be used to calculate the effect of nonperturbative anomalies on the effective action of a non-linear $`\sigma `$ model induced by fermions.
We consider an effective action $`S_{eff}(n)`$ of Dirac fermions on the three-dimensional sphere $`S^3`$ coupled to a background chiral field $`n`$:
$$e^{S_{eff}(n)}=𝑑\psi 𝑑\overline{\psi }\mathrm{exp}\left(_{S^3}d^3x\overline{\psi }\left[i\overline{)}+im\widehat{n}\right]\psi \right).$$
(1)
Here $`\overline{)}=\gamma ^\mu _\mu `$ where the $`\gamma ^\mu `$ are three-dimensional gamma-matrices which can be chosen, e.g., to be Pauli matrices, $`\widehat{n}=\stackrel{}{n}\stackrel{}{\tau }`$ with $`\stackrel{}{n}S^2`$, $`\stackrel{}{n}^2=1`$, $`\stackrel{}{\tau }`$ is a set of Pauli matrices acting in the isospace, and we use a Euclidian formulation.
We calculate the variation of the effective action $`S_{eff}(n)=\mathrm{ln}detD`$, $`D=i\overline{)}+im\widehat{n}`$ with respect to $`n`$.
$$\delta S_{eff}=\text{Tr}\left[\delta DD^{}(DD^{})^1\right],$$
(2)
where $`\delta D=im\delta \widehat{n}`$ and $`D^{}=i\overline{)}im\widehat{n}`$. Then we use $`DD^{}=^2+m^2+m\overline{)}\widehat{n}`$ and expand in $`m\overline{)}\widehat{n}`$ obtaining $`(DD^{})^1=G_0G_0m\overline{)}\widehat{n}G_0+G_0(m\overline{)}\widehat{n}G_0)^2+\mathrm{}`$, where $`G_0=\frac{1}{^2+m^2}`$.
Calculating traces and leaving only first nonzero orders in $`1/m`$ of real and imaginary parts of the effective action we obtain (only trace in isospace is left)
$`\delta S_{eff}`$ $`=`$ $`\delta S_{Re}+\delta S_{Im},`$ (3)
$`\delta S_{Re}`$ $`=`$ $`{\displaystyle \frac{|m|}{8\pi }}{\displaystyle d^3x\text{tr}(_\mu \delta \widehat{n})(_\mu \widehat{n})}+\mathrm{},`$ (4)
$`\delta S_{Im}`$ $`=`$ $`i{\displaystyle \frac{\text{sgn}(m)}{32\pi }}{\displaystyle d^3xϵ^{\mu \nu \lambda }\text{tr}(\widehat{n}\delta \widehat{n}_\mu \widehat{n}_\nu \widehat{n}_\lambda \widehat{n})}.`$ (5)
Now our goal is to restore the effective action from its variation (3-5). For the real part we have
$$S_{Re}=\frac{|m|}{16\pi }d^3x\text{tr}(_\mu \widehat{n})^2+\mathrm{}.$$
(6)
It is straightforward to check that the variation of the imaginary part of the action (5) is identically zero and naively $`S_{Im}=0`$. However, because of the nontrivial homotopy group $`\pi _3(S^2)=Z`$, the configurations of $`n`$ are divided into topological classes labeled by an integer-valued Hopf invariant $`H(n)`$. If $`S_{Im}H(n)`$ then $`\delta S_{Im}=0`$ and we can not find the “Hopf term” from our perturbative calculation. The presence of the Hopf term in the effective action is of great importance because it changes spin and statistics of solitons of NL$`\sigma `$M.
In the following we show that the correct result for the imaginary part is
$$S_{Im}=i\pi \text{sgn}(m)H(n).$$
(7)
The value of the coefficient in front of the Hopf invariant corresponds to the Fermi-Dirac statistics of solitons which agrees with their fermionic charge.
To find the imaginary part of the effective action we generalize the model (1), changing the size of isospace and replacing $`\widehat{n}`$ by $`\widehat{n}=2zz^{}1`$ with $`z^t=(z_1,z_2,\mathrm{},z_{M+1})`$ – complex vector with unit modulus $`z^{}z=1`$. The $`(M+1)\times (M+1)`$ matrix $`\widehat{n}`$ does not depend on the phase of $`z`$, i.e., $`z`$ should be considered as a CP<sup>M</sup> field. We note that for the particular configuration $`z^t=(z_1,z_2,0,\mathrm{},0)`$ the fermions with isospace indices higher than $`2`$ are decoupled from $`z`$ and do not contribute $`z`$-dependent terms into the effective action. As a consequence we can obtain the effective action for (1) by restricting the effective action of the CP<sup>M</sup> model to particular configurations $`z^t=(z_1,z_2,0,\mathrm{},0)`$.
It is easy to check that $`\widehat{n}^2=(2zz^{}1)^2=1`$ and perturbative results (3-5) are still valid. However, for $`M>1`$ the homotopy group $`\pi _3(\text{CP}^\mathrm{M})=0`$ and there are no topologically nontrivial configurations of $`\widehat{n}`$. The effective action can be found perturbatively! Substituting $`\widehat{n}=2zz^{}1`$ into (5) we obtain after some algebra
$$\delta S_{Im}=i\frac{\text{sgn}(m)}{2\pi }d^3xϵ^{\mu \nu \lambda }\delta a_\mu _\nu a_\lambda ,$$
(8)
where $`a_\mu =z^{}(i_\mu )z`$. From here we obtain
$$S_{Im}=i\frac{\text{sgn}(m)}{4\pi }d^3xϵ^{\mu \nu \lambda }a_\mu _\nu a_\lambda .$$
(9)
Restricting (9) to particular CP$`{}_{}{}^{1}=S^2`$ configurations we have (7) with the well-known expression for the Hopf invariant
$$H(n)=\frac{1}{4\pi ^2}d^3xϵ^{\mu \nu \lambda }a_\mu _\nu a_\lambda ,$$
(10)
where $`a_\mu =z^{}(i_\mu )z`$ with a two-component $`z`$ and $`\stackrel{}{n}=z^{}\stackrel{}{\tau }z`$.
Combining (6) and (7) we obtain for the effective action of the CP<sup>M</sup> non-linear $`\sigma `$-model induced by Dirac fermions
$$S_{eff}=d^3x\left\{\frac{|m|}{16\pi }\text{tr}(_\mu \widehat{n})^2+i\frac{\text{sgn}(m)}{4\pi }ϵ^{\mu \nu \lambda }a_\mu _\nu a_\lambda \right\},$$
(11)
where we kept only the imaginary part of the action and the terms of order $`m`$ in the real part. The second term of (11) is a perturbative “Chern-Simons” term in the case of $`M>1`$. It becomes a nonperturbative (global) Hopf term in the case of $`M=1`$.
In conclusion, we have derived the effective action for the CP<sup>M</sup> non-linear $`\sigma `$-model induced by Dirac fermions on a three-dimensional sphere. We have shown that this effective action has a nontrivial topological term which (for $`M=1`$) is equal to the well-known Hopf term for the O(3) non-linear $`\sigma `$-model. We used the method of embedding known for global anomalies in gauge theories.
We would like to point out some differences between global terms for the $`\sigma `$-model and for the gauge field models. In the case of gauge fields there always exists a direct interpolation between configurations which differ by non-trivial gauge transformation. Therefore, the question of the relative phase of fermionic determinants for those configurations is well-defined. In the case of the NL$`\sigma `$M there is no such direct interpolation between configurations from different topological classes. Therefore, strictly speaking, these topological classes can be weighed in the partition function corresponding to (1) with arbitrary weights. The imaginary part of an effective action, or the relative phase of determinants with chiral field configurations belonging to different topological classes, is not defined. However, in a realistic physical model there might be some regularization which allows one to connect field configurations from different topological classes. E.g., if a theory is defined on a lattice, the singular processes changing topological classes are allowed. Another possibility is to make a constraint, $`\widehat{n}^2=1`$, soft. Then at some points in space-time $`\widehat{n}=0`$, the target manifold is not $`S^2`$, and there are no distinct topological classes anymore.
In this sense the “embedding method” we used is not just a technical trick but is essentially a method of regularization which allows us to interpolate between different topological classes of NL$`\sigma `$M.
We thank M. Braverman for discussing mathematics. We would also like to thank P.A. Lee, X.-G. Wen, and especially P.B. Wiegmann for many helpful discussions. We appreciate the hospitality of Aspen Center for Physics where part of this work has been done. This research has been supported by NSF DMR 9813764.
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# Critical properties of projected 𝑆𝑂(5) models at finite temperatures
## I Introduction
The $`SO(5)`$ theory of high-T<sub>c</sub> superconductivity has been introduced as a concept to unify antiferromagnetism (AF) and d-wave superconductivity (SC) under a common symmetry principle. In order to study the physical consequences, and to make predictions to compare with experiments, several exact $`SO(5)`$ -symmetric models with small symmetry-breaking terms have been proposed and investigated in detail . However, a shortcoming of these models, and of an exact $`SO(5)`$ theory in general is that they are inconsistent with the antiferromagnetic gap at half filling, one of the most important features of the high-T<sub>c</sub> cuprates . This can be understood by the fact that an $`SO(5)`$ transformation “rotates” spin into charge and, thus, a requirement for an exactly $`SO(5)`$ -invariant system would be to have the same charge and spin gap. This is in contradiction with the experimental situation in the high-T<sub>c</sub> materials, where a large charge gap of some eV is present in the AF state at half-filling, while spin-wave excitations are ungapped. The introduction of a small symmetry-breaking term , while on the one hand correctly selecting the AF state at half-filling and shifting the AF-SC transition to finite doping, does not introduce a charge gap of the correct order of magnitude. In contrast, in a weakly-coupled Hubbard ladder model, a spin and a charge gap of the same size are present and, in fact, it has been shown that $`SO(5)`$ symmetry is dynamically restored at half filling .
In order to cure this problem at strong coupling as well, a class of $`SO(5)`$ models – “projected” $`SO(5)`$ models – has been introduced where the Mott-Hubbard gap is taken into account by means of a Gutzwiller projection, whereby doubly-occupied states are projected out . In that , it was shown that, despite the symmetry-breaking effects of the projection, static correlation functions remain exactly $`SO(5)`$ symmetric within a mean-field approximation. This is due to the fact that, neglecting dynamic effects, the Hamiltonian is manifestly $`SO(5)`$ invariant. However, dynamic effects breaking the $`SO(5)`$ symmetry become important whenever quantum fluctuations are taken into account . In another paper , it was shown that the projection is crucial in order to correctly relate the $`d`$-wave superconducting gap at finite doping with the $`d`$-wave modulation of the AF gap observed at half filling by ARPES experiments .
In a microscopic physical system, one would in general expect $`SO(5)`$ symmetry to be explicitly broken by several terms. This is certainly the case for the Hubbard model, for example. However, it often occurs in nature that a symmetry, which is broken on the microscopic level, is then restored in the long-wavelength limit. Concerning $`SO(5)`$ , this has been shown to happen for quite generic ladder systems of the Hubbard type . Recently, Murakami and Nagaosa argued that the bicritical point of the AF to SC transition in the organic superconductor $`\kappa `$-(BEDT-TTF)<sub>2</sub>X (Bis(ethylenedithio)tetrathiafulvalene) with X=Cu\[N(CN)<sub>2</sub>\]Cl, shows $`SO(5)`$ critical exponents . In fact, one of the scenarios suggested by Zhang is that there might be a direct first-order AF to SC transition terminating at a finite-temperature bicritical point, where the $`SO(5)`$ symmetry is asymptotically restored at long wavelengths . These ideas are very interesting from an experimental point of view, and open the possibility of an explicit test of $`SO(5)`$ symmetry, via a direct “measurement of the number 5”. This could be done, as suggested in Ref. , by measuring the critical exponents of the AF to SC transition, which, given the spatial dimensionality, should only depend on the number of components $`n`$ of the order parameter. On the other hand, it is well known that for $`n>n_c4`$, the $`SO(5)`$ symmetric fixed point is unstable towards a so-called biconical fixed point . However, since $`n=5`$ is close to $`n_c`$, it turns out that the stable biconical fixed point only breaks the symmetry by about $`20\%`$.
However, the situation of the high-T<sub>c</sub> materials is quite delicate. As discussed above, the Mott-Hubbard gap plays an important role, and it produces a substantial breaking of the $`SO(5)`$ symmetry. In Ref. , it was shown that in the extreme case of a Gutzwiller projection, a degree of freedom is eliminated completely and the real and imaginary part of the local superconducting parameter become conjugate variables. Therefore, it is not clear whether such a projected $`SO(5)`$ symmetry can become asymptotically a complete $`SO(5)`$ symmetry in the neighborhood of some critical point.
In this paper, we show why symmetry-breaking effects due to the projection are asymptotically irrelevant in the neighborhood of a finite-temperature critical point. However, two kinds of symmetry-breaking effects tend to prevent $`SO(5)`$ symmetry from being restored asymptotically. One is related with the different mobilities of hole pairs and magnons ($`\eta 1`$ below), and the second one is due to the renormalization effects from quantum fluctuation at an intermediate length scale. The common tendency of these effects is to draw the system into a region of instability, where the two AF/normal(N) and SC/N transitions become first order before merging at the AF/SC/N triple point . However, when the first effect is large, quantum fluctuations tend to take the system back to the $`SO(5)`$ point.
This paper is organized as follows: In Sec. II we start from the projected $`SO(5)`$ -symmetric model (allowing for a symmetry-breaking term $`\eta `$ in the mobilities), and treat it by a slave-boson functional-integral approach in order to deal with the hard-core constraint. The important result is that at the AF-SC transition, the classical part of the action of the projected model preserves its $`SO(5)`$ structure at $`\eta =1`$, despite the symmetry-breaking terms arising from the projection. These terms only appear in the quantum-mechanical (i. e., time derivative) part of the action. This fact gives a rigorous justification for the much used semiclassical description of the high-T<sub>c</sub> materials via a $`SO(5)`$ -symmetric model despite of the presence of the large Hubbard gap. In Sec. III, we derive the associate effective Ginzburg-Landau model by integrating out the momenta conjugate to the AF superspin variables . We study the properties of such model in the neighborhood of the AF/SC/N triple point and discuss the possibility of $`SO(5)`$ symmetry restoring at long wavelengths.
In Sec. IV, we evaluate the corrections to the effective classical action due to the so far neglected quantum fluctuations. For small temperatures, these mainly affect the magnon-magnon scattering, thus breaking the $`SO(5)`$ symmetry.
Finally, in Sec. V we draw our conclusions. Some details of the calculations are given in the appendices.
## II The model
We start from the effective bosonic model introduced in Refs. , which describes low-energy bosonic excitations of “blocks”, (also referred to as “sites”) labeled by the coordinate $`x`$, consisting of a rung in a 1-D ladder or of a $`2\times 2`$ plaquette in a 2-D system.
$`H`$ $`=`$ $`\overline{\mathrm{\Delta }}_s{\displaystyle \underset{x}{}}t_\alpha ^{}(x)t_\alpha (x)+\overline{\mathrm{\Delta }}_c{\displaystyle \underset{x}{}}t_i^{}(x)t_i(x)`$ (1)
$``$ $`\overline{J}_s{\displaystyle \underset{<xx^{}>}{}}n_\alpha (x)n_\alpha (x^{})\overline{J}_c{\displaystyle \underset{<xx^{}>}{}}n_i(x)n_i(x^{})`$ (2)
In this paper, we shall use similar conventions as in Ref. , where the indices $`a,b,..`$ are the $`SO(5)`$ superspin indices and take the values $`1,2,3,4,5`$ (in some cases, they might also include the “hole” index $`h`$), $`\alpha ,\beta ,..=2,3,4`$ (corresponding to $`x,y,z`$) denote the spin indices and $`i,j=1,5`$ denote the charge indices, and repeated indices are implicitly summed over. A boldface sign indicates the vector as a whole. Here, $`\overline{\mathrm{\Delta }}_s`$ is the energy required to produce a magnon excitation, i. e. to replace a singlet with a triplet in a block, while $`\overline{\mathrm{\Delta }}_c`$ is the energy required in order to produce a particle or hole pair. It is clear that $`\overline{\mathrm{\Delta }}_c`$ is of the order of the Mott-Hubbard gap and, thus, $`\overline{\mathrm{\Delta }}_c\overline{\mathrm{\Delta }}_s`$. $`\overline{J}_s`$ and $`\overline{J}_c`$ describe the hybridization of these excitations between nearest-neighbor sites and are related to their mobility. The Hamiltonian Eq. (2) acts on a “vacuum” $`|\mathrm{\Omega }`$, which is a kind of “RVB” state consisting of a product state of half-filled singlet states $`|\mathrm{\Omega }(x)`$ in each block. On the other hand, the five-fold states $`t_a^{}(x)|\mathrm{\Omega }(x)`$ describe the triplet magnon states (for $`a=2,3,4`$), and the $`d`$-wave hole and particle pair states on a block ($`a=1,5`$. More specifically, one can define the charge eigenoperators $`t_h`$ and $`t_p`$ as
$$t_1=\frac{1}{\sqrt{2}}(t_h+t_p)t_5=\frac{1}{i\sqrt{2}}(t_ht_p),$$
(3)
where $`t_h^{}`$ is the creation operator for a hole pair and $`t_p^{}`$ is the creation operator for a particle pair. In Eq. (2), the $`n_a`$ play the role of the “displacement” coordinates of the local harmonic oscillators, while we will denote with $`p_a`$ the conjugate momenta, and we have the transformation to canonical variables:
$$t_a=\frac{1}{\sqrt{2}}\left(n_a+ip_a\right).$$
(4)
Due to their microscopic origin these bosonic states are hard-core bosons, in the sense that at most one boson can reside on each site. Mathematically, this is expressed by the condition
$$t_a^{}(x)t_a(x)1.$$
(5)
The particle and hole density is controlled by the chemical potential $`\mu `$, which couples to the Hamiltonian via a term
$$H_\mu =2\mu \underset{x}{}\left[t_p^{}(x)t_p(x)t_h^{}(x)t_h(x)\right].$$
(6)
In the presence of this chemical potential term, the gap energy of the hole and particle pairs become $`\overline{\mathrm{\Delta }}_c+2\mu `$ and $`\overline{\mathrm{\Delta }}_c2\mu `$ respectively. A (negative) chemical potential of the order of the charge gap $`\overline{\mathrm{\Delta }}_c/2`$ is needed to induce an AF-SC transition in this system. Near such a transition point, the gap energy of the hole pair $`\overline{\mathrm{\Delta }}_c+2\mu `$ can be comparable to the (local) spin gap $`\overline{\mathrm{\Delta }}_s`$, while the gap towards a particle pair excitation is pushed up and becomes of the order of twice the charge gap. Since this is a very large energy scale, we can safely project this excitation out of the spectrum in the low–energy limit, by requiring that the condition
$$t_p(x)|\mathrm{\Psi }=0$$
(7)
is fulfilled at every site $`x`$. The new Hamiltonian takes the form
$`H=\overline{\mathrm{\Delta }}_s{\displaystyle \underset{x,\alpha }{}}t_\alpha ^{}(x)t_\alpha (x)+(\overline{\mathrm{\Delta }}_c+2\mu ){\displaystyle \underset{x}{}}t_h^{}(x)t_h(x)`$ (8)
$`\overline{J}_s{\displaystyle \underset{<x,x^{}>,\alpha }{}}n_\alpha (x)n_\alpha (x^{})`$ (9)
$`\overline{J}_c/2{\displaystyle \underset{<x,x^{}>}{}}(t_h^{}(x)t_h(x^{})+h.c.).`$ (10)
In Ref. it was shown that the constraint Eq. (7) can be enforced by introducing canonical commutation rules between the two variables $`n_1`$ and $`n_5`$, i. e.
$$[n_1,n_5]=i/2,$$
(11)
and therefore we can identify $`\sqrt{2}n_1`$ with the “hole displacement” $`n_h`$ and $`\sqrt{2}n_5`$ with its conjugate momentum $`p_h`$. The $`SO(5)`$ structure of the Hamiltonian becomes now clear if one introduces the superspin vector
$$m_a(\eta n_h,n_2,n_3,n_4,\eta p_h),$$
(12)
where, for convenience, we have absorbed the different mobility for hole pairs and magnons $`\eta \sqrt{\frac{\overline{J}_c}{2\overline{J}_s}}`$ into the definition of the superspin. Carrying out the transformation to canonical variables Eq. (4), the Hamiltonian Eq. (8) now takes the simple form
$`H`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }_s}{2}}{\displaystyle \underset{x}{}}p_\alpha (x)^2+{\displaystyle \frac{\mathrm{\Delta }_s}{2}}{\displaystyle \underset{x}{}}m_\alpha (x)^2+{\displaystyle \frac{\mathrm{\Delta }_c}{2}}{\displaystyle \underset{x}{}}m_i(x)^2`$ (13)
$``$ $`J{\displaystyle \underset{<xx^{}>}{}}m_a(x)m_a(x^{}),`$ (14)
where we have further redefined
$$\mathrm{\Delta }_c\frac{\overline{\mathrm{\Delta }}_c+2\mu }{\eta ^2}\mathrm{\Delta }_s\overline{\mathrm{\Delta }}_s\text{ , and}J\overline{J}_s.$$
(15)
The anisotropy in superspin space due to $`\eta `$ reflects now into the constraint, as we will see in Eq. (21) below. If one forgets for a moment the connection between coordinates $`m_a`$ and their conjugate momenta, the Hamiltonian Eq. (13) becomes exactly $`SO(5)`$ invariant under rotation of the superspin Eq. (12) at the AF-SC transition point $`\mathrm{\Delta }_s=\mathrm{\Delta }_c`$, which is reached by changing the chemical potential $`\mu `$, i. e. at the AF-SC transition. If one further has $`\eta =1`$, i. e. $`2\overline{J}_s=\overline{J}_c`$, the constraint is invariant as well, and one apparently has a complete $`SO(5)`$ symmetric model (cf. Ref. ). More specifically, one would like to $`SO(5)`$ “rotate” just the $`m_a`$ coordinates, leaving the conjugate coordinates to the magnon part $`p_\alpha `$ unrotated. This is possible, for example, in a classical ensemble, where, due to Liouville’s theorem, expectation values are evaluated as $`_idp_idq_i\mathrm{exp}H[p_i,q_i]`$, and rotations of the $`q_i`$ only leaves the measure invariant. Of course, this does not hold for dynamics, which is affected by the relation between the two “superconducting” canonically conjugate components $`m_1/\eta `$ and $`m_5/\eta `$, and between the AF components $`m_\alpha `$ and their conjugate momenta $`p_\alpha `$. Thus, $`SO(5)`$ symmetry is broken in dynamics, as pointed out in Ref. . Unfortunately, the relation between conjugate variables is also important in quantum-mechanical static averages, so that ground-state or finite-temperature averages are generally expected to break the symmetry when the full quantum problem is taken into account.
In order to understand the nature of the symmetry-breaking terms, it is convenient to go over to a functional-integral representation of the partition function for the Hamiltonian Eq. (8). The hard-core constraints can be conveniently taken care of by means of a slave-boson representation , where the boson operator $`e(x)`$ labeling “empty” sites is introduced. The detailed procedure is shown in Appendix A. After this transformation, the action takes the form
$$S=S_{QM}+S_{CL},$$
(16)
where
$$S_{QM}=_0^\beta 𝑑\tau \underset{x}{}\left[ip_\alpha (x,\tau )\dot{m}_\alpha (x,\tau )\frac{i}{\eta ^2}m_5(x,\tau )\dot{m}_1(x,\tau )\right],$$
(17)
($`\dot{m}_a`$ indicates the time derivative of $`m_a`$) has the well-known form $`p\dot{q}`$ of the Feynman path integral, the $`i`$ coming from the imaginary-time representation. Moreover,
$`S_{CL}={\displaystyle _0^\beta }d\tau \{{\displaystyle \frac{\mathrm{\Delta }_s}{2}}{\displaystyle \underset{x}{}}p_\alpha (x,\tau )^2+{\displaystyle \frac{\mathrm{\Delta }_s}{2}}{\displaystyle \underset{x}{}}m_\alpha (x,\tau )^2+{\displaystyle \frac{\mathrm{\Delta }_c}{2}}{\displaystyle \underset{x}{}}m_i(x,\tau )^2`$ (18)
$``$ $`J{\displaystyle \underset{<xx^{}>}{}}e(x,\tau )m_a(x,\tau )e(x^{},\tau )m_a(x^{},\tau )\},`$ (19)
where we have to replace
$$e(x,\tau )=\sqrt{1\frac{p_\alpha (x,\tau )^2}{2}\frac{m_\alpha (x,\tau )^2}{2}\frac{m_i(x,\tau )^2}{2\eta ^2}},$$
(20)
which implicitly includes the condition
$$\frac{p_\alpha (x,\tau )^2}{2}+\frac{m_\alpha (x,\tau )^2}{2}+\frac{m_i(x,\tau )^2}{2\eta ^2}1,$$
(21)
and where we have already carried out the transformation to canonical coordinates, Eq. (4), for the corresponding fields . Eq. (18) is the correct classical limit of a projected, i. e. of the physical $`SO(5)`$ model. Notice that the effects of the hard-core constraint is to introduce a renormalization of the boson hopping, and to bound the superspin magnitude, without, however, fixing its length . Thus, the requirement that the superspin magnitude be unity should not be taken as a rigorous constraint of the $`SO(5)`$ theory, at least not of the projected one (which is the physical one). On the other hand, one expects that in the homogeneous ordered phase this constraint might be a good assumption. A similar result has been shown by Wegner , namely, that the orthogonality constraint in the exact $`SO(5)`$ model is not a rigorous constraint, but it is favored at high temperature, as it maximizes the entropy.
Eq. (16) clearly identifies the $`SO(5)`$ -symmetry breaking terms. The classical action $`S_{CL}`$ is exactly $`SO(5)`$ invariant at the AF-SC transition ($`\mathrm{\Delta }_s=\mathrm{\Delta }_c`$) and for $`\eta =1`$, while apparently incurable symmetry-breaking terms come from the time-derivative terms in $`S_{QM}`$. More specifically, with these values of the parameters, if one carries out an $`SO(5)`$ rotation within the superspin vector, Eq. (12), $`S_{CL}`$ remains invariant, while $`S_{QM}`$ is changed. If quantum fluctuations are neglected, one can choose time-independent fields and set Eq. (17) to zero. In this case, any equilibrium expectation value is exactly $`SO(5)`$ invariant. More specifically, let us take a generic $`SO(5)`$ rotation matrix $`R(𝐧)=\mathrm{exp}in_a\mathrm{\Gamma }_a`$ parametrized by the vector $`𝐧`$ ($`\mathrm{\Gamma }_a`$ are the $`SO(5)`$ generators ), and $`f[𝐦(x),𝐩(x)]`$ is a function of the superspin vector $`𝐦(x)`$, and, possibly, of $`𝐩(x)`$. Then, the classical expectation values $`<>_{CL}`$ have the property
$$<f[𝐦(x),𝐩(x)]>_{CL}=<f[𝐑𝐦(x),𝐩(x)]>_{CL},$$
(22)
which is the requirement of $`SO(5)`$ invariance. Notice that the $`p_\alpha `$ should not be rotated, while in an exact $`SO(5)`$ model they should.
The question is: when is it justified to neglect the time dependence of the fields ?. This is allowed at moderately high temperatures, more precisely, at temperatures much larger than $`v/\xi `$ (in units of $`k_B=\mathrm{}=1`$), where $`\xi `$ is the correlation length and $`v`$ is a typical velocity, in our case equal to $`Ja`$, $`a`$ being the lattice spacing. This means that neglecting $`S_{QM}`$ is exactly justified when $`\xi `$ becomes infinite, i. e. in the neighborhood of a finite-temperature critical point, as a possible (finite temperature) multicritical point at which the AF/N and the SC/N transition lines merge into a first-order line . Moreover, this critical point is indeed a good candidate for a possible asymptotic restoring of the complete $`SO(5)`$ symmetry even in the presence, microscopically, of a projected $`SO(5)`$ symmetry. This is very important as it would mean that the large-energy symmetry-breaking effect of the Mott-insulator gap would be exactly compensated at this critical point. This is analogous to the well-known situation for the antiferromagnetic spin-flop transition where a system with uniaxial anisotropy restores $`SO(3)`$ symmetry at the bicritical point. However, there are some important differences with respect to the spin-flop transition, as we will show in the next Sections. Moreover, notice that due to the symmetry breaking term, Eq. (17), it is unlikely that $`SO(5)`$ symmetry can be restored if the AF-SC transition is controlled by a quantum-critical point. Since we are interested in finite-temperature critical points, we will restrict to the case of three spatial dimensions $`D`$.
## III Effective Ginzburg-Landau action
In this Section, we study the action Eq. (13) in more detail. We first integrate out the momenta $`p_\alpha `$ and obtain an effective action restricted to the superspin variables. For temperatures smaller than the singlet-triplet splitting $`\mathrm{\Delta }_s`$, one can restrict to a Gaussian integration of the momenta, i. e., consider only quadratic terms in $`p_\alpha `$. Carrying out such an expansion, one obtains
$$S_{CL}=S_{pm}+S_m+𝒪(p_\alpha ^4),$$
(23)
where, leaving the $`\tau `$ dependence implicit
$`S_m={\displaystyle _0^\tau }d\tau [{\displaystyle \frac{\mathrm{\Delta }_s}{2}}{\displaystyle \underset{x}{}}m_\alpha (x)^2+{\displaystyle \frac{\mathrm{\Delta }_c}{2}}{\displaystyle \underset{x}{}}m_i(x)^2`$ (24)
$``$ $`J{\displaystyle \underset{<xx^{}>}{}}r(x)m_a(x)r(x^{})m_a(x^{})],`$ (25)
$$S_{pm}=_0^\tau 𝑑\tau \underset{x}{}\frac{\mathrm{\Delta }_s}{2}𝒜(x)p_\beta (x)^2,$$
(26)
where we have defined
$$𝒜(x)1+\frac{2J}{\mathrm{\Delta }_s}\frac{m_a(x)}{4r(x)}\underset{d}{\overset{nn}{}}m_a(x+d)r(x+d),$$
(27)
$$r(x)\sqrt{1\frac{m_\alpha (x)^2}{2}\frac{m_i(x)^2}{2\eta ^2}},$$
(28)
and the sum $`_d^{nn}`$ extends over nearest-neighbor sites.
It is now convenient to reabsorb the $`𝐦`$-dependent coefficient $`𝒜(x)`$ of the $`𝐩^2`$ term into the definition of the momenta $`𝐩`$. This is done in order to avoid the appearance of terms depending on the amplitude of the imaginary-time slice in the effective action. Furthermore, in order to avoid a $`𝐦`$-dependent Jacobian due to the transformation, it is convenient to transform the $`𝐦`$-coordinates in such a way that the Jacobian remains unity. The general procedure is illustrated in Appendix B. Up to second order in $`𝐦^2`$, the new $`𝐦^{}`$ coordinates are related with the old ones via
$$m_a(x)=m_a^{}(x)(1+\frac{3J}{7\mathrm{\Delta }_s}|𝐦^{}(x)|^2).$$
(29)
After this transformation, the integration of the $`p_\alpha ^{}(x)p_\alpha (x)\sqrt{𝒜(x)}`$ only affects $`S_{QM}`$, and one obtains a new QM action in the form
$$S_{QM}^{}=_0^\beta 𝑑\tau 𝑑x\left[\frac{\dot{m}_\alpha ^2}{2\mathrm{\Delta }_s𝒜(x)}\frac{i}{\eta ^2}m_5(x)\dot{m}_1(x)\right],$$
(30)
where the transformation Eq. (29) should be inserted, and we have absorbed the unit cell volume $`𝒱=a^3`$ in the definition of the fields by renaming $`m_a^2/𝒱m_a^2`$.
Thus, the total effective $`SO(5)`$ action restricted to the superspin variables is given by Eq. (24) plus Eq. (30). The transformation Eq. (29) must still be carried out on the $`m`$ variables, but, due to the fact that the coefficient $`\frac{3J}{7\mathrm{\Delta }_s}`$ is small at the transition, this does not change the result significantly. On the other hand, it is important to take into account the effects of the hard-core constraint, which introduces the transformation Eq. (28), and, implicitly, the restriction of the superspin within a $`5`$-dimensional hypersphere (or a ellipsoid, if $`\eta 1`$) . Thus, $`S_m`$ ( Eq. (24)) gives the first effective classical functional microscopically derived from an $`SO(5)`$ model , where the physics of the Mott insulating gap has been properly taken into account via the projection. This is the appropriate functional which should be used for physical predictions of the $`SO(5)`$ theory, consistent with the gap.
Close to the phase transitions, it is more convenient to derive a Ginzburg-Landau form for the action, obtained, as usual, by expanding in powers of the field $`𝐦`$ and keeping only lowest-order gradient terms. After inserting Eq. (29) and dropping the prime indices in the fields $`𝐦`$, we obtain
$`S_{CL}^{}={\displaystyle _0^\beta }d\tau {\displaystyle }dx\{{\displaystyle \frac{r_s}{2}}m_\alpha (x)^2+{\displaystyle \frac{r_c}{2}}m_i(x)^2+{\displaystyle \frac{\rho }{2}}(\stackrel{}{}m_a(x))^2`$ (31)
$`+{\displaystyle \frac{u_s}{8}}\left({\displaystyle \underset{\alpha }{}}m_\alpha (x)^2\right)^2+{\displaystyle \frac{u_c}{8}}\left({\displaystyle \underset{i}{}}m_i(x)^2\right)^2+{\displaystyle \frac{u_{cs}}{4}}\left({\displaystyle \underset{\alpha }{}}m_\alpha (x)^2\right)\left({\displaystyle \underset{i}{}}m_i(x)^2\right)\},`$ (32)
with the parameters
$`{\displaystyle \frac{r_{s/c}}{2}}=({\displaystyle \frac{\mathrm{\Delta }_{s/c}}{2}}DJ)`$ (33)
$`{\displaystyle \frac{\rho }{2}}={\displaystyle \frac{Ja^2}{2}}`$ (34)
$`{\displaystyle \frac{u_s}{8}}={\displaystyle \frac{𝒱J}{2}}(D+{\displaystyle \frac{3r_s}{7\mathrm{\Delta }_s}})`$ (35)
$`{\displaystyle \frac{u_c}{8}}={\displaystyle \frac{𝒱J}{2}}({\displaystyle \frac{D}{\eta ^2}}+{\displaystyle \frac{3r_c}{7\mathrm{\Delta }_s}})`$ (36)
$`u_{cs}={\displaystyle \frac{u_c+u_s}{2}}.`$ (37)
The critical properties of the model Eq. (31) have been analyzed in several works . Its phase diagram is determined by two relevant parameters, the first one $`r_sr_c\mathrm{\Delta }_s\mathrm{\Delta }_c`$ controls the transition between the AF and the SC phases, while the other $`\mathrm{min}(r_s,r_c)`$ controls the second-order transition between the appropriate ordered (AF or SC) and the normal phase. At the transition point $`r_sr_c0`$, there are two competing fixed point controlling the transition , the Heisenberg bicritical fixed point (in this specific case, the $`SO(5)`$ fixed point), and the biconical tetracritical fixed point. According to the $`ϵ`$-expansion, the latter fixed point turns out to be the stable one for $`n>n_c4O(ϵ)`$. This means that, in general, the model Eq. (31), which has $`n=5`$, is expected to flow to this latter fixed point and not to the $`SO(5)`$ -symmetric one for $`u_su_cu_{cs}`$. On the other hand, since $`n=5`$ is not very far away from $`n_c`$, the stable biconical fixed point is approximately $`SO(5)`$ invariant with symmetry-breaking terms of the order of $`20\%`$. Moreover, there is a plane in the $`u_s,u_c,u_{cs}`$ space, given by the condition $`u_{cs}^2=u_cu_s`$, from which the system flows to the $`SO(5)`$ point . This is due to the fact that a scale transformation of, say, the SC components $`m_i^2m_i^2u_s/u_c`$ of the order parameter would yield again an $`SO(5)`$ -symmetric interaction of the form $`u|𝐦|^4`$. The asymmetry would then be transfered into different susceptibilities $`\rho _s,\rho _c`$ for the AF and for the SC order parameters. However, it has been shown in Refs. that the different in the susceptibilities is an irrelevant parameter.
In our case, we have
$$\mathrm{\Delta }u^2u_{cs}^2u_cu_s=\left(\frac{u_cu_s}{2}\right)^20$$
(38)
which means that the $`SO(5)`$ symmetric fixed point is never reached, except when the equal sign holds, i. e., when $`\eta 1`$ (at the transition $`r_s=r_c`$). On the other hand, we expect on physical grounds the mobility of the hole pairs to be smaller than that of the magnons, and, thus, $`\eta `$ to be smaller than $`1`$. Unfortunately, for the case Eq. (38), the couplings flow away into a region of instability. The common interpretation is that the AF/N and SC/N transitions become first order as well (fluctuation-induced first-order transition), at least close enough to the AF/SC/N triple point .
This fact seems in contrast with the apparent observation of bicritical behavior with $`SO(5)`$ critical exponents in the organic superconductor $`\kappa `$-(BEDT-TTF)<sub>2</sub>X (see Refs. ), by Murakami and Nagaosa . There may be several ways to understand this. One possibility is that other effects not considered here, such as, e. g., Coulomb interactions, fermionic excitations , or quantum effects, as discussed in Sec. IV, counterbalance this effect and draw the system back to the domain of attraction of the biconical fixed point. As discussed above, the differences between the $`SO(5)`$ and the biconical fixed point are only about $`20\%`$, so that they might be not observable experimentally. Alternatively, since the flow would cross the $`SO(5)`$ plane, it could produce $`SO(5)`$ exponents at intermediate length scales. On the other hand, Hu et al. , observe a coexistence region of AF and SC for the $`SO(5)`$ -anisotropic case, which could be possibly identified with the biconical phase. Their result could be due to the fact that they consider a different $`c`$-axis anisotropy ($`\chi `$ in footnote ), for the AF and for the SC variables.
## IV Quantum corrections
Even when considering a classical (i. e., finite-temperature) critical point, the quantum-mechanical symmetry-breaking terms $`S_{QM}`$ although irrelevant in the RG sense, contribute to the RG flow up to a certain length scale of the order of $`v/T`$. Since $`S_{QM}`$ breaks the $`SO(5)`$ symmetry, it is expected, during this initial renormalization process, to introduce symmetry-breaking terms in $`S_{CL}`$. Therefore, even when $`S_{CL}`$ is $`SO(5)`$ symmetric at the microscopic scale, the renormalized $`S_{CL}`$ at the scale $`\xi v/T`$ will probably break the symmetry. In this Section, we evaluate these symmetry-breaking terms originating from $`S_{QM}`$, or, more precisely, from the time dependence of the fields.
In order to evaluate these effects, we separate the fields into their static and dynamic parts, and integrate out the latter. Since we are working at finite temperature, we have to integrate out the components of the fields with Matsubara frequencies $`\omega _n=2\pi nT`$ with $`n0`$. In order to obtain an analytic expression for these corrections, we restrict to one-loop contributions and take just the leading low-temperature terms.
We first diagonalize the non-interacting (quadratic) part of the action Eq. (30) plus Eq. (31) by Fourier transform. We can neglect the corrections to $`S_{QM}^{}`$ due to the transformation to the primed variables Eq. (29), as it introduces irrelevant quartic time-derivative terms. In Fourier space, the action takes the usual form
$$S_{QM}^{}+S_{CL}^{}=\frac{1}{2}_km_a(k)\left[G(k)^1\right]_{ab}m_b(k)+\frac{1}{8}_{k_1,k_2,k_3}m_a(k_1)m_a(k_2)u_{ab}m_b(k_3)m_b(k_1k_2,k_3),$$
(39)
where we have introduced the shorthand notation $`k(k,\omega )`$, and $`_k\frac{1}{\beta }_\omega ^\mathrm{\Lambda }\frac{d^3k}{(2\pi )^3}`$, with $`\mathrm{\Lambda }1/a`$ a short-distance cutoff for $`k`$. In Eq. (39), the nonzero elements of the (non interacting) Green’s functions read
$$G(k)_{\alpha \beta }=\frac{\delta _{\alpha ,\beta }}{r_s+\frac{\omega ^2}{\mathrm{\Delta }_s}+\rho k^2},$$
(40)
$$G(k)_{1,1}=G(k)_{5,5}=\frac{\rho k^2+r_c}{(\rho k^2+r_c)^2+\frac{\omega ^2}{\eta ^4}},$$
(41)
and
$$G(k)_{5,1}=G(k)_{1,5}=\frac{\frac{\omega }{\eta ^2}}{(\rho k^2+r_c)^2+\frac{\omega ^2}{\eta ^4}},$$
(42)
and the interaction parameters are $`u_{\alpha ,\beta }=u_s`$, $`u_{i,j}=u_c`$, and $`u_{i,\alpha }=u_{cs}`$.
At one loop, integration of the $`\omega 0`$ fields only changes the parameters $`r`$, and $`u`$, similarly to conventional field theory. In the $`T0`$ limit, the change of the former is finite, and merely shifts the transition point. On the other hand, the changes $`\delta u_{ab}`$ in the interaction parameters $`u_{ab}`$ grow logarithmically with decreasing temperature at the critical point. We will, thus, restrict to evaluation of these corrections. These are given by the sum of the usual “loop” diagrams, which give
$$\delta u_{ab}=\frac{1}{2}\underset{c}{}u_{ac}u_{cb}I_{cc}2u_{ab}^2I_{ab}u_{ab}\left(I_{aa}u_{aa}+I_{bb}u_{bb}\right),$$
(43)
where the integrals $`I_{ab}`$ are given by
$$I_{ab}=_{k,\omega 0}G(k)_{aa}G(k)_{bb}.$$
(44)
In Eq. (43) and Eq. (44), we have neglected contributions from nondiagonal parts of Green’s functions Eq. (42), as they only give finite contributions to integrals of the form Eq. (44) in the low-temperature limit. The same holds for integrals containing at least one Green’s function of the superconducting fields Eq. (41). This is due to the fact that for these fields the (bare) dynamical critical exponent $`z`$ is equal to $`2`$, and it does not produce divergences in $`D=3`$. This in turns occurs because the two components of the SC order parameter are canonically conjugate, while the AF ones have independent massive ones. Therefore, we will consider only the divergent contribution
$$I_{\alpha ,\beta }=I_s\frac{1}{8\pi ^2}\sqrt{\frac{\mathrm{\Delta }_s}{\rho ^3}}\mathrm{ln}\frac{\mathrm{\Lambda }\sqrt{\rho \mathrm{\Delta }_s}}{2\pi T},$$
(45)
where we have assumed that we lie outside of the region of influence of the quantum critical fixed point, i. e. $`T\sqrt{r_{c/s}\mathrm{\Delta }_s}`$. Replacing Eq. (45) in Eq. (43), we obtain for the leading contributions
$`\delta u_s={\displaystyle \frac{11}{2}}u_s^2I_s`$ (46)
$`\delta u_c={\displaystyle \frac{3}{2}}u_{cs}^2I_s`$ (47)
$`\delta u_{cs}={\displaystyle \frac{5}{2}}u_su_{cs}I_s.`$ (48)
As expected, quantum fluctuations draw the system away from the $`SO(5)`$ -invariant point even in the case where $`\eta =1`$. This can be seen by adding these corrections to an initially $`SO(5)`$ -invariant system with $`u_c=u_s=u_{cs}=u`$. At the lowest order in the $`u_a`$, the renormalized parameters $`u_a^{}=u_a+\delta u_a`$ obey the relation
$$\mathrm{\Delta }u^{}_{}{}^{}2u_{cs}^{}_{}{}^{}2u_s^{}u_c^{}=2I_su^3>0,$$
(49)
i. e., as in the case of $`\eta 1`$, Eq. (38), the system is drawn into the instability region where a fluctuation-induced first-order transition is expected. This indicates that quantum fluctuations and anisotropy $`\eta 1`$ cooperate in the same direction and draw the system into the instability region, where no finite fixed point is expected. However, for the case where the $`u_a`$ are different, one obtains
$$\mathrm{\Delta }u^{}_{}{}^{}2\mathrm{\Delta }u^2=\frac{u_s}{2}\left(7u_{cs}^211u_cu_s\right)I_s.$$
(50)
Further inserting the values of the $`u_a`$ from Eq. (33) with $`\eta 1`$ (we fix ourselves at the triple point $`r_s=r_c`$), Eq. (50) becomes negative for $`\eta <x_c`$, or $`\eta >1/x_c`$ with $`x_c0.498`$. Therefore, for large difference in the mobilities $`\eta `$, quantum fluctuations tend to shift the renormalized parameters back towards the domain of attraction of the biconical and of the $`SO(5)`$ fixed point.
## V Conclusions
In conclusion, we have analyzed the properties of a projected $`SO(5)`$ model which takes into account the high-energy physics of the Mott-insulating gap. As already pointed out in Ref. , the chemical potential can always be shifted to the AF-SC transition point in order to cancel the symmetry-breaking terms produced by the gap in the classical part of the action. On the other hand, symmetry-breaking terms due to the projection show up in the quantum-mechanical part of the action, as a conjugacy relation between the superconducting components of the superspin vector. A further source of symmetry breaking is due to the different mobility of the hole pairs and of the magnons parametrized by $`\eta 1`$.
Close to the AF/SC/N finite-temperature multicritical point, the quantum effects due to the projection are irrelevant, although subleading symmetry-breaking corrections appear at intermediate length scales. When considered separately, these symmetry-breaking effects both draw the RG flow into a region of instability with first order transitions and no $`SO(5)`$ symmetry. On the other hand, for strong anisotropies $`\eta 0.5`$, quantum corrections partly cancel the symmetry-breaking effects.
There are possibly other effects, such as Coulomb interaction, or fermionic excitations, which can possibly take the system back into the domain of attraction of the biconical fixed point, where $`SO(5)`$ symmetry is only broken by $`20\%`$. Notice that, since the order parameter must be rescaled in order to reach this fixed point, the (possibly approximate) $`SO(5)`$ symmetry reached at this critical point is renormalized, in the sense of Ref. . This means, for example, that the $`SO(5)`$ picture would be consistent with different absolute magnitudes of the SC and AF gaps, as observed experimentally .
## Acknowledgments
This paper is dedicated to Professor Franz Wegner on the occasion of his 60<sup>th</sup> birthday.
We acknowledge many enlightening and pleasant discussions with S. C. Zhang. This work was partially supported by the DFG (HA 1537/17-1).
## A Exact slave-boson treatment of the constraint
The hard-core constraint Eq. (5) becomes (after projecting out the electron pairs)
$$Q(x)=t_\alpha ^{}(x)t_\alpha (x)+t_h^{}(x)t_h(x)+e^{}(x)e(x)1=0.$$
(A1)
The “physical” bosonic operators are then obtained as usual by the replacement
$$t_a(x)t_a(x)e^{}(x),$$
(A2)
(including $`a=h`$) so that the constraint is now conserved by the Hamiltonian. Within the functional integral, the constraint Eq. (A1) can be enforced as usual by adding a “Lagrange multiplier” term $`i_x\lambda (x)Q(x)`$ and integrating over all $`\lambda (x)`$. The partition function can thus be written in terms of an integral over bosonic fields
$$𝒵=𝒟t_\alpha ^{}𝒟t_\alpha 𝒟t_h^{}𝒟t_h𝒟e^{}𝒟e𝑑\lambda \mathrm{exp}S^{},$$
(A3)
with the action
$`S^{}={\displaystyle _0^\beta }d\tau \{{\displaystyle \underset{x}{}}[t_\alpha ^{}(x,\tau )({\displaystyle \frac{}{\tau }}+i\lambda (x))t_\alpha (x,\tau )+t_h^{}(x,\tau )({\displaystyle \frac{}{\tau }}+i\lambda (x))t_h(x,\tau )`$ (A4)
$`+e^{}(x,\tau )({\displaystyle \frac{}{\tau }}+i\lambda (x))e(x,\tau )i\lambda (x)]+H(\tau )\},`$ (A5)
where $`H(\tau )`$ is obtained by replacing Eq. (A2) in Eq. (8) and by replacing all bosonic operators with the corresponding fields at the imaginary time $`\tau `$ (since the Hamiltonian is already normal ordered). In principle, one should take a discretization of the time variable and consider the continuum limit only at the end of the calculation . Notice that the integration of $`\lambda `$ would not give a constraint like Eq. (A1) for the bosonic fields at all imaginary times. Nevertheless, one can proceed in the usual way by carrying out the gauge transformation
$`e(x,\tau )=\overline{e}(x,\tau )e^{i\theta (x,\tau )}`$ (A6)
$`t_a(x,\tau )=\overline{t}_a(x,\tau )e^{i\theta (x,\tau )}`$ (A7)
$`\lambda (x)=\overline{\lambda }(x,\tau )\dot{\theta }(x,\tau ),`$ (A8)
where $`\overline{e}(x,\tau )=|e(x,\tau )|`$. In this way, we can restrict to real values of the boson filed $`e`$ and absorb the time dependence of its phase into a (now) time-dependent $`\lambda `$. Integration over $`\lambda (x,\tau )`$ now leads to the enforcement of the constraint via the $`\delta `$ function (for simplicity, we drop the bar everywhere)
$$\underset{x,\tau }{}\delta \left[|t_\alpha (x,\tau )|^2+|t_h(x,\tau )|^2+e(x,\tau )^21\right]$$
(A9)
at all imaginary times. Integration over $`e(x,\tau )`$ allows one to replace it everywhere in the Hamiltonian, leading to the new action Eq. (16) with Eq. (18).
## B Integration of the momenta
The $`𝐩`$-dependent part of the action has the general form (Cf. Eq. (17) and Eq. (26))
$$S_p=_0^\beta 𝑑\tau 𝑑x\mathrm{\Delta }A(|𝐦(x)|^2)p_\alpha (x)^2ip_\alpha (x)B(x),$$
(B1)
where $`A`$ is a function of the superspin’s magnitude squared (for simplicity, we neglect gradient terms). In order to absorb the coefficient $`A`$, we define new momentum variables
$$p_\alpha ^{}(x)=p_\alpha (x)\sqrt{A(|𝐦(x)|^2)}.$$
(B2)
However, since we don’t want to produce an $`m`$-dependent Jacobian, we carry out a similar transformation for the $`m`$ variables as
$$m_a^{}(x)=m_a(x)g[|𝐦(x)|^2],$$
(B3)
where $`g`$ is chosen in order to have a Jacobian equal to $`1`$. This requirement gives the differential equation
$$A(|𝐦|^2)^{3/2}\left[g(|𝐦|^2)^n2|𝐦|^2g(|𝐦|^2)^{n1}g^{}(|𝐦|^2)\right]=1,$$
(B4)
$`n`$ ($`=5`$) being the number of components of the superspin $`𝐦`$. The solution of this equation is
$$(\sqrt{r}g(r))^n=\frac{n}{2}r^{n/21}A(r)^{3/2}𝑑r,$$
(B5)
where $`r=|𝐦|^2`$. Upon restricting to the lowest order of Eq. (27), $`A(r)=1+\frac{J}{2\mathrm{\Delta }_s}r+𝒪(r^2)`$, we obtain
$$g[|𝐦(x)|^2](1\frac{3J}{7\mathrm{\Delta }_s}|𝐦(x)|^2),$$
(B6)
and its inverse Eq. (29).
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# Reflected Iron Line From a Source Above a Kerr Black Hole Accretion Disc
## 1 Introduction
It is now commonly agreed that Active Galactic Nuclei (AGN) are composed of a central supermassive black hole closely surrounded by a thin equatorial accretion disc, the emission process being controlled by viscous transport of angular momentum (see e.g. Rees 1984; Blandford and Rees 1992). The disc material is relatively cold ($`10^6`$ K) compared to the hard X-ray corona ($`10^8`$ K) by which it is surrounded. The corona is responsible for the observed X-ray continuum emission which follows a power-law of index $`\mathrm{\Gamma }1.7`$. X-ray observation of Seyfert-1 galaxies revealed that there exist spectral features and other deviations from a simple power-law (e.g. Pounds et al. 1989; Nandra et al. 1989; Matsuoka et al. 1990). In particular a strong iron line is observed at 6.4 keV together with an excess in the continuum emission at energies higher than $``$4 keV. This is well explained by Compton reflection re-processing of the hard X-ray power-law continuum onto the accretion disc (see e.g George & Fabian 1991; Matt, Perola & Piro 1991 ). The observed iron line is the result of fluorescent processes and the excess in the continuum is accounted for by the presence of a continuum component in the reflected emission. The detailed spectrum of the iron line is of great importance since it allows one to probe regions of the accretion disc as close as a few Schwarzschild radii from the central black hole (see e.g. Tanaka et al. 1995). In the particular case of MCG–6-30-15, the lineshape yields evidence that the emission occurs as close as half a Schwarzschild radius, giving strong support for the presence of a rotating Kerr black hole (Tanaka et al. 1995, Dabrowski et al. 1997).
The fluorescent iron line profile is well predicted by models which assume a power-law $`ϵ(r)r^\alpha `$ for the fluorescent emissivity on the accretion disc, where $`\alpha 2`$ (see e.g. Tanaka et al. 1995; Bromley, Chen & Miller 1997; Dabrowski et al. 1997; Reynolds & Begelman 1997; Cadez et al. 1998). In such cases the hard X-ray corona is assumed to be an extended source located above and below the disc so that approximately half of the emitted power illuminates the disc and is responsible for the reflected continuum component as well as observed fluorescent lines (George & Fabian 1991). The remaining half escapes to infinity and constitutes the major part of the observed continuum spectra. The predicted iron line equivalent width (EW) is in this case estimated to be $``$100 - 200 eV (e.g. George & Fabian 1991, Matt et al. 1991). However, as discussed in Nandra et al. (1997), for a sample of 18 Seyfert-1 galaxies, observations reveal equivalent widths of 300-600 eV, or even as high as 1 keV in the case of MCG–6-30-15. Part of the EW enhancement may be explained by an overabundance of iron in the accretion disc material (e.g. Lee et al. 1998, 1999). However, other means of enhancement need to be found in order to account for large values of the iron line EW (Reynolds & Fabian 1997). Following the work of Martocchia & Matt (1996) and very recently Martocchia, Karas & Matt (1999), in this paper we propose to investigate the possible EW enhancement due to the gravitational lensing of the primary light rays which illuminate the disc and therefore drive the fluorescent process. The primary hard X-ray source is assumed to be a static point source located above the accretion disc. Situations where the source is located on or off the axis of rotation are both considered here.
Our predictions are based upon Monte Carlo simulations which account for the primary X-ray emission as well the reflected iron line in a consistent manner. It is assumed here that the primary continuum emission dominates over the reflected continuum emission so that the reflected continuum is neglected while calculating the line EW (see Section 4). Details of our fully relativistic theoretical approach are given in Section 2. Some results are given in the form of iron line profiles and equivalent widths in Section 3 for a primary source located on the axis of rotation. In Section 4, this assumption is relaxed and we investigate line profiles and equivalent widths for arbitrary positions of the source above the disc. Finally our results are discussed in Section 5.
## 2 Reflection Model In the Kerr Metric
The black holes considered here belong to the Kerr family of solutions to Einstein’s equation and the metric employed is in the Boyer-Lindquist form, which is defined by
$`ds^2`$ $`=`$ $`dt^2\rho ^2\left({\displaystyle \frac{dr^2}{\mathrm{\Delta }}}+d\theta ^2\right)\left(r^2+L^2\right)\mathrm{sin}^2\theta d\varphi ^2`$ (1)
$`{\displaystyle \frac{2Mr}{\rho ^2}}\left(L\mathrm{sin}^2\theta d\varphi dt\right)^2,`$
where
$$\rho ^2=r^2+L^2\mathrm{cos}^2\theta ,$$
(2)
$$\mathrm{\Delta }=r^22Mr+L^2.$$
(3)
Here $`t`$, $`r`$, $`\theta `$ and $`\varphi `$ are the space-time coordinates, $`M`$ is the geometric mass, the mass of the black hole being $`Mc^2/G`$ and the quantity $`LM`$ is the angular momentum of the hole as measured at infinity (e.g. D’Inverno 1992). The gravity of the disc itself is assumed to be negligible and therefore the whole spacetime is described by the Boyer-Lindquist metric. We employ natural units $`G=c=\mathrm{}=1`$, unless stated otherwise.
### 2.1 Photon Path
Throughout this paper we use the gauge-theoretic formalism of Lasenby, Doran & Gull (1998) which allows fully relativistic and covariant calculations to be made in the context of an easy-to-handle flat space. It is important to note that predictions for measurable quantities agree exactly with those of General Relativity for this case. Equations relevant to an equivalent approach to the problem, but following a more standard general relativistic formalism, are presented in Appendix A.
The gauge-theoretic approach to gravity of Lasenby et al. (1998) employs the language of geometric algebra, which seems to most clearly expose the physics involved. In particular, the geometric algebra of spacetime (Hestenes 1966), known as the Space Time Algebra (STA), is generated by a set of four orthonormal vectors {$`\gamma _\mu `$}, $`\mu =0\mathrm{}3`$, satisfying
$$\gamma _\mu \gamma _\nu =\eta _{\mu \nu }=\mathrm{diag}(+).$$
(4)
In order to investigate problems in relation with the Kerr geometry, we also introduce a suitable set of coordinates $`t`$, $`r`$, $`\theta `$ and $`\varphi `$ associated with the polar coordinate system and defined in terms of the fixed {$`\gamma _\mu `$} frame. The corresponding coordinate frame $`e_\mu _\mu x`$ is
$`e_t`$ $``$ $`\gamma _0`$ (5)
$`e_r`$ $``$ $`\mathrm{sin}\theta \mathrm{cos}\varphi \gamma _1+\mathrm{sin}\theta \mathrm{sin}\varphi \gamma _2+\mathrm{cos}\theta \gamma _3`$ (6)
$`e_\theta `$ $``$ $`r\left(\mathrm{cos}\theta \mathrm{cos}\varphi \gamma _1+\mathrm{cos}\theta \mathrm{sin}\varphi \gamma _2\mathrm{sin}\theta \gamma _3\right)`$ (7)
$`e_\varphi `$ $``$ $`r\mathrm{sin}\theta \left(\mathrm{sin}\varphi \gamma _1+\mathrm{cos}\varphi \gamma _2\right).`$ (8)
Note that $`e_\theta `$ and $`e_\varphi `$ are not unit vectors and differ by factors $`r`$ and $`r\mathrm{sin}\theta `$ from the usual spherical-polar basis vectors, respectively. This is a direct consequence of their definition in terms of the {$`\gamma _\mu `$} frame (see Lasenby et al. 1998).
The translation and rotation gauge fields corresponding to the Boyer-Lindquist metric (1) are given in Doran, Lasenby & Gull (1996).
We start by defining an observer with covariant 4-velocity equal to $`\gamma _0`$. Hereafter this observer will be called the $`\gamma _0`$-observer. One can show that, at infinity, the $`\gamma _0`$-observer is in a flat Minkowski space-time and at rest with respect to the central black hole. Elsewhere its function is to provide a useful frame in which to express quantities. In order to parameterise the energy and direction of the photon path, let us define the photon 4-momentum $`p`$ in the {$`\gamma _\mu `$} frame by
$`p`$ $`=`$ $`\mathrm{\Phi }\gamma _0+\mathrm{\Phi }\mathrm{sin}\theta _p\mathrm{cos}\varphi _p\gamma _1+\mathrm{\Phi }\mathrm{sin}\theta _p\mathrm{sin}\varphi _p\gamma _2+`$ (9)
$`\mathrm{\Phi }\mathrm{cos}\theta _p\gamma _3.`$
The photon energy as measured by the $`\gamma _0`$-observer is given by $`p\gamma _0`$ which, from equation (9), is equal to $`\mathrm{\Phi }`$. Similarly, as illustrated in Fig. 1, the angles $`\theta _p`$ and $`\varphi _p`$ correspond to the usual spherical-polar angles in the local spatial frame of the $`\gamma _0`$-observer and define the direction at which photons are received (or emitted) by the $`\gamma _0`$-observer.
We note that $`p`$ is guaranteed to be null since $`pp=0`$, which is what is required for a massless particle. The photon trajectory itself is parameterised by the affine parameter $`\lambda `$. The spacetime position of the photon ($`t`$, $`r`$, $`\theta `$, $`\varphi `$) is defined by the position vector
$$x=te_t+re_r,$$
(10)
where $`t`$, $`r`$, $`\theta `$ and $`\varphi `$ are scalar functions of $`\lambda `$. The photon geodesic equations are given in Appendix B. The equations we obtain are all first order ordinary differential equations for $`t`$, $`r`$, $`\theta `$, $`\varphi `$, $`\mathrm{\Phi }`$, $`\theta _p`$ and $`\varphi _p`$. As mentioned in the appendix, the attractive aspect of this approach is that the geodesic equations have been obtained without making use of conserved quantities along the photon path such as the energy, angular momentum or Carter constant. Instead these quantities can be used a posteriori as an effective check of both analytical and numerical results (see Appendix B).
### 2.2 Accretion Disc
We now employ a similar method as in Section 2.1 to study the properties of an accretion disc rotating around a Kerr black hole. In particular we are interested in the velocity and energy of the accreted material. The assumptions are the following:
(i) The accretion disc is a thin disc and lies in the equatorial plane ($`\theta =\pi /2`$), perpendicular to the axis of rotation of the central Kerr black hole.
(ii) The orbits are stable circular orbits. In the same way as we defined the photon momentum (9), we start here by defining the 4-velocity vector $`v_d`$ of a particle in circular orbit of radius $`r`$ around the black hole. Since the motion occurs in the equatorial plane, we can parameterise $`v_d`$ as a function of $`\gamma _0`$, $`\gamma _1`$ and $`\gamma _2`$ only:
$$v_d=\mathrm{cosh}U\gamma _0+\left(\mathrm{sinh}U\mathrm{cos}\omega \tau \right)\gamma _1+\left(\mathrm{sinh}U\mathrm{sin}\omega \tau \right)\gamma _2,$$
(11)
where $`\omega `$ represents the angular velocity and $`\tau `$ is the proper time along the world-line of the particle in circular orbit. $`U`$ is a function of $`r`$, $`L`$ and $`M`$, but not a function of $`\tau `$ since the orbit is circular. The parameterisation in $`\mathrm{cosh}`$ and $`\mathrm{sinh}`$ has been chosen to ensure that
$$v_dv_d=\mathrm{cosh}^2U\mathrm{sinh}^2U=1,$$
(12)
as required for a massive particle. In the {$`e_\mu `$} frame, the velocity takes the simple form:
$$v_d=\mathrm{cosh}Ue_t+\frac{\mathrm{sinh}U}{r}e_\varphi .$$
(13)
We note that the $`\gamma _0`$-observer is not comoving with the particle in orbit. Let us however assume that, at a given time $`t`$, the $`\gamma _0`$-observer is at the position ($`r`$, $`\theta =\pi /2`$, $`\varphi `$), i.e. instantaneously at the same position as the orbiting particle (see Fig. 2).
According to the $`\gamma _0`$-observer, the energy per unit mass of the particle is given by $`E_{\gamma _0}=v_d\gamma _0=\mathrm{cosh}U`$. The particle velocity has a magnitude $`V_{\gamma _0}=\mathrm{tanh}U`$ and makes an angle $`\omega \tau =\pi /2+\varphi `$ from the $`\gamma _1`$ axis in the local spatial frame ($`\gamma _1`$, $`\gamma _2`$, $`\gamma _3`$) of the $`\gamma _0`$-observer. When solving the dynamical equations, we find
$$\mathrm{tanh}U=\frac{L\pm \sqrt{Mr}}{\mathrm{\Delta }^{\frac{1}{2}}},$$
(14)
and therefore
$`\mathrm{cosh}U`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Delta }^{\frac{1}{2}}}{\sqrt{r^23Mr\pm 2L\sqrt{Mr}}}},`$
$`\mathrm{sinh}U`$ $`=`$ $`{\displaystyle \frac{L\pm \sqrt{Mr}}{\sqrt{r^23Mr\pm 2L\sqrt{Mr}}}}.`$ (15)
In the case where $`L>0`$, the $`+`$ sign in (14) and (15) is for orbits co-rotating with the spinning hole, while the $``$ sign is for counter-rotating orbits. If $`L<0`$ the significance of the signs is reversed.
Finally, following the same approach as in Appendix B, the conserved energy per unit mass $`E_r`$ and angular momentum $`J_r`$ of a circular orbit can be evaluated by using the Killing vectors $`K_t`$ (equation 51) and $`K_\varphi `$ (equation 52). We find
$$E_r=K_tv_d=E_{\gamma _0}\frac{r^22Mr\pm L\sqrt{Mr}}{r\mathrm{\Delta }^{\frac{1}{2}}}$$
(16)
and
$$J_r=K_\varphi v_d=E_{\gamma _0}\frac{L\mathrm{\Delta }^{\frac{1}{2}}+L^2V_{\gamma _0}}{r}+E_{\gamma _0}V_{\gamma _0}r.$$
(17)
The conserved quantity $`E_r`$ in (16) is related to the local energy $`E_{\gamma _0}`$ by a redshift factor, which tends to unity at infinity. Therefore $`E_r`$ is usually called the energy at infinity. Regarding $`J_r`$ in (17), the first term is a purely gravitational effect due to the dragging effect of the spinning hole, while the second term is more directly related to the Newtonian expression, $`mVr`$, of the angular momentum. We note that the condition for stability is given by $`dE_r/dr>0`$, which gives the following condition for the minimum stable orbit $`r_{\mathrm{ms}}`$:
$$\mathrm{tanh}U(r_{\mathrm{ms}})=\frac{1}{2}.$$
(18)
It is perhaps somewhat surprising that the minimum stable orbit is for a local velocity of one half the speed of light. In this work, the accretion disc extends from $`r=r_{\mathrm{in}}=r_{\mathrm{ms}}`$ up to the arbitrary large outer radius $`r=r_{\mathrm{out}}=5000r_g`$, where $`r_g=GM/c^2`$ is the gravitational radius.
### 2.3 Primary Source and Distant Observer
The primary X-ray source is assumed to be a static point source and, as shown in Fig. 3, is located above the disc at an arbitrary position ($`\rho _s`$, $`\varphi _s`$, $`z_s`$), where $`\rho `$, $`\varphi `$ and $`z`$ are the usual cylindrical polar coordinates. The 4-velocity $`v_s`$ of the source has to be proportional to the Killing vector $`K_t`$ (equation 51) so that the source can remain static (see e.g. Misner et al. 1973). Requiring that $`v_sv_s=1`$, we find
$$v_s=\frac{1}{\sqrt{1\frac{2Mr}{\rho ^2}}}K_t$$
(19)
The source is emitting isotropically in its proper frame with a rate of emission defined by
$$N_{\nu _s}\nu _s^\mathrm{\Gamma }\mathrm{photon}\mathrm{s}^1\mathrm{Hz}^1,$$
(20)
where $`\mathrm{\Gamma }`$ is the photon power-law index. Throughout this work we assume $`\mathrm{\Gamma }=1.7`$ (Mushotzky 1982; Turner & Pounds 1989).
Predictions regarding the iron line flux and equivalent width are those measured by a distant static observer. Since $`K_t=\gamma _0`$ for large $`r`$, the 4-velocity of the distant observer is
$$v_o=\gamma _0.$$
(21)
As illustrated in Fig. 3, the position of the distant observer is given by $`r_o`$, $`\theta _o`$, $`\varphi _o`$ with collecting area on the $`r=r_o`$ sphere equal to
$$\mathrm{\Delta }A=r_o^2\mathrm{sin}\theta _o\mathrm{\Delta }\varphi _o\mathrm{\Delta }\theta _o.$$
(22)
### 2.4 Monte Carlo Simulations
The Monte Carlo simulations we carry out in this work have the advantage of computing both the continuum emission and the reflected iron line flux as measured by a distant observer. In this manner the equivalent width of the line can be estimated in a consistent manner. A simulation consists of sending $`N`$ photons from the source isotropically with respect to its proper frame. $`N_c`$ of these photons will reach directly the observer’s collecting area and therefore contribute to the continuum flux. $`N_l`$ of them will be received after having hit the disc and being reprocessed into iron line photons. These contribute to the observed flux of the line. For each photon received on the collecting area, the total spectrum is computed by considering the total power in a succession of frequency bins of extent $`\mathrm{\Delta }\nu `$. The following two sections (2.4.1 and 2.4.2) describe how the total power in each bin is estimated, first for the continuum emission then the reflected emission.
#### 2.4.1 Continuum Emission
Let us consider a photon emitted by the source with 4-momentum $`p_s`$ (equation 9) which reaches directly the collecting area at radius $`r_o`$ with 4-momentum $`p_o`$. The redshift measured by the distant observer is therefore
$$1+z_{so}=\frac{p_sv_s}{p_ov_o},$$
(23)
Since $`v_o=\gamma _0`$ the denominator is simply the value of $`\mathrm{\Phi }`$ at the observer, $`\mathrm{\Phi }_o`$. The numerator is calculated using (19) and (9). We find
$$1+z_{so}=\frac{\mathrm{\Phi }_s\left[\mathrm{\Delta }^{\frac{1}{2}}L\mathrm{sin}\theta \mathrm{sin}\theta _p\mathrm{sin}(\varphi _p\varphi )\right]_s}{\mathrm{\Phi }_o},$$
(24)
where the subscript $`s`$ indicates that the expression is evaluated at the source.
We define $`W_\nu ^c`$ to be the continuum total power received in the collecting area of the distant observer in the range of frequency \[$`\nu `$, $`\nu +\mathrm{\Delta }\nu `$\]. We have
$$W_\nu ^c=\underset{N_c}{}_\nu ^{\nu +\mathrm{\Delta }\nu }W_{\nu _o}d\nu _o,$$
(25)
where the sum is for each received continuum photon, $`W_{\nu _o}`$ is the received power per unit frequency attached to the photon, and the subscript $`o`$ denotes quantities as measured by the distant observer. Expressing quantities on the right hand side of equation (25) in terms of quantities defined at the source, we obtain
$$W_\nu ^c=\underset{N_c}{}\left(\frac{1}{1+z_{so}}\right)^2_{\nu (1+z_{so})}^{(\nu +\mathrm{\Delta }\nu )(1+z_{so})}\nu _s^{\mathrm{\Gamma }+1}d\nu _s.$$
(26)
#### 2.4.2 Reflected Iron Line
We now consider a photon that leaves the source with a 4-momentum $`p_s`$ and hits the disc with a 4-momentum $`p_d`$. The redshift measured by the co-rotating observer is
$$1+z_{sd}=\frac{p_sv_s}{p_dv_d}.$$
(27)
The numerator is given in (24) and the denominator can be calculated using (9) and (13). We find
$$1+z_{sd}=\frac{\mathrm{\Phi }_s\left[\mathrm{\Delta }^{\frac{1}{2}}L\mathrm{sin}\theta \mathrm{sin}\theta _p\mathrm{sin}(\varphi _p\varphi )\right]_s}{\mathrm{\Phi }_d\left[\mathrm{cosh}U\mathrm{sinh}U\mathrm{sin}\theta _p\mathrm{sin}(\varphi _p\varphi )\right]_d},$$
(28)
where the subscripts $`s`$ and $`d`$ indicate that the expression is evaluated at the source and at the disc, respectively. A similar approach can be followed for the second part of the trajectory, from the disc to the collecting area of the distant observer. We find
$$1+z_{do}=\frac{\mathrm{\Phi }_d\left[\mathrm{cosh}U\mathrm{sinh}U\mathrm{sin}\theta _p\mathrm{sin}(\varphi _p\varphi )\right]_d}{\mathrm{\Phi }_o}.$$
(29)
The radiation processes occurring on the disc are described in detail in George & Fabian (1991). Here, we are simplifying the situation significantly by ignoring the reflected continuum emission so that only the reflected iron line is taken into account (see Section 5 for a discussion of the validity of this approach). We define $`N_{\nu _d}^{in}(\nu _d,\theta _i)`$ to be the total number of photons impinging on the disc with energy $`\nu _d`$ and incidence angle $`\theta _i`$ per unit frequency and per unit time, as measured by a co-rotating observer (see Fig. 3). Following George & Fabian (1991), the number of fluorescent photons per unit frequency and per unit time which are able to escape the disc is
$$N_{\nu _d}^{out}=N_{\nu _d}^{in}(\nu _d,\theta _i)g(\theta _i)f(\nu _d),$$
(30)
where,
$`g(\theta _i)`$ $`=`$ $`\left(6.55.6\mathrm{cos}\theta _i+2.2\mathrm{cos}^2\theta _i\right)\times 10^2,`$ (31)
$`f(\nu _d)`$ $`=`$ $`7.4\times 10^2+2.5\mathrm{exp}\left({\displaystyle \frac{\nu _d1.8}{5.7}}\right).`$ (32)
This analytical approximation is valid for $`\nu _t<\nu _d<\nu _m`$, where $`\nu _t=7.1\mathrm{keV}`$ is the energy threshold for trigerring the $`6.4\mathrm{keV}`$ iron fluorescent emission, and $`\nu _m=30\mathrm{keV}`$. As given in Ghisellini, Haardt & Matt (1994), the intensity of the emerging fluorescent emission is not isotropic and is here assumed to be proportional to $`\mathrm{cos}\theta _e\mathrm{ln}(1+1/\mathrm{cos}\theta _e)`$, where $`\theta _e`$ is the outgoing inclination angle as measured in the co-rotating frame. Therefore $`N_{\nu _d}^{out}`$ follows the angular distribution
$$N_{\nu _d}^{out}(\theta _e)\mathrm{d}\theta _e=2\mathrm{cos}\theta _e\mathrm{ln}\left(1+\frac{1}{\mathrm{cos}\theta _e}\right)\mathrm{sin}\theta _e\mathrm{d}\theta _e.$$
(33)
We now define $`W_\nu ^l`$ to be the iron line total power received in the collecting area of the distant observer in the range of frequency \[$`\nu `$, $`\nu +\mathrm{\Delta }\nu `$\]. Following a similar reasoning as in Section 2.4.1 we obtain
$$W_\nu ^l=\underset{N_l}{}\delta _\nu \left(\frac{1}{1+z_{do}}\right)^2N^{out}\nu _\alpha ,$$
(34)
where $`\nu _\alpha =6.4\mathrm{keV}`$, $`\delta _\nu =1`$ if $`\nu <\nu _\alpha /(1+z_{do})<\nu +\mathrm{\Delta }\nu `$, $`\delta _\nu =0`$ otherwise, and
$$N^{out}=g(\theta _i)_{\nu _t}^{\nu _m}f(\nu _d)N_{\nu _d}^{in}d\nu _d.$$
(35)
When expressed in the source proper frame, we have
$`W_\nu ^l={\displaystyle \underset{N_l}{}}\delta \nu \left({\displaystyle \frac{1}{1+z_{do}}}\right)^2{\displaystyle \frac{1}{1+z_{sd}}}\nu _\alpha g(\theta _i)\times `$
$`{\displaystyle _{\nu _t(1+z_{sd})}^{\nu _m(1+z_{sd})}}f\left({\displaystyle \frac{\nu _s}{1+z_{sd}}}\right)\nu _s^\mathrm{\Gamma }d\nu _s.`$ (36)
## 3 Source on Axis
When discussing the fluorescent iron line profile and equivalent width, many authors have considered the situation where the primary X-ray source is located above the disc, on the axis of rotation (e.g. Martocchia et al. 1999, Matt et al. 1992, Martocchia & Matt 1996; Reynolds & Begelman 1997, Young, Ross & Fabian 1998; Reynolds et al. 1999). This assumption has the advantage of simplifying both calculations and analysis since the axial symmetry of the problem is conserved. In this section we also follow this assumption, but it will be relaxed in Section 4.
### 3.1 Gravitational Lensing Effect
As pointed out by Martocchia & Matt (1996), if the primary source is very close to the hole, a significant fraction of the primary photon rays will be lensed towards the disc, decreasing considerably the direct flux relatively to the reflected flux, as observed by a distant observer. In order to quantify this phenomena, the ratio of the solid angle formed by the rays hitting the disc to the solid angle formed by the rays escaping to infinity ($`\mathrm{\Omega }_{\mathrm{disc}}/\mathrm{\Omega }_{\mathrm{escape}}`$) is plotted versus the height $`h=z_s`$ of the source in Fig. 4. The effect is negligible in the case of a Schwarzschild black hole since a large amount of the lensed radiation is lost within the innermost stable orbit of the disc (i.e. 6 $`r_g`$). In the case of a maximally rotating Kerr black hole, where the inner radius of the disc can be as small as 1 $`r_g`$, the effect becomes significant for a sufficiently close source, typically $`h<6r_g`$ as seen on Fig. 4. However, in the case of a source at height $`h=2r_g`$, $`\mathrm{\Omega }_{\mathrm{disc}}`$ is only a factor of $`5.5`$ bigger than $`\mathrm{\Omega }_{\mathrm{escape}}`$. In this case we find that almost half of the rays bent towards the equatorial plane do not actually hit the disc but are strongly lensed towards the event horizon, resulting in a significant damping of the effect.
### 3.2 Line Profiles
Examples of line profiles obtained with our Monte-Carlo simulations for various values of $`h`$ are given in Figs. 5, 6, 7 and 8 for the maximum Kerr and the Schwarzschild cases. The predicted line shapes are for distant observers at inclination angles $`\theta _0=30`$ degrees and $`\theta _0=60`$ degrees.
As already noticed in previous studies (e.g. Fabian et al. 1989, Dabrowski et al. 1997), the line extends further towards the blue shifted part of the spectrum for larger inclination angles. Indeed, for large $`\theta _0`$ the blue part of the line is boosted by the Doppler effect due to the rotation of the disc. On the other hand, the extension towards the red shifted part of the spectrum is controlled by the strength of the gravitational redshift which increases as the fluorescent emission happens closer to the central black hole. This is very clearly illustrated here for the maximally rotating Kerr case (Figs 7, 8), where $`\nu /\nu _\alpha `$ could be as small as $`0.16`$ for a source located at $`h=1.5r_g`$. Indeed, as the source gets closer to the hole it illuminates more intensively the inner part of the disc, resulting in strong gravitational redshift of the fluorescent line. In the Schwarzschild case as well (Figs. 5, 6) the reddening of the line increases as the source gets closer to the hole. However, because the inner radius of the disc is $`6r_g`$ as opposed to $`1.24r_g`$ in the maximally rotating Kerr case, less strong reddening is predicted in the Schwarzschild case. For example, one can notice very little difference between the predicted lines for $`h=5r_g`$ and $`h=2.5r_g`$. In particular the red wing is not more extended in the $`h=2.5r_g`$ case because most of the illumination is focused within the minimum stable orbit.
### 3.3 Equivalent Width
The equivalent width (EW) of the iron fluorescent line is estimated as follows:
$$\mathrm{EW}=\frac{_0^{\mathrm{}}W_\nu ^ld\nu }{W_{\nu =\nu _\alpha }^c}.$$
(37)
Results are given in Figs. 9 and 10 as a function of the height of the primary source in the maximally rotating Kerr case ($`L=0.998M`$), in the Schwarzschild case ($`L=0`$) and in an intermediate case where $`L=0.5M`$. Fig. 9 is for a distant observer with inclination $`\theta _0=30`$ degrees and Fig. 10 is for $`\theta _0=60`$ degrees.
We expect the strength of the line to be larger as the source gets closer to the hole since the primary emission would preferentially be focused towards the plane of the disc. This effect is clearly visible in the case of a maximally rotating Kerr black hole (solid lines in Figs.9 & 10). However, for lower values of the black hole spin and consequently higher values of the inner radius of the disc, this EW enhancement effect is lessened, as predicted in Section 3.1. In cases where the source is so close to the hole that most of the lensed illumination is lost within the inner radius of the disc, the effect is actually reversed and the EW diminishes as the primary source gets closer to the hole. In particular, no EW enhancement is predicted in the case of a Schwarzschild black hole. For a distant observer at inclination $`\theta _o=30`$ degrees the maximum EW is predicted to be $``$290 eV, while the EW obtained for the classical case where no lensing effect comes into play (i.e. large $`h`$) is $``$165 eV. Therefore, even for a primary source located just above the event horizon of a maximally rotating Kerr black hole, the EW enhancement is less than a factor of two. As seen in Fig.10, the gravitational lensing effect becomes more effective when observed from a large inclination angle. Indeed, the line is strengthened by the lensing effect for both low and high inclination angle points of view. However, due to the same effect, the continuum emission received at large $`\theta _o`$ is strongly reduced compared to that received at smaller $`\theta _o`$. For an observer at inclination $`\theta _o=60`$ degrees, EWs as high as $``$1.2 keV are predicted, which represents an EW enhancement factor of about 8. This strong EW enhancement effect can also be seen in Fig. 8, where the area under each line is proportional to the line EW. However in the case of MCG–6-30-15 the disc is thought to have a inclination angle of about 30 degrees (e.g. Tanaka et al. 1995; Dabrowski et al. 1997).
## 4 Source off Axis
We have seen in Section 3 that the EW enhancement is limited by the fact that a significant proportion of the primary illumination is ‘lost’ within the inner radius of the disc instead of contributing to the strengthen of the iron line. This limiting effect is particularly important when the primary source lies above the disc on the axis of rotation. Stronger EW enhancement should be obtained by relaxing this assumption and allowing the primary source not to lie right above the event horizon of the hole. In this section, we investigate this possibility in the maximally rotating Kerr solution ($`L=0.998M`$).
Since the axial symmetry of the system is here broken, we need to define the distant observer by both the inclination $`\theta _o`$ and the azimuthal angle $`\varphi _o`$. For simplicity, here we fix the observer at $`\varphi _o=\pi /2`$ while allowing the source to be located at various azimuths $`\varphi _s`$ above the disc, as suggested in Fig. 3. Line spectra are given in Fig. 11 as a function of $`\varphi _s`$ for sources located at height $`z_s=1r_g`$ above the plane of the disc and radius $`\rho _s=2.5r_g`$ from the axis of rotation, close to the ergosphere radius at this latitude ($`1.93r_g`$).
Since the illumination takes place above the inner parts of the disc, all these lines display strongly redshifted features. However, as one can notice on Fig. 11, the shape of the spectral line varies significantly as a function of $`\varphi _s`$. For $`\varphi _s=0`$ the line is mostly redshifted and its flux peaks at around 4.5 keV, while for $`\varphi _s=\pi `$ the line is double peaked around 1.8 keV and 6.5 keV. In the former case the primary source is located above the approaching part of the disc where the emission is significantly boosted by the Doppler effect. In the latter case, the strongly redshifted red wing is the result of Doppler shift from the receding part of the disc, while the blue wing is due to rays that have reached the approaching side of the disc while spiralling down on to the disc. Lines for $`\varphi _s=\pi /2`$ and $`\varphi _s=3\pi /2`$ are hybrids of the above two cases.
Since the illumination occurs right above the Doppler boosted region of the disc for $`\varphi _s=0`$, the predicted line EW is much larger in this geometrical situation than for other values of $`\varphi _s`$. This is illustrated in Fig. 11 where the spectral line integrals are proportional to their EW. For the rest of this section we assume that the primary source is always located at this privileged azimuthal angle $`\varphi _s=0`$. Figs. 12 and 13 give the predicted iron line EW as a function of the primary source height $`z_s`$ and radius $`\rho _s`$.
In the case of an observer at inclination $`\theta _o`$=30 degrees, the observed EW can be much larger than for a source located on the axis of rotation. For example, as seen on Fig.12, for a primary source located at ($`\rho _s=2.5r_g`$, $`\varphi _s=0`$, $`z_s=1r_g`$) the predicted EW is $``$1 keV and if the source is as close as $`z_s=0.5r_g`$ to the disc the EW can be as high as $``$1.5 keV. In this latter case, the EW is enhanced by a factor of $``$10 of its classical value. However, in order to reach such high EW the primary source needs to be located above the very inner parts of the disc. For example, in the case of a source located at ($`\rho _s=5r_g`$, $`\varphi _s=0`$, $`z_s=0.5r_g`$) the predicted line EW is only $``$500 eV. As a consequence, very little EW enhancement is expected in the case of slowly rotating black holes. As expected, for an observer at inclination $`\theta _o`$=60 degrees the predicted line EW is much larger than for $`\theta _o`$=30 degrees. As seen on Fig.13, in the extreme case where the primary source is located at ($`\rho _s=2.5r_g`$, $`\varphi _s=0`$, $`z_s=0.5r_g`$), the observed EW can be as high as $``$5.5 keV, which represents an EW enhancement factor of about 35.
## 5 Discussion
In Section 2 we have presented a new fully relativistic approach to calculating photon paths and energies in the Boyer-Lindquist metric. The geodesic equations obtained in this approach are simple first order ordinary differential equations which allow easy and stable numerical integration. Furthermore, the way the photon 4-momentum is parameterised (equation 9) allows us to be very clear regarding observable quantities such as the frequency and incoming/outgoing angles of the light ray, as measured by an observer. We believe that our theoretical approach is an improvement compared to more standard methods usually employed in similar applications (e.g. Reynolds et al. 1999).
We found in Sections 3 and 4 that large values of the observed iron line equivalent width (EW) may be explained by general relativistic effects when the primary X-ray source lies above the very inner part of an accretion disc (within $``$5 gravitational radii) around a maximally rotating central black hole. When the primary source is located on the axis of rotation, EWs up to $``$300 eV only can be accounted for by this model for a disc at inclination 30 degrees. This is a much lower prediction than that of the early work of Martocchia & Matt (1996) who found EWs up to $``$1500 eV for a comparable case. Our results are however in good agreement with the very recent work of Martocchia et al. (1999). At inclination 60 degrees EWs up to $``$1.2 keV are predicted.
By allowing the primary source to be located off the axis of rotation, we found that much larger values of the EW are obtained when the source is located above the approaching side of the disc (up to $``$1.5 keV and $``$5.5 keV for inclination 30 and 60 degrees respectively). However, in the most extreme cases, a reflected continuum component should also be included in the model, in addition to the reflected iron line (Lightman & White 1988; Pounds et al. 1990; Matt et al. 1991; George & Fabian 1991 ). When the reflected emission is not enhanced by the lensing of the primary light rays, the reflected continuum accounts for about 10 percent of the observed continuum flux at 6.4 keV. One would therefore expect the reflected continuum to dominate when the enhancement factor reaches $``$10. Beyond this value, the continuum flux observed at 6.4 keV increases together with the iron line flux and therefore the line should no longer gain in strength. As a result one should regard values of Fig. 13 as indicative of the strength of the lensing enhancement effect, while the actual observed iron line EW should not be larger than $``$ 1600 eV (i.e. 10 times its classical value).
Results found in Section 4 may be used as arguments towards explaining the large iron line EW observed in Seyfert-1 galaxies. Such an explanation may be valid for the EW enhancement observed in the case of MCG–6-30-15 where it is likely that strong emission processes take place very close to the hole, at least during specific periods (see Iwasawa et al. 1996, 1999). However, as highlighted by Reynolds & Fabian (1997), in the case of more typical objects it is less probable that strong emission would occur so close to the hole. Reynolds & Fabian (1997) found that the relative motion between the disc and the source may significantly affect the EW of the line, simply by special relativistic arguments. This could allow EW enhancement even when the emission occurs further away from the hole. Such an effect could be investigated in our model by allowing the source to have a peculiar velocity. For example, a simple model where the source moves with a constant angular velocity could be assumed (Yu & Lu 1999).
## Acknowledgements
YD would like to thank Trinity Hall, Cambridge for support in the form of a Research Fellowship. The authors also thank Andy Fabian and Anthony Challinor for useful discussions.
## Appendix A General Relativistic Approach
This work has been carried out using the gauge-theoretic formalism described in detail in Lasenby et al. (1998). However, the theoretical calculations presented in Section 2 and Appendix B could in principle be carried out following a more standard general relativistic approach by making use of orthogonal tetrads (see e.g. Misner et al. 1973). A suitable tetrad frame in the Boyer-Lindquist metric (equation 1) can be defined as follows:
$`\left(e_t\right)^\mu `$ $`=`$ $`({\displaystyle \frac{r^2+L^2}{\rho \mathrm{\Delta }^{\frac{1}{2}}}},0,0,{\displaystyle \frac{L}{\rho \mathrm{\Delta }^{\frac{1}{2}}}})`$ (38)
$`\left(e_r\right)^\mu `$ $`=`$ $`(0,{\displaystyle \frac{\mathrm{\Delta }^{\frac{1}{2}}}{\rho }},0,0)`$ (39)
$`\left(e_\theta \right)^\mu `$ $`=`$ $`(0,0,{\displaystyle \frac{r}{\rho }},0)`$ (40)
$`\left(e_\varphi \right)^\mu `$ $`=`$ $`({\displaystyle \frac{Lr\mathrm{sin}^2\theta }{\rho }},0,0,{\displaystyle \frac{r}{\rho }}),`$ (41)
where $`\mu =0\mathrm{}3`$ are the usual coordinate indexes, corresponding here to the Boyer-Lindquist coordinates ($`t`$, $`r`$, $`\theta `$, $`\varphi `$). This tetrad can be identified with the coordinate frame {$`e_\mu `$} defined by equations (5) to (8). The orthonormal tetrad corresponding to the {$`\gamma _\mu `$} frame (equation 4) is given by
$`\left(\gamma _0\right)^\mu `$ $`=`$ $`({\displaystyle \frac{r^2+L^2}{\rho \mathrm{\Delta }^{\frac{1}{2}}}},0,0,{\displaystyle \frac{L}{\rho \mathrm{\Delta }^{\frac{1}{2}}}})`$ (42)
$`\left(\gamma _1\right)^\mu `$ $`=`$ $`({\displaystyle \frac{L\mathrm{sin}\theta \mathrm{sin}\varphi }{\rho }},{\displaystyle \frac{\mathrm{\Delta }^{\frac{1}{2}}\mathrm{sin}\theta \mathrm{cos}\varphi }{\rho }},`$ (43)
$`{\displaystyle \frac{\mathrm{cos}\theta \mathrm{cos}\varphi }{\rho }},{\displaystyle \frac{\mathrm{sin}\varphi }{\rho \mathrm{sin}\theta }})`$
$`\left(\gamma _2\right)^\mu `$ $`=`$ $`({\displaystyle \frac{L\mathrm{sin}\theta \mathrm{cos}\varphi }{\rho }},{\displaystyle \frac{\mathrm{\Delta }^{\frac{1}{2}}\mathrm{sin}\theta \mathrm{sin}\varphi }{\rho }},`$ (44)
$`{\displaystyle \frac{\mathrm{cos}\theta \mathrm{sin}\varphi }{\rho }},{\displaystyle \frac{\mathrm{cos}\varphi }{\rho \mathrm{sin}\theta }})`$
$`\left(\gamma _3\right)^\mu `$ $`=`$ $`(0,{\displaystyle \frac{\mathrm{\Delta }^{\frac{1}{2}}\mathrm{cos}\theta }{\rho }},{\displaystyle \frac{\mathrm{sin}\theta }{\rho }},0).`$ (45)
The equations given in Section 2 can now be translated into this general relativistic context. For example equation (9) of Section 2.1 should be read as follows
$`p^\mu `$ $`=`$ $`\mathrm{\Phi }\left(\gamma _0\right)^\mu +\mathrm{\Phi }\mathrm{sin}\theta _p\mathrm{cos}\varphi _p\left(\gamma _1\right)^\mu +`$ (46)
$`\mathrm{\Phi }\mathrm{sin}\theta _p\mathrm{sin}\varphi _p\left(\gamma _2\right)^\mu +\mathrm{\Phi }\mathrm{cos}\theta _p\left(\gamma _3\right)^\mu .`$
While the photon geodesic equations $`\dot{t}`$, $`\dot{r}`$, $`\dot{\theta }`$, $`\dot{\varphi }`$ given in Appendix B are obtained by the following equation
$$p^\mu =\frac{\left(\mathrm{d}x\right)^\mu }{\mathrm{d}\lambda },$$
(47)
where
$$\left(\mathrm{d}x\right)^\mu =(\mathrm{d}t,\mathrm{d}r,\mathrm{d}\theta ,\mathrm{d}\varphi ).$$
(48)
## Appendix B Photon Geodesic Equations
When solving the geodesic equations in the Boyer-Lindquist metric for the photon 4-momentum given in Section 2.1 (see Lasenby et al. 1998), we find
$`{\displaystyle \frac{\mathrm{d}t}{\mathrm{d}\lambda }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }}{\rho \mathrm{\Delta }^{\frac{1}{2}}}}\left[L^2+r^2+L\mathrm{\Delta }^{\frac{1}{2}}\mathrm{sin}\theta \mathrm{sin}\theta _p\mathrm{sin}(\varphi _p\varphi )\right],`$
$`{\displaystyle \frac{\mathrm{d}r}{\mathrm{d}\lambda }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }\mathrm{\Delta }^{\frac{1}{2}}}{\rho }}\left[\mathrm{cos}\theta _p\mathrm{cos}\theta +\mathrm{sin}\theta \mathrm{sin}\theta _p\mathrm{cos}(\varphi _p\varphi )\right],`$ (49)
$`{\displaystyle \frac{\mathrm{d}\theta }{\mathrm{d}\lambda }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }}{\rho }}\left[\mathrm{cos}\theta \mathrm{sin}\theta _p\mathrm{cos}(\varphi _p\varphi )\mathrm{sin}\theta \mathrm{cos}\theta _p\right],`$
$`{\displaystyle \frac{\mathrm{d}\varphi }{\mathrm{d}\lambda }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }}{\rho \mathrm{\Delta }^{\frac{1}{2}}\mathrm{sin}\theta }}\left[\mathrm{\Delta }^{\frac{1}{2}}\mathrm{sin}\theta _p\mathrm{sin}(\varphi _p\varphi )+L\mathrm{sin}\theta \right],`$
and for the photon 4-momentum
$`{\displaystyle \frac{\mathrm{d}\mathrm{\Phi }}{\mathrm{d}\lambda }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }^2}{\rho ^3}}\left\{\left[\mathrm{cos}\theta \mathrm{cos}\theta _p+\mathrm{sin}\theta \mathrm{sin}\theta _p\mathrm{cos}(\varphi _p\varphi )\right](2\rho ^2r^2r\mathrm{\Delta }^{\frac{1}{2}}+B)L^2\mathrm{cos}\theta \mathrm{cos}\theta _p\right\},`$
$`{\displaystyle \frac{\mathrm{d}\varphi _p}{\mathrm{d}\lambda }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }}{\rho ^3}}\left\{B{\displaystyle \frac{\mathrm{sin}(\varphi _p\varphi )\mathrm{sin}\theta }{\mathrm{sin}\theta _p}}+LA\left[{\displaystyle \frac{\mathrm{cos}(\varphi _p\varphi )\mathrm{sin}2\theta }{2\mathrm{tan}\theta _p}}\mathrm{cos}\theta ^2\right]+L\rho ^2\mathrm{\Delta }^{\frac{1}{2}}\right\},`$ (50)
$`{\displaystyle \frac{\mathrm{d}\theta _p}{\mathrm{d}\lambda }}`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Phi }\mathrm{cos}\theta }{\rho ^3}}\left\{B\left[\mathrm{tan}\theta \mathrm{cos}(\varphi _p\varphi )\mathrm{cos}\theta _p\mathrm{sin}\theta _p\right]+LA\mathrm{sin}\theta \mathrm{sin}(\varphi _p\varphi )\right\},`$
where
$`A`$ $`=`$ $`2L\mathrm{sin}\theta _p\mathrm{sin}(\varphi _p\varphi )\mathrm{sin}\theta +2\mathrm{\Delta }^{\frac{1}{2}},`$
$`B`$ $`=`$ $`Ar+\rho ^2\mathrm{\Delta }^{\frac{1}{2}}\left(Mr\mathrm{\Delta }^{\frac{1}{2}}\right).`$
We note the very simple form of the ordinary first order differential equations obtained for $`\dot{t}`$, $`\dot{r}`$, $`\dot{\theta }`$ and $`\dot{\varphi }`$, where overdots denote differentiation with respect to $`\lambda `$. The $`\dot{\mathrm{\Phi }}`$, $`\dot{\theta }_p`$ and $`\dot{\varphi }_p`$ expressions are slightly more complex but still are first order ordinary differential equations. This approach is therefore useful because of the computational advantage of first order equations in terms of efficiency and stability.
An interesting aspect of the geodesic equations (B) and (B) is that they do not rely on conserved quantities such as the energy, angular momentum or Carter constant along the photon path. As a consequence, one can actually compute these quantities and verify that they are indeed conserved along the trajectory. This allows one to check easily the analytical expressions (B) and (B) as well as the corresponding numerical code. The conserved energy $`E_p`$ and angular momentum $`J_p`$ of a photon geodesic are evaluated by making use of the Killing vectors $`K_t`$ and $`K_\varphi `$ associated with the temporal and axial symmetries of (1) respectively. We have
$$K_t=\frac{\mathrm{\Delta }^{\frac{1}{2}}}{\rho }e_t\frac{L}{\rho r}e_\varphi .$$
(51)
and
$$K_\varphi =\frac{L\mathrm{\Delta }^{\frac{1}{2}}\mathrm{sin}^2\theta }{\rho }e_t+\frac{r^2+L^2}{\rho r}e_\varphi .$$
(52)
Therefore
$$E_p=pK_t=\mathrm{\Phi }\left[\frac{L\mathrm{sin}\theta \mathrm{sin}\theta _p\mathrm{sin}(\varphi _p\varphi )+\mathrm{\Delta }^{\frac{1}{2}}}{\rho }\right]$$
(53)
and
$$J_p=pK_\varphi =\mathrm{\Phi }\mathrm{sin}\theta \left[\frac{(L^2+r^2)\mathrm{sin}\theta _p\mathrm{sin}(\varphi _p\varphi )+L\mathrm{\Delta }^{\frac{1}{2}}\mathrm{sin}\theta }{.}\rho \right]$$
(54)
The derivatives $`\dot{E_p}`$ and $`\dot{J_p}`$ with respect to the affine parameter $`\lambda `$ can be obtained. As expected, substituting for the expressions of $`\dot{r}`$, $`\dot{\theta }`$, $`\dot{\varphi }`$, $`\dot{\mathrm{\Phi }}`$, $`\dot{\theta }_p`$ and $`\dot{\varphi }_p`$ found above leads to
$$\frac{dE_p}{d\lambda }=\frac{dJ_p}{d\lambda }=0.$$
(55)
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# DESY 00-047 ISSN 0418-9833 hep-ph/0005060 May 2000 Gauge-Independent 𝑊-Boson Partial Decay Widths
## 1 Introduction
The discovery of the $`W`$ boson at the CERN S$`p\overline{p}`$S collider in 1983 was one of the great successes of the standard model (SM) of the electroweak interactions. The properties of the $`W`$ boson, including its partial decay widths, have been extensively studied at the CERN S$`p\overline{p}`$S collider and the Fermilab Tevatron, and since 1995 also at the CERN Large Electron-Positron Collider (LEP2). On the theoretical side, the one-loop QED and electroweak radiative corrections to the partial decay widths of the $`W`$ boson were calculated in the light-fermion approximation in Refs. , respectively, and for finite fermion masses in Ref. . The QCD corrections, which are present for the hadronic decay modes, were computed at one loop for arbitrary quark masses in Refs. . The two- and three-loop QCD corrections for massless quarks may be extracted from Ref. , and the respective terms proportional to $`m_q^2/M_W^2`$, where $`m_q`$ is a generic quark mass, may be found in Ref. .
In the calculation of the electroweak corrections, the treatment of the Cabibbo-Kobayashi-Maskawa (CKM) mixing matrix , which rotates the weak eigenstates of the quark fields into their mass eigenstates, deserves special attention. In the approximation of neglecting the down-quark masses against the $`W`$-boson mass, the CKM matrix can be taken to be unity, so that it does not need to be renormalized. In fact, this avenue was taken in Refs. . As a matter of principle, however, the CKM matrix elements must be renormalized because they are parameters of the bare Lagrangian. This was realized for the Cabibbo angle in the SM with two fermion generations in a pioneering paper by Marciano and Sirlin . A compact and plausible on-shell renormalization prescription for the CKM matrix of the three-generation SM was proposed in Ref. on the basis of a detailed inspection of the ultraviolet singularities that are left over in the one-loop expressions for the hadronic decay widths of the $`W`$ boson if one does not include renormalization constants for the CKM matrix elements appearing in the tree-level formulas. Since the analysis of Ref. was performed in a specific gauge, namely in ’t Hooft-Feynman gauge, it could not be checked if the final one-loop results are gauge independent as required. A recent analysis in $`R_\xi `$ gauge has revealed that this is not the case . In fact, the finite parts of renormalization constants for the CKM matrix elements as defined in Ref. do depend on the gauge parameter $`\xi _W`$ associated with the $`W`$ boson, and so do the renormalized CKM matrix elements. In Refs. , an alternative renormalization prescription for the CKM matrix was proposed which, at one loop, is consistent with the relevant Ward-Takahashi identities and avoids this problem.
In this paper, we explicitly calculate the $`W`$-boson partial decay widths at one loop in the on-shell renormalization scheme supplemented by the alternative renormalization prescription for the CKM matrix . For simplicity, we henceforth refer to this scheme as the on-shell scheme. As in Ref. , we keep the full dependence on the fermion masses. We work in $`R_\xi `$ gauge, with arbitrary gauge parameters $`\xi _W`$, $`\xi _Z`$, and $`\xi _A`$, and verify that the final results are indeed independent of them. Using up-to-date information on the input parameters, we present quantitative predictions for the various leptonic and hadronic decay widths of the $`W`$ boson, which can be readily confronted with precise experimental data from the Tevatron and from LEP2. Furthermore, we establish the relationships between the CKM matrix elements renormalized according to the modified minimal subtraction ($`\overline{\mathrm{MS}}`$) scheme and their counterparts in the on-shell scheme. We also provide similar relationships appropriate for the Wolfenstein parameterization. From these relationships, we recover the beta functions of the CKM matrix elements , which may be relevant for studies within grand unified theories.
This paper is organized as follows. In Sec. 2, we establish our formalism and outline our calculation. In Sec. 3, we exhibit the relationships between the CKM matrix elements defined in the on-shell and $`\overline{\mathrm{MS}}`$ schemes and extract their beta functions. In Sec. 4, we present our quantitative predictions. Our conclusions are contained in Sec. 5. In the Appendix, we list the fermion two-point functions at one loop in $`R_\xi `$ gauge.
## 2 Analytical Results
We now describe our analytical analysis, largely adopting the notations from Ref. . We consider the two-particle decay of the $`W^+`$ boson to generic leptons or quarks,
$$W^+(k)f_i(p_1)\overline{f}_j^{}(p_2),$$
(1)
where $`f_i=\nu _e,\nu _\mu ,\nu _\tau ,u,c`$, $`f_j^{}=e,\mu ,\tau ,d,s,b`$, the bar indicates antiparticles, and the four-momenta are specified in parentheses. We wish to calculate the electroweak and QCD one-loop corrections to the partial decay width of process (1), with finite fermion masses $`m_{f,i}`$ and $`m_{f^{},j}`$ and general CKM matrix $`V_{ij}`$. The result for the charge-conjugate process, $`W^{}\overline{f}_if_j^{}`$, will be the same. The relevant standard matrix elements read
$`_1^\sigma `$ $`=`$ $`\overline{u}(p_1)\overline{)}\epsilon (k)\omega _\sigma v(p_2),`$
$`_2^\sigma `$ $`=`$ $`\overline{u}(p_1)\omega _\sigma v(p_2)\epsilon (k)p_1,`$ (2)
where $`\omega _\pm =(1\pm \gamma _5)/2`$, $`\epsilon (k)`$ is the polarization four-vector of the $`W^+`$ boson, and $`\overline{u}(p_1)`$ and $`v(p_2)`$ are the spinors of the fermions $`f_i`$ and $`\overline{f}_j^{}`$, respectively. It is convenient to define
$`G_1^{}`$ $`=`$ $`{\displaystyle \underset{\mathrm{pol}}{}}_1^{}_1^{}=2M_W^2m_{f,i}^2m_{f^{},j}^2{\displaystyle \frac{\left(m_{f,i}^2m_{f^{},j}^2\right)^2}{M_W^2}},`$
$`G_1^+`$ $`=`$ $`{\displaystyle \underset{\mathrm{pol}}{}}_1^{}_1^+=6m_{f,i}m_{f^{},j},`$
$`G_2^{}`$ $`=`$ $`{\displaystyle \underset{\mathrm{pol}}{}}_1^{}_2^{}={\displaystyle \frac{m_{f,i}}{2}}{\displaystyle \frac{\kappa ^2(M_W^2,m_{f,i}^2,m_{f^{},j}^2)}{M_W^2}},`$
$`G_2^+`$ $`=`$ $`{\displaystyle \underset{\mathrm{pol}}{}}_1^{}_2^+={\displaystyle \frac{m_{f^{},j}}{2}}{\displaystyle \frac{\kappa ^2(M_W^2,m_{f,i}^2,m_{f^{},j}^2)}{M_W^2}},`$ (3)
where $`\kappa (x,y,z)=\sqrt{x^2+y^2+z^22(xy+yz+zx)}`$ is the Källén function and it issummed over the polarization states of the $`W^+`$ boson and the spin states of the fermions $`f_i`$ and $`f_j^{}`$. Let us first deal with the quark case, which is more involved. In the Born approximation, the transition ($`T`$) matrix element of process (1) then reads
$$_0^{Wf_if_j^{}}=\frac{eV_{ij}}{\sqrt{2}s_w}_1^{},$$
(4)
where $`e`$ is the electron charge magnitude and $`s_w^2=1c_w^2=1M_W^2/M_Z^2`$. Thus, the partial decay width is given by
$$\mathrm{\Gamma }_0^{Wf_if_j^{}}=\frac{N_C^f\alpha |V_{ij}|^2}{24s_w^2M_W^3}\kappa (M_W^2,m_{f,i}^2,m_{f^{},j}^2)G_1^{},$$
(5)
where $`N_C^f=3`$ and $`\alpha =e^2/(4\pi )`$ is Sommerfeld’s fine-structure constant.
The one-loop-corrected $`T`$-matrix element of process (1) emerges from Eq. (4) by including the renormalization constants for the parameters $`e`$, $`s_w`$, and $`V_{ij}`$, those for the $`W^+`$, $`f_i`$, and $`f_j^{}`$ fields, and the proper vertex correction. It reads
$`_1^{Wf_if_j^{}}`$ $`=`$ $`{\displaystyle \frac{eV_{ij}}{\sqrt{2}s_w}}\{_1^{}[1+{\displaystyle \frac{\delta e}{e}}{\displaystyle \frac{\delta s_w}{s_w}}+{\displaystyle \frac{\delta V_{ij}}{V_{ij}}}+{\displaystyle \frac{1}{2}}\delta Z_W+{\displaystyle \frac{1}{2}}{\displaystyle \underset{k}{}}(\delta Z_{ik}^{f,L}V_{kj}+V_{ik}\delta Z_{kj}^{f^{},L})]`$ (6)
$`+{\displaystyle \underset{a=1}{\overset{2}{}}}{\displaystyle \underset{\sigma =\pm }{}}_a^\sigma \delta F_a^\sigma (M_W,m_{f,i},m_{f^{},j})\},`$
where $`\delta F_a^\sigma `$ are electroweak form factors. The various renormalization constants may be expressed in terms of the unrenormalized, one-particle-irreducible two-point functions of the gauge bosons and fermions, as
$`{\displaystyle \frac{\delta e}{e}}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{s_w}{c_w}}\delta Z_{ZA}+\delta Z_{AA}\right)`$ (7)
$`=`$ $`{\displaystyle \frac{s_w}{c_w}}{\displaystyle \frac{\mathrm{\Pi }^{ZA}(0)}{M_Z^2}}+{\displaystyle \frac{1}{2}}{\displaystyle \frac{\mathrm{\Pi }^{AA}(k^2)}{k^2}}|_{k^2=0},`$
$`{\displaystyle \frac{\delta s_w}{s_w}}`$ $`=`$ $`{\displaystyle \frac{c_w^2}{2s_w^2}}\left({\displaystyle \frac{\delta M_W^2}{M_W^2}}{\displaystyle \frac{\delta M_Z^2}{M_Z^2}}\right)`$ (8)
$`=`$ $`{\displaystyle \frac{c_w^2}{2s_w^2}}\stackrel{~}{\text{Re}}\left({\displaystyle \frac{\mathrm{\Pi }^{WW}\left(M_W^2\right)}{M_W^2}}{\displaystyle \frac{\mathrm{\Pi }^{ZZ}\left(M_Z^2\right)}{M_Z^2}}\right),`$
$`\delta V_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{2}}\{{\displaystyle \underset{ki}{}}{\displaystyle \frac{m_{f,i}^2+m_{f,k}^2}{m_{f,i}^2m_{f,k}^2}}[\mathrm{\Sigma }_{ik}^{f,L}(0)+2\mathrm{\Sigma }_{ik}^{f,S}(0)]V_{kj}`$ (9)
$`{\displaystyle \underset{kj}{}}V_{ik}{\displaystyle \frac{m_{f^{},k}^2+m_{f^{},j}^2}{m_{f^{},k}^2m_{f^{},j}^2}}[\mathrm{\Sigma }_{kj}^{f^{},L}(0)+2\mathrm{\Sigma }_{kj}^{f^{},S}(0)]\},`$
$`\delta Z_W`$ $`=`$ $`\stackrel{~}{\text{Re}}{\displaystyle \frac{\mathrm{\Pi }^{WW}(k^2)}{k^2}}|_{k^2=M_W^2},`$ (10)
$`\delta Z_{ij}^{f,L}`$ $`=`$ $`{\displaystyle \frac{2}{m_{f,i}^2m_{f,j}^2}}\stackrel{~}{\text{Re}}[m_{f,i}^2\mathrm{\Sigma }_{ij}^{f,L}\left(m_{f,j}^2\right)+m_{f,i}m_{f,j}\mathrm{\Sigma }_{ij}^{f,R}\left(m_{f,j}^2\right)`$ (11)
$`+(m_{f,i}^2+m_{f,j}^2)\mathrm{\Sigma }_{ij}^{f,S}\left(m_{f,j}^2\right)],ij,`$
$`\delta Z_{ii}^{f,L}`$ $`=`$ $`\stackrel{~}{\text{Re}}\mathrm{\Sigma }_{ii}^{f,L}\left(m_{f,i}^2\right)m_{f,i}^2{\displaystyle \frac{}{p^2}}\stackrel{~}{\text{Re}}\left[\mathrm{\Sigma }_{ii}^{f,L}(p^2)+\mathrm{\Sigma }_{ii}^{f,R}(p^2)+2\mathrm{\Sigma }_{ii}^{f,S}(p^2)\right]_{p^2=m_{f,i}^2}.`$ (12)
Here, $`\mathrm{\Pi }^{BB^{}}(k^2)`$ is the transverse coefficient of the two-point function, at four-momentum $`k`$, of the gauge bosons $`B`$ and $`B^{}`$, and $`\mathrm{\Sigma }_{ij}^{f,L/R/S}(p^2)`$ are the left-handed, right-handed, and scalar coefficients of the two-point function, at four-momentum $`p`$, of the fermions $`f_i`$ and $`f_j`$, respectively. The latter are listed in Eq. (A.2) of the Appendix. The symbol $`\stackrel{~}{\text{Re}}`$ takes the dispersive parts of the loop integrals appearing in the two-point functions and commutes with complex-valued parameters, such as $`V_{ij}`$. As a consequence, $`\delta Z_{ij}^{f,L}=\delta Z_{ji}^{f,L}`$ is obtained from $`\delta Z_{ji}^{f,L}`$ by complex conjugation of the CKM matrix elements contained therein. From Eq. (A.2) it hence follows that $`\delta Z_{ij}^{f,L}`$ emerges from $`\delta Z_{ij}^{f,L}`$ through the interchange of $`m_{f,i}`$ and $`m_{f,j}`$. In particular, we have $`\delta Z_{ii}^{f,L}=\delta Z_{ii}^{f,L}`$. Notice that the vertex correction only depends linearly on $`V_{ij}`$, which is factored out in Eq. (6).<sup>1</sup><sup>1</sup>1The last term in Eq. (9.6) of Ref. should be multiplied with $`V_{ij}`$.
All the renormalization constants appearing in Eq. (6) are ultraviolet (UV) divergent. If the renormalized parameters $`e`$, $`s_w`$, $`V_{ij}`$ are to represent physical observables, they must be gauge independent, and so must the respective renormalization constants. This is well established for $`\delta e/e`$ and $`\delta s_w/s_w`$ given by Eqs. (7) and (8), respectively, while an appropriate definition of $`\delta V_{ij}/V_{ij}`$, namely Eq. (9), was proposed only recently . On the other hand, the field renormalization constants in Eq. (6) are gauge dependent and, with the exception of $`\delta Z_{ij}^{f,L}`$ with $`ij`$, also infrared (IR) divergent. Since $`_1^+`$ and $`_2^\pm `$ do not yet appear in Eq. (4), the respective form factors are finite and gauge independent. However, $`\delta F_1^{}`$ is IR and UV divergent and gauge dependent. The right-hand side of Eq. (6) is UV finite and gauge independent , but it is IR divergent. This IR divergence is cancelled in the one-loop expression for the partial decay width by the real bremsstrahlung correction $`\delta _\mathrm{b}^{\mathrm{ew}}`$. Also including the one-loop QCD correction $`\delta ^{\mathrm{QCD}}`$, we have
$$\mathrm{\Gamma }_1^{Wf_if_j^{}}=\mathrm{\Gamma }_0^{Wf_if_j^{}}\left(1+\delta ^{\mathrm{ew}}+\delta ^{\mathrm{QCD}}\right),$$
(13)
where
$$\delta ^{\mathrm{ew}}=\delta _{\mathrm{virt}}^{\mathrm{ew}}+\delta _\mathrm{b}^{\mathrm{ew}},$$
(14)
with the virtual electroweak correction
$`\delta _{\mathrm{virt}}^{\mathrm{ew}}`$ $`=`$ $`2{\displaystyle \frac{\delta e}{e}}2{\displaystyle \frac{\delta s_w}{s_w}}+2{\displaystyle \frac{\delta V_{ij}}{V_{ij}}}+\delta Z_W+{\displaystyle \underset{k}{}}\left(\delta Z_{ik}^{f,L}V_{kj}+V_{ik}\delta Z_{kj}^{f^{},L}\right)`$ (15)
$`+{\displaystyle \frac{2}{G_1^{}}}{\displaystyle \underset{a=1}{\overset{2}{}}}{\displaystyle \underset{\sigma =\pm }{}}G_a^\sigma \delta F_a^\sigma (M_W,m_{f,i},m_{f^{},j}).`$
Note that $`\delta _\mathrm{b}^{\mathrm{ew}}`$ is gauge independent because $`\mathrm{\Gamma }_0^{Wf_if_j^{}}\delta _\mathrm{b}^{\mathrm{ew}}`$ represents the partial decay width of the physical process $`W^+f_i\overline{f}_j^{}\gamma `$ in the Born approximation. Thus, $`\delta ^{\mathrm{ew}}`$ is finite and gauge independent, as it must be because it exhausts the leading-order electroweak correction to the physical process $`W^+f_i\overline{f}_j^{}(\gamma )`$. The QCD correction $`\delta ^{\mathrm{QCD}}`$ may be obtained from $`\delta ^{\mathrm{ew}}`$ by retaining only the terms containing $`Q_f^2`$, $`Q_f^{}^2`$, or $`Q_fQ_f^{}`$, setting $`Q_f=Q_f^{}=1`$, replacing $`\alpha `$ with the strong-coupling constant $`\alpha _s^{(n_f)}(\mu )`$, and including the overall colour factor $`C_F=4/3`$. Also $`\delta ^{\mathrm{QCD}}`$ is IR and UV finite and gauge independent, as it must be because it exhausts the leading-order QCD correction to the physical process $`W^+f_i\overline{f}_j^{}(g)`$. In the limit $`m_{f,i}=m_{f^{},j}=0`$, the well-known correction factor $`\delta ^{\mathrm{QCD}}=1+\alpha _s^{(n_f)}(\mu )/\pi `$ is recovered.
In the lepton case, we have $`N_C^f=1`$ and $`\delta ^{\mathrm{QCD}}=0`$. Furthermore, we have $`V_{ij}=\delta _{ij}`$ and $`\delta V_{ij}=0`$ if we assume the neutrinos to be massless.
We computed all ingredients of Eq. (13) in $`R_\xi `$ gauge, with arbitrary gauge parameters $`\xi _W`$, $`\xi _Z`$, and $`\xi _A`$, for finite fermion masses $`m_{f,i}`$ and $`m_{f^{},j}`$ and general CKM matrix $`V_{ij}`$. We regularized the UV divergences by means of dimensional regularization (DR), in $`D=42ϵ`$ space-time dimensions, and the IR ones by introducing an infinitesimal photon mass $`\lambda `$. To guarantee the correctness of our results, we chose two independent approaches. The first approach was based on the program packages FeynArts and FeynCalc , which are written in Mathematica. FeynArts generates the relevant Feynman diagrams and translates them into $`T`$-matrix elements, in a format which is readable by FeynCalc. FeynCalc then simplifies the expressions and decomposes them into the standard one-loop scalar integrals $`A_0`$, $`B_0`$, and $`C_0`$. The second approach was to essentially perform all the calculations by hand using well-tested custom-made programs, written in FORM , in the intermediate steps. Both approaches led to identical results. Our results for the gauge-boson self-energies fully agree with those listed in Eqs. (7)–(10) of Ref. and will not be presented here. The corresponding formulas in ’t Hooft-Feynman gauge, with $`\xi _W=\xi _Z=\xi _A=1`$, may be found in Appendix B of Ref. . Generic expressions for the fermion self-energies, which were originally derived in Ref. , are specified in Eq. (A.2) of the Appendix. The diagonal fermion wave-function renormalization constants $`\delta Z_{ii}^{f,L}`$ suffer from IR divergences. Therefore, we retained the infinitesimal photon mass $`\lambda `$ in those parts of Eq. (A.2), from which IR divergences may arise. In the limit $`\lambda =0`$, Eq. (A.2) agrees with Eqs. (B1)–(B3) of Ref. .<sup>2</sup><sup>2</sup>2The difference may be traced to the terms in Eq. (A.2) that contain $`B_0(p^2,\sqrt{\xi _A}\lambda ,m_{f,i})`$. If we put $`\xi _A=1`$ in the argument of this function, then Eq. (A.2) coincides with Eqs. (B1)–(B3) of Ref. . This difference becomes relevant when IR-divergent quantities such as $`\delta Z_{ii}^{f,L}`$ are to be calculated, but it is immaterial for the purposes of Ref. . The corresponding formulas in ’t Hooft-Feynman gauge may be found in Appendix A of Ref. . We do not display our analytical results for the form factors $`\delta F_a^\sigma `$ because they are somewhat lengthy. They can be compared with the literature in the limiting cases $`\xi _W=\xi _Z=\xi _A=1`$ or $`m_{f,i}=m_{f^{},j}=0`$ . In the first case, we find agreement with Eqs. (27)–(29) of Ref. .<sup>3</sup><sup>3</sup>3The function in the seventh line of Eq. (29) should carry the superscript “$`\sigma `$” instead of “$``$”. In the second case, only the form factor $`\delta F_1^{}`$ survives, as may be seen from Eq. (3), and we find agreement with Eq. (33) of Ref. . We verified the expressions for $`\delta _\mathrm{b}^{\mathrm{ew}}`$ and $`\delta ^{\mathrm{QCD}}`$ listed in Eqs. (35) and (37) of Ref. , respectively.
At this point, we should comment on a very recent paper in which a new renormalization prescription for the CKM matrix is proposed. The quantity $`T_1`$, defined in Eq. (4) of Ref. , corresponds to our quantity $`_1^{Wf_if_j^{}}`$, defined in Eq. (6) above. The authors of Ref. claim that the finite part of $`_1^{Wf_if_j^{}}`$ becomes gauge dependent if $`\delta V_{ij}`$ is omitted. In order to substantiate this claim, they introduce, in Eq. (23), the auxiliary quantity $`\delta X_{ud}`$, which is to represent the difference between $`_1^{Wf_if_j^{}}`$ and its counterpart for $`V_{ij}=\delta _{ij}`$. In our notation, this quantity reads
$$\delta X_{ij}=\frac{1}{2}V_{ij}\left(\delta Z_{ii}^{f,L}\delta Z_{ii[1]}^{f,L}+\delta Z_{jj}^{f^{},L}\delta Z_{jj[1]}^{f^{},L}\right)+\frac{1}{2}\left(\underset{ki}{}\delta Z_{ik}^{f,L}V_{kj}+\underset{kj}{}V_{ik}\delta Z_{kj}^{f^{},L}\right),$$
(16)
where the subscript “$`[1]`$” indicates that the identification $`V_{ij}=\delta _{ij}`$ is to be made. They find that this quantity is gauge dependent, and propose to define $`\delta V_{ij}=\delta X_{ij}`$. From our above discussion, it is clear that $`\delta V_{ij}`$ must be gauge independent, in order for the renormalized parameters $`V_{ij}`$ to be gauge independent. Otherwise, the latter would not qualify as physical observables. Furthermore, we verified, by inspecting the analytic expressions, that $`_1^{Wf_if_j^{}}+e_1^{}\delta V_{ij}/\left(\sqrt{2}s_w\right)`$ is gauge independent, in accordance with Refs. . This implies that the quantity $`\delta X_{ij}`$, defined in Eq. (16), is also gauge independent, as we explicitly checked. Finally, we remark that the expression for $`T_1`$ given in Eq. (24) of Ref. differs from that given in Eq. (4) ibidem by finite terms because $`\delta g/g`$ and $`\delta Z_W`$ do depend on $`V_{ij}`$.
Equation (13) is formulated in the pure on-shell renormalization scheme, which uses $`\alpha `$ and the physical particle masses as basic parameters. In this scheme, large electroweak corrections arise from fermion loop contributions to the renormalizations of $`\alpha `$ and $`s_w`$. As in any charged-current process, these corrections can be greatly reduced by parameterizing the lowest-order result with Fermi’s coupling constant $`G_F`$ and $`M_W`$ instead of $`\alpha `$ and $`s_w`$. This can be achieved with the aid of the relationship
$$G_F=\frac{\pi \alpha }{\sqrt{2}s_w^2M_W^2}\frac{1}{1\mathrm{\Delta }r},$$
(17)
where $`\mathrm{\Delta }r`$ contains those radiative corrections to the muon decay width which the SM introduces on top of the purely photonic corrections from within Fermi’s model. At one loop, we have
$`\mathrm{\Delta }r`$ $`=`$ $`{\displaystyle \frac{\mathrm{\Pi }^{WW}(0)\stackrel{~}{\text{Re}}\mathrm{\Pi }^{WW}\left(M_W^2\right)}{M_W^2}}+{\displaystyle \frac{c_w^2}{s_w^2}}\stackrel{~}{\text{Re}}\left({\displaystyle \frac{\mathrm{\Pi }^{WW}\left(M_W^2\right)}{M_W^2}}{\displaystyle \frac{\mathrm{\Pi }^{ZZ}\left(M_Z^2\right)}{M_Z^2}}\right)+2{\displaystyle \frac{c_w}{s_w}}{\displaystyle \frac{\mathrm{\Pi }^{ZA}(0)}{M_Z^2}}`$ (18)
$`+{\displaystyle \frac{\mathrm{\Pi }^{AA}(k^2)}{k^2}}|_{k^2=0}+{\displaystyle \frac{\alpha }{4\pi s_w^2}}\left[\left({\displaystyle \frac{7}{2s_w^2}}2\right)\mathrm{ln}c_w^2+6\right].`$
The last term herein represents the vertex and box corrections to the muon decay width in ’t Hooft-Feynman gauge. Thus, the $`\mathrm{\Pi }^{BB^{}}`$ functions in Eq. (18) have to be evaluated in this gauge, too. We recall that Eq. (18) is gauge independent and finite. The quantity $`\mathrm{\Pi }^{AA}(k^2)/k^2|_{k^2=0}`$ receives important contributions from the light-quark flavours, which cannot be reliably predicted in perturbative QCD. This problem is usually circumvented by relating the finite and gauge-independent quantity
$$\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}=\left[\frac{\mathrm{\Pi }^{AA}(k^2)}{k^2}|_{k^2=0}\frac{\mathrm{\Pi }^{AA}\left(M_Z^2\right)}{M_Z^2}\right]_{udscb}$$
(19)
via a subtracted dispersion relation to experimental data on the total cross section of inclusive hadron production in $`e^+e^{}`$ annihilation. In our numerical analysis, we substitute $`\alpha =\sqrt{2}G_Fs_w^2M_W^2/\pi `$ in Eqs. (5) and (14) and, in turn, include the term $`\mathrm{\Delta }r`$, evaluated from Eq. (18), within the parentheses on the right-hand side of Eq. (13). Then, the quantity $`\mathrm{\Pi }^{AA}(k^2)/k^2|_{k^2=0}`$ exactly cancels, so that the theoretical uncertainty in $`\mathrm{\Delta }\alpha _{\mathrm{had}}^{(5)}`$ does not affect our results.
## 3 $`\overline{\mathrm{𝐌𝐒}}`$ Definition of the CKM Matrix
The relationship between the on-shell and $`\overline{\mathrm{𝐌𝐒}}`$ definitions of the CKM matrix may be conveniently revealed by considering the identity
$`𝑽_{𝒊𝒋}^\mathrm{𝟎}`$ $`\mathbf{=}`$ $`𝑽_{𝒊𝒋}\mathbf{+}𝜹𝑽_{𝒊𝒋}`$ (20)
$`\mathbf{=}`$ $`\overline{𝑽}_{𝒊𝒋}\mathbf{+}𝜹\overline{𝑽}_{𝒊𝒋}\mathbf{,}`$
where the superscript “0” labels bare quantities, and $`\overline{\mathrm{𝐌𝐒}}`$ quantities are marked by a bar. By definition, $`𝜹\overline{𝑽}_{𝒊𝒋}`$ is the UV-divergent part of $`𝜹𝑽_{𝒊𝒋}`$, proportional to $`\mathrm{𝟏}\mathbf{/}\mathit{ϵ}\mathbf{+}\mathrm{𝐥𝐧}\mathbf{(}\mathrm{𝟒}𝝅\mathbf{)}\mathbf{}𝜸_𝑬`$, where $`𝜸_𝑬`$ is Euler’s constant. We thus obtain the relationship
$$\overline{𝑽}_{𝒊𝒋}\mathbf{(}𝝁\mathbf{)}\mathbf{=}𝑽_{𝒊𝒋}\mathbf{+}𝚫𝑽_{𝒊𝒋}\mathbf{(}𝝁\mathbf{)}\mathbf{,}$$
(21)
where
$$𝚫𝑽_{𝒊𝒋}\mathbf{(}𝝁\mathbf{)}\mathbf{=}𝜹𝑽_{𝒊𝒋}\mathbf{}𝜹\overline{𝑽}_{𝒊𝒋}$$
(22)
is a finite shift, which depends on the ’t Hooft mass scale $`𝝁`$ of DR. Inserting Eq. (A.2) in Eq. (9) and performing the $`\overline{\mathrm{𝐌𝐒}}`$ subtraction, we find
$`𝚫𝑽_{𝒊𝒋}\mathbf{(}𝝁\mathbf{)}`$ $`\mathbf{=}`$ $`{\displaystyle \underset{𝒌\mathbf{}𝒊}{\mathbf{}}}{\displaystyle \underset{𝒍}{\mathbf{}}}𝑽_{𝒊𝒍}𝑽_{𝒍𝒌}^{\mathbf{}}𝑽_{𝒌𝒋}{\displaystyle \frac{𝒎_{𝒇\mathbf{,}𝒌}^\mathrm{𝟐}\mathbf{+}𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}{𝒎_{𝒇\mathbf{,}𝒌}^\mathrm{𝟐}\mathbf{}𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}}𝒇\mathbf{(}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒍}\mathbf{)}`$ (23)
$`\mathbf{+}{\displaystyle \underset{𝒌\mathbf{}𝒋}{\mathbf{}}}{\displaystyle \underset{𝒍}{\mathbf{}}}𝑽_{𝒊𝒌}𝑽_{𝒌𝒍}^{\mathbf{}}𝑽_{𝒍𝒋}{\displaystyle \frac{𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}^\mathrm{𝟐}\mathbf{+}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒋}^\mathrm{𝟐}}{𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}^\mathrm{𝟐}\mathbf{}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒋}^\mathrm{𝟐}}}𝒇\mathbf{(}𝒎_{𝒇\mathbf{,}𝒍}\mathbf{)}\mathbf{,}`$
where
$$𝒇\mathbf{(}𝒎\mathbf{)}\mathbf{=}\frac{𝜶}{\mathrm{𝟏𝟔}𝝅𝒔_𝒘^\mathrm{𝟐}}\frac{𝒙}{\mathrm{𝟐}}\mathbf{\left[}\mathrm{𝟑}\mathrm{𝐥𝐧}\frac{𝝁^\mathrm{𝟐}}{𝑴_𝑾^\mathrm{𝟐}}\mathbf{}\frac{\mathrm{𝟓}𝒙\mathbf{}\mathrm{𝟏𝟏}}{\mathrm{𝟐}\mathbf{(}𝒙\mathbf{}\mathrm{𝟏}\mathbf{)}}\mathbf{+}\frac{\mathrm{𝟑}𝒙\mathbf{(}𝒙\mathbf{}\mathrm{𝟐}\mathbf{)}}{\mathbf{(}𝒙\mathbf{}\mathrm{𝟏}\mathbf{)}^\mathrm{𝟐}}\mathrm{𝐥𝐧}𝒙\mathbf{\right]}_{𝒙\mathbf{=}\frac{𝒎^\mathrm{𝟐}}{𝑴_𝑾^\mathrm{𝟐}}}\mathbf{.}$$
(24)
The $`𝝁`$ dependence of $`\overline{𝑽}_{𝒊𝒋}\mathbf{(}𝝁\mathbf{)}`$ is described by the $`𝜷`$ function
$$𝜷_{𝑽_{𝒊𝒋}}\mathbf{=}\frac{𝒅}{𝒅\mathrm{𝐥𝐧}𝝁}\overline{𝑽}_{𝒊𝒋}\mathbf{(}𝝁\mathbf{)}\mathbf{.}$$
(25)
Inserting Eq. (21) into Eq. (25), we obtain the one-loop expression
$`𝜷_{𝑽_{𝒊𝒋}}`$ $`\mathbf{=}`$ $`{\displaystyle \frac{\mathrm{𝟑}𝜶}{\mathrm{𝟏𝟔}𝝅𝒔_𝒘^\mathrm{𝟐}𝑴_𝑾^\mathrm{𝟐}}}\mathbf{[}{\displaystyle \underset{𝒌\mathbf{}𝒊}{\mathbf{}}}{\displaystyle \underset{𝒍}{\mathbf{}}}𝑽_{𝒊𝒍}𝑽_{𝒍𝒌}^{\mathbf{}}𝑽_{𝒌𝒋}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒍}^\mathrm{𝟐}{\displaystyle \frac{𝒎_{𝒇\mathbf{,}𝒌}^\mathrm{𝟐}\mathbf{+}𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}{𝒎_{𝒇\mathbf{,}𝒌}^\mathrm{𝟐}\mathbf{}𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}}`$ (26)
$`\mathbf{+}{\displaystyle \underset{𝒌\mathbf{}𝒋}{\mathbf{}}}{\displaystyle \underset{𝒍}{\mathbf{}}}𝑽_{𝒊𝒌}𝑽_{𝒌𝒍}^{\mathbf{}}𝑽_{𝒍𝒋}𝒎_{𝒇\mathbf{,}𝒍}^\mathrm{𝟐}{\displaystyle \frac{𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}^\mathrm{𝟐}\mathbf{+}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒋}^\mathrm{𝟐}}{𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}^\mathrm{𝟐}\mathbf{}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒋}^\mathrm{𝟐}}}\mathbf{]}\mathbf{,}`$
which agrees with the one given in Ref. .
The standard parameterization of $`𝑽_{𝒊𝒋}`$ utilizes three angles, $`𝜽_{\mathrm{𝟏𝟐}}`$, $`𝜽_{\mathrm{𝟐𝟑}}`$, and $`𝜽_{\mathrm{𝟏𝟑}}`$, and one phase, $`𝜹_{\mathrm{𝟏𝟑}}`$ . A popular approximation that emphasizes the hierarchy in the size of the angles, $`\mathrm{𝐬𝐢𝐧}𝜽_{\mathrm{𝟏𝟐}}\mathbf{}\mathrm{𝐬𝐢𝐧}𝜽_{\mathrm{𝟐𝟑}}\mathbf{}\mathrm{𝐬𝐢𝐧}𝜽_{\mathrm{𝟏𝟑}}`$, is due to Wolfenstein , where one sets $`𝝀\mathbf{=}\mathrm{𝐬𝐢𝐧}𝜽_{\mathrm{𝟏𝟐}}`$, the sine of the Cabibbo angle, and then writes the other CKM matrix elements in terms of powers of $`𝝀`$. Through $`𝑶\mathbf{(}𝝀^\mathrm{𝟑}\mathbf{)}`$, one has
$$𝑽\mathbf{=}\mathbf{\left(}\begin{array}{ccc}\mathrm{𝟏}\mathbf{}𝝀^\mathrm{𝟐}\mathbf{/}\mathrm{𝟐}& 𝝀& 𝑨𝝀^\mathrm{𝟑}\mathbf{(}𝝆\mathbf{}𝒊𝜼\mathbf{)}\\ \mathbf{}𝝀& \mathrm{𝟏}\mathbf{}𝝀^\mathrm{𝟐}\mathbf{/}\mathrm{𝟐}& 𝑨𝝀^\mathrm{𝟐}\\ 𝑨𝝀^\mathrm{𝟑}\mathbf{(}\mathrm{𝟏}\mathbf{}𝝆\mathbf{}𝒊𝜼\mathbf{)}& \mathbf{}𝑨𝝀^\mathrm{𝟐}& \mathrm{𝟏}\end{array}\mathbf{\right)}\mathbf{,}$$
(27)
where $`𝑨`$, $`𝝆`$, and $`𝜼`$ are real numbers, which were intended to be of order unity. The relationships between the parameters $`𝝀`$, $`𝑨`$, $`𝝆`$, and $`𝜼`$ in the on-shell scheme and their counterparts in the $`\overline{\mathrm{𝐌𝐒}}`$ scheme may be obtained with the aid of Eq. (21) and read
$`\overline{𝝀}`$ $`\mathbf{=}`$ $`𝝀\mathbf{[}\mathrm{𝟏}\mathbf{+}\mathbf{(}\mathrm{𝟏}\mathbf{}𝝀^\mathrm{𝟐}\mathbf{)}𝑭\mathbf{(}𝒎_𝒖\mathbf{,}𝒎_𝒄\mathbf{,}𝒎_𝒅\mathbf{,}𝒎_𝒔\mathbf{)}\mathbf{]}\mathbf{,}`$
$`\overline{𝑨}`$ $`\mathbf{=}`$ $`𝑨\mathbf{[}\mathrm{𝟏}\mathbf{}\mathrm{𝟐}𝑭\mathbf{(}𝒎_𝒖\mathbf{,}𝒎_𝒄\mathbf{,}𝒎_𝒅\mathbf{,}𝒎_𝒔\mathbf{)}\mathbf{+}𝑭\mathbf{(}𝒎_𝒄\mathbf{,}𝒎_𝒕\mathbf{,}𝒎_𝒔\mathbf{,}𝒎_𝒃\mathbf{)}\mathbf{]}\mathbf{,}`$
$`\overline{𝝆}`$ $`\mathbf{=}`$ $`𝝆\mathbf{[}\mathrm{𝟏}\mathbf{}𝑭\mathbf{(}𝒎_𝒖\mathbf{,}𝒎_𝒄\mathbf{,}𝒎_𝒅\mathbf{,}𝒎_𝒔\mathbf{)}\mathbf{}𝑭\mathbf{(}𝒎_𝒄\mathbf{,}𝒎_𝒕\mathbf{,}𝒎_𝒔\mathbf{,}𝒎_𝒃\mathbf{)}\mathbf{+}𝑭\mathbf{(}𝒎_𝒖\mathbf{,}𝒎_𝒕\mathbf{,}𝒎_𝒅\mathbf{,}𝒎_𝒃\mathbf{)}\mathbf{]}`$
$`\mathbf{+}\mathbf{\left(}{\displaystyle \frac{𝒎_𝒖^\mathrm{𝟐}\mathbf{+}𝒎_𝒄^\mathrm{𝟐}}{𝒎_𝒖^\mathrm{𝟐}\mathbf{}𝒎_𝒄^\mathrm{𝟐}}}\mathbf{}{\displaystyle \frac{𝒎_𝒖^\mathrm{𝟐}\mathbf{+}𝒎_𝒕^\mathrm{𝟐}}{𝒎_𝒖^\mathrm{𝟐}\mathbf{}𝒎_𝒕^\mathrm{𝟐}}}\mathbf{\right)}\mathbf{[}𝒇\mathbf{(}𝒎_𝒅\mathbf{)}\mathbf{}𝒇\mathbf{(}𝒎_𝒔\mathbf{)}\mathbf{]}`$
$`\mathbf{+}\mathbf{\left(}{\displaystyle \frac{𝒎_𝒔^\mathrm{𝟐}\mathbf{+}𝒎_𝒃^\mathrm{𝟐}}{𝒎_𝒔^\mathrm{𝟐}\mathbf{}𝒎_𝒃^\mathrm{𝟐}}}\mathbf{}{\displaystyle \frac{𝒎_𝒅^\mathrm{𝟐}\mathbf{+}𝒎_𝒃^\mathrm{𝟐}}{𝒎_𝒅^\mathrm{𝟐}\mathbf{}𝒎_𝒃^\mathrm{𝟐}}}\mathbf{\right)}\mathbf{[}𝒇\mathbf{(}𝒎_𝒄\mathbf{)}\mathbf{}𝒇\mathbf{(}𝒎_𝒕\mathbf{)}\mathbf{]}\mathbf{,}`$
$`\overline{𝜼}`$ $`\mathbf{=}`$ $`𝜼\mathbf{[}\mathrm{𝟏}\mathbf{}𝑭\mathbf{(}𝒎_𝒖\mathbf{,}𝒎_𝒄\mathbf{,}𝒎_𝒅\mathbf{,}𝒎_𝒔\mathbf{)}\mathbf{}𝑭\mathbf{(}𝒎_𝒄\mathbf{,}𝒎_𝒕\mathbf{,}𝒎_𝒔\mathbf{,}𝒎_𝒃\mathbf{)}\mathbf{+}𝑭\mathbf{(}𝒎_𝒖\mathbf{,}𝒎_𝒕\mathbf{,}𝒎_𝒅\mathbf{,}𝒎_𝒃\mathbf{)}\mathbf{]}\mathbf{,}`$ (28)
where
$$𝑭\mathbf{(}𝒎_\mathrm{𝟏}\mathbf{,}𝒎_\mathrm{𝟐}\mathbf{,}𝒎_\mathrm{𝟑}\mathbf{,}𝒎_\mathrm{𝟒}\mathbf{)}\mathbf{=}\mathbf{\left\{}\frac{𝒎_\mathrm{𝟏}^\mathrm{𝟐}\mathbf{+}𝒎_\mathrm{𝟐}^\mathrm{𝟐}}{𝒎_\mathrm{𝟏}^\mathrm{𝟐}\mathbf{}𝒎_\mathrm{𝟐}^\mathrm{𝟐}}\mathbf{[}𝒇\mathbf{(}𝒎_\mathrm{𝟑}\mathbf{)}\mathbf{}𝒇\mathbf{(}𝒎_\mathrm{𝟒}\mathbf{)}\mathbf{]}\mathbf{+}\frac{𝒎_\mathrm{𝟑}^\mathrm{𝟐}\mathbf{+}𝒎_\mathrm{𝟒}^\mathrm{𝟐}}{𝒎_\mathrm{𝟑}^\mathrm{𝟐}\mathbf{}𝒎_\mathrm{𝟒}^\mathrm{𝟐}}\mathbf{[}𝒇\mathbf{(}𝒎_\mathrm{𝟏}\mathbf{)}\mathbf{}𝒇\mathbf{(}𝒎_\mathrm{𝟐}\mathbf{)}\mathbf{]}\mathbf{\right\}}\mathbf{.}$$
(29)
Exploiting the strong hierarchy among the up- and down-type quark masses, which satisfy $`𝒎_𝒖\mathbf{}𝒎_𝒄\mathbf{}𝒎_𝒕`$ and $`𝒎_𝒅\mathbf{}𝒎_𝒔\mathbf{}𝒎_𝒃`$, respectively, and the significant mass splittings within the second and third quark generations, $`𝒎_𝒔\mathbf{}𝒎_𝒄`$ and $`𝒎_𝒃\mathbf{}𝒎_𝒕`$, respectively, we can derive the following approximation formulas:
$`\overline{𝝀}`$ $`\mathbf{=}`$ $`𝝀\mathbf{\left[}\mathrm{𝟏}\mathbf{+}\mathbf{(}\mathrm{𝟏}\mathbf{}𝝀^\mathrm{𝟐}\mathbf{)}𝒇\mathbf{(}𝒎_𝒄\mathbf{)}\mathbf{\right]}\mathbf{,}`$
$`\overline{𝑨}`$ $`\mathbf{=}`$ $`𝑨\mathbf{\left[}\mathrm{𝟏}\mathbf{+}𝒇\mathbf{(}𝒎_𝒕\mathbf{)}\mathbf{\right]}\mathbf{,}`$
$`\overline{𝝆}`$ $`\mathbf{=}`$ $`𝝆\mathbf{,}`$
$`\overline{𝜼}`$ $`\mathbf{=}`$ $`𝜼\mathbf{.}`$ (30)
For $`𝝁\mathbf{=}𝑴_𝑾`$, these approximations agree with the exact results, evaluated from Eq. (28), within an error of less than $`\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟓}`$.
## 4 Numerical Results
We now describe our numerical analysis of Eq. (13), with the modifications specified at the end of Sec. 2. We evaluated the $`𝑨_\mathrm{𝟎}`$, $`𝑩_\mathrm{𝟎}`$, and $`𝑪_\mathrm{𝟎}`$ functions with the aid of the Fortran program package FF , which is embedded into the Mathematica environment through the program package LoopTools . We performed several checks for the correctness and the stability of our numerical results. As we demonstrated in Sec. 2, Eq. (13) is IR and UV finite and gauge independent, as a consequence of cancellations among the various terms contained in Eqs. (14) and (15). We can check the IR finiteness and gauge independence numerically by varying the IR regulator $`𝝀`$ and the gauge parameters $`𝝃_𝑾`$, $`𝝃_𝒁`$, and $`𝝃_𝑨`$, respectively. In the physical limit $`𝑫\mathbf{}\mathrm{𝟒}`$, the UV divergences appear as terms proportional to $`\mathrm{𝟏}\mathbf{/}\mathit{ϵ}\mathbf{+}\mathrm{𝐥𝐧}\mathbf{(}\mathrm{𝟒}𝝅\mathbf{)}\mathbf{}𝜸_𝑬`$, which are accompanied by a term $`\mathrm{𝐥𝐧}\mathbf{(}𝝁^\mathrm{𝟐}\mathbf{/}𝑴^\mathrm{𝟐}\mathbf{)}`$, with $`𝑴`$ being a characteristic mass scale of the considered loop integral. Although we nullified $`\mathrm{𝟏}\mathbf{/}\mathit{ϵ}\mathbf{+}\mathrm{𝐥𝐧}\mathbf{(}\mathrm{𝟒}𝝅\mathbf{)}\mathbf{}𝜸_𝑬`$ in our computer program, we can check the UV finiteness numerically by varying the ’t Hooft mass scale $`𝝁`$. A further check on the stability of our numerical analysis can be obtained by directly evaluating the two- and three-point tensor integrals with the aid of LoopTools instead of applying the Passarino-Veltman reduction algorithm . Our numerical analysis passed all these checks. Finally, we managed to reproduce the numerical results of Ref. after adopting the definition of $`𝜹𝑽_{𝒊𝒋}`$, the choice of gauge, and the values of the input parameters from there.
We use the following input parameters :
$`\begin{array}{ccc}𝜶\mathbf{=}\mathrm{𝟏}\mathbf{/}\mathbf{137.03599976}\mathbf{,}\hfill & 𝑮_𝑭\mathbf{=}\mathbf{1.16639}\mathbf{\times }\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟓}\text{GeV}^\mathbf{}\mathrm{𝟐}\mathbf{,}\hfill & 𝜶_𝒔^{\mathbf{(}\mathrm{𝟓}\mathbf{)}}\mathbf{(}𝑴_𝒁\mathbf{)}\mathbf{=}\mathbf{0.1181}\mathbf{,}\hfill \\ 𝑴_𝑾\mathbf{=}\mathbf{80.419}\text{GeV}\mathbf{,}\hfill & 𝑴_𝒁\mathbf{=}\mathbf{91.1871}\text{GeV}\mathbf{,}\hfill & \\ 𝒎_𝒆\mathbf{=}\mathbf{0.510998902}\text{MeV}\mathbf{,}\hfill & 𝒎_𝝁\mathbf{=}\mathbf{105.658357}\text{MeV}\mathbf{,}\hfill & 𝒎_𝝉\mathbf{=}\mathbf{1777.03}\text{MeV}\mathbf{,}\hfill \\ 𝒎_𝒖\mathbf{=}\mathrm{𝟑}\text{MeV}\mathbf{,}\hfill & 𝒎_𝒅\mathbf{=}\mathrm{𝟔}\text{MeV}\mathbf{,}\hfill & 𝒎_𝒔\mathbf{=}\mathrm{𝟏𝟐𝟑}\text{MeV}\mathbf{,}\hfill \\ 𝒎_𝒄\mathbf{=}\mathbf{1.5}\text{GeV}\mathbf{,}\hfill & 𝒎_𝒃\mathbf{=}\mathbf{4.6}\text{GeV}\mathbf{,}\hfill & 𝒎_𝒕\mathbf{=}\mathbf{174.3}\text{GeV}\mathbf{,}\hfill \\ 𝒔_{\mathrm{𝟏𝟐}}\mathbf{=}\mathbf{0.223}\mathbf{,}\hfill & 𝒔_{\mathrm{𝟐𝟑}}\mathbf{=}\mathbf{0.040}\mathbf{,}\hfill & 𝒔_{\mathrm{𝟏𝟑}}\mathbf{=}\mathbf{0.004}\mathbf{.}\hfill \end{array}`$ (37)
Here, $`𝒎_𝒖`$, $`𝒎_𝒅`$, and $`𝒎_𝒔`$ correspond to current-quark masses, and $`𝒎_𝒄`$, $`𝒎_𝒃`$, and $`𝒎_𝒕`$ to pole masses. Since the contributions from the light quarks $`𝒒\mathbf{=}𝒖\mathbf{,}𝒅\mathbf{,}𝒔\mathbf{,}𝒄\mathbf{,}𝒃`$ come with the suppression factors $`𝒎_𝒒^\mathrm{𝟐}\mathbf{/}𝑴_𝑾^\mathrm{𝟐}`$, the considerable uncertainties in the values of $`𝒎_𝒒`$ do not jeopardize the reliability of our theoretical predictions. We take the neutrinos to be massless. For simplicity, we assume that $`𝜹_{\mathrm{𝟏𝟑}}\mathbf{=}\mathrm{𝟎}`$. Then, the values for $`𝒔_{𝒊𝒋}\mathbf{=}\mathrm{𝐬𝐢𝐧}𝜽_{𝒊𝒋}`$ provided in the last row of Eq. (37) lead to
$`\begin{array}{ccc}𝑽_{𝒖𝒅}\mathbf{=}\mathbf{0.975}\mathbf{,}\hfill & 𝑽_{𝒖𝒔}\mathbf{=}\mathbf{0.223}\mathbf{,}\hfill & 𝑽_{𝒖𝒃}\mathbf{=}\mathbf{0.004}\mathbf{,}\hfill \\ 𝑽_{𝒄𝒅}\mathbf{=}\mathbf{}\mathbf{0.223}\mathbf{,}\hfill & 𝑽_{𝒄𝒔}\mathbf{=}\mathbf{0.974}\mathbf{,}\hfill & 𝑽_{𝒄𝒃}\mathbf{=}\mathbf{0.040}\mathbf{,}\hfill \\ 𝑽_{𝒕𝒅}\mathbf{=}\mathbf{0.005}\mathbf{,}\hfill & 𝑽_{𝒕𝒔}\mathbf{=}\mathbf{}\mathbf{0.040}\mathbf{,}\hfill & 𝑽_{𝒕𝒃}\mathbf{=}\mathbf{0.999}\mathbf{.}\hfill \end{array}`$ (41)
These values approximately satisfy the unitarity condition $`𝑽_{𝒊𝒌}𝑽_{𝒌𝒋}^{\mathbf{}}\mathbf{=}𝜹_{𝒊𝒋}`$.<sup>4</sup><sup>4</sup>4Notice that the indices of $`V_{ij}`$ refer to generations rather than quark flavours. For example, we have $`V_{12}=V_{us}`$, $`V_{21}=V_{cd}`$, and $`V_{12}^{}=V_{21}^{}=V_{cd}^{}`$. We evaluate $`𝜶_𝒔^{\mathbf{(}𝒏_𝒇\mathbf{)}}\mathbf{(}𝝁\mathbf{)}`$ appearing in $`𝜹^{\mathrm{𝐐𝐂𝐃}}`$ at the renormalization scale $`𝝁\mathbf{=}𝑴_𝑾`$ with $`𝒏_𝒇\mathbf{=}\mathrm{𝟓}`$ active quark flavours from the one-loop relation
$$𝜶_𝒔^{\mathbf{(}𝒏_𝒇\mathbf{)}}\mathbf{(}𝑴_𝑾\mathbf{)}\mathbf{=}\frac{𝜶_𝒔^{\mathbf{(}𝒏_𝒇\mathbf{)}}\mathbf{(}𝑴_𝒁\mathbf{)}}{\mathrm{𝟏}\mathbf{+}𝜶_𝒔^{\mathbf{(}𝒏_𝒇\mathbf{)}}\mathbf{(}𝑴_𝒁\mathbf{)}𝜷_\mathrm{𝟎}\mathrm{𝐥𝐧}\mathbf{(}𝑴_𝑾^\mathrm{𝟐}\mathbf{/}𝑴_𝒁^\mathrm{𝟐}\mathbf{)}\mathbf{/}𝝅}\mathbf{,}$$
(42)
where $`𝜷_\mathrm{𝟎}\mathbf{=}\mathrm{𝟏𝟏}\mathbf{/}\mathrm{𝟒}\mathbf{}𝒏_𝒇\mathbf{/}\mathrm{𝟔}`$. For the Higgs-boson mass, we consider the values $`𝑴_𝑯\mathbf{=}\mathrm{𝟏𝟎𝟎}`$, 250, and 600 GeV.
We now present our numerical results. We first investigate the quantitative significance of the definition of $`𝜹𝑽_{𝒊𝒋}`$. Toward this end, we compare, in Table 1, our results for the partial widths of the various hadronic $`𝑾`$-boson decay channels with those obtained in ’t Hooft-Feynman gauge with the definition of $`𝜹𝑽_{𝒊𝒋}`$ proposed in Ref. , assuming $`𝑴_𝑯\mathbf{=}\mathrm{𝟐𝟓𝟎}`$ GeV. The relative deviations are largest for the final states involving the $`𝒃`$ quark, where they are of order $`𝜶𝒎_𝒃^\mathrm{𝟐}\mathbf{/}\mathbf{(}𝝅𝑴_𝑾^\mathrm{𝟐}\mathbf{)}\mathbf{}\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟓}`$. Although small against the present experimental accuracies , they are of the same order as the entire shifts due to the renormalization of the CKM matrix . We stress that the numbers in the second column of Table 1 do depend on the choice of gauge. However, this gauge dependence turns out to be feeble. In Table 2, we present our tree-level and one-loop results for the leptonic and hadronic partial decay widths of the $`𝑾`$ boson, assuming $`𝑴_𝑯\mathbf{=}\mathrm{𝟏𝟎𝟎}`$, 250, or 600 GeV. In the leptonic channels, the radiative corrections are nearly flavour independent and amount to approximately $`\mathbf{}\mathbf{0.3}\mathbf{\%}`$. In the hadronic channels, the corrections range between 3.5% and 3.8% and are dominantly of QCD origin. In all cases, the $`𝑴_𝑯`$ dependence is feeble, of relative order $`\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟓}`$. Finally, we determine the uncertainties in our theoretical prediction for the total $`𝑾`$-boson decay width $`𝚪_𝑾`$ due to the errors on our input parameters. Specifically, the variations of $`𝑮_𝑭`$, $`𝜶_𝒔^{\mathbf{(}\mathrm{𝟓}\mathbf{)}}\mathbf{(}𝑴_𝒁\mathbf{)}`$, $`𝑴_𝑾`$, $`𝑴_𝒁`$, $`𝒎_𝒄`$, $`𝒎_𝒃`$, $`𝒎_𝒕`$, $`𝒔_{\mathrm{𝟏𝟐}}`$, $`𝒔_{\mathrm{𝟐𝟑}}`$, and $`𝒔_{\mathrm{𝟏𝟑}}`$ by $`\mathbf{\pm }\mathrm{𝟏}\mathbf{\times }\mathrm{𝟏𝟎}^{\mathbf{}\mathrm{𝟏𝟎}}\text{GeV}^\mathbf{}\mathrm{𝟐}`$, $`\mathbf{\pm }\mathbf{0.002}`$, $`\mathbf{\pm }\mathrm{𝟑𝟖}`$ MeV, $`\mathbf{\pm }\mathbf{2.1}`$ MeV, $`\mathbf{\pm }\mathbf{0.1}`$ GeV, $`\mathbf{\pm }\mathbf{0.2}`$ GeV, $`\mathbf{\pm }\mathbf{5.1}`$ GeV, $`\mathbf{\pm }\mathbf{0.004}`$, $`\mathbf{\pm }\mathbf{0.003}`$, and $`\mathbf{\pm }\mathbf{0.002}`$ shift $`𝚪_𝑾`$ by $`\mathbf{\pm }\mathbf{0.02}`$, $`\mathbf{\pm }\mathbf{0.9}`$, $`\mathbf{\pm }\mathbf{3.0}`$, $`\mathbf{\pm }\mathbf{7.5}\mathbf{\times }\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟒}`$, $`\mathbf{\pm }\mathbf{0.02}`$, $`\mathbf{\pm }\mathrm{𝟏}\mathbf{\times }\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟒}`$, $`\mathbf{\pm }\mathbf{0.033}`$, $`\mathbf{\pm }\mathrm{𝟑}\mathbf{\times }\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟓}`$, $`\mathbf{\pm }\mathrm{𝟑}\mathbf{\times }\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟒}`$, and $`\mathbf{\pm }\mathrm{𝟑}\mathbf{\times }\mathrm{𝟏𝟎}^\mathbf{}\mathrm{𝟓}`$ MeV, respectively. The residual parametric uncertainties are marginal.
## 5 Conclusions
We calculated the partial decay widths of the $`𝑾`$ boson at one loop in the SM using the on-shell scheme endowed with the gauge-independent definition of the CKM matrix $`𝑽_{𝒊𝒋}`$ recently proposed in Refs. . Working in $`𝑹_𝝃`$ gauge, we explicitly verified that the final expressions are independent of the gauge parameters. In particular, the renormalization constant $`𝜹𝑽_{𝒊𝒋}\mathbf{/}𝑽_{𝒊𝒋}`$ (9) and the one-loop amplitude $`𝓜_\mathrm{𝟏}^{𝑾𝒇_𝒊𝒇_𝒋^{\mathbf{}}}`$ (6) of the $`𝑾`$-boson decay to quarks (1) with this renormalization constant removed are separately gauge independent. In this respect, we disagree with the findings of Ref. . The difference between our analysis and the corresponding one with the gauge-dependent definition of $`𝜹𝑽_{𝒊𝒋}\mathbf{/}𝑽_{𝒊𝒋}`$ from Ref. is of the same order as the entire effect due to the renormalization of the CKM matrix, but it is small compared to the present experimental precision. Furthermore, we established the relationship between the on-shell and $`\overline{\mathrm{𝐌𝐒}}`$ definitions of the CKM matrix, both in its generic form and in the Wolfenstein parameterization . As a by-product of our analysis, we recovered the beta function of the CKM matrix .
Note added
In the meantime, a revised version of Ref. has appeared, in which the weaknesses pointed out in Sec. 2 have been remedied.
Acknowledgements
We are grateful to K.-P. O. Diener for checking our results for the triangle diagrams and for verifying the expressions for $`𝜹_𝐛^{\mathrm{𝐞𝐰}}`$ and $`𝜹^{\mathrm{𝐐𝐂𝐃}}`$ listed in Eqs. (35) and (37) of Ref. , respectively. We enjoyed fruitful discussions with P. A. Grassi. M. S. would like to thank T. Hahn and G. Weiglein for valuable advice concerning FeynArts and FeynCalc. This work was supported in part by the Deutsche Forschungsgemeinschaft through Grant No. KN 365/1-1, by the Bundesministerium für Bildung und Forschung through Grant No. 05 HT9GUA 3, and by the European Commission through the Research Training Network Quantum Chromodynamics and the Deep Structure of Elementary Particles under Contract No. ERBFMRX-CT98-0194.
## Appendix A Fermion Two-Point Functions
The unrenormalized two-point function describing the fermion transition $`𝒇_𝒋\mathbf{}𝒇_𝒊`$, at four-momentum $`𝒑`$, may be decomposed as
$$𝚪_{𝒊𝒋}^𝒇\mathbf{(}𝒑\mathbf{)}\mathbf{=}𝒊\mathbf{\left[}𝜹_{𝒊𝒋}\mathbf{\left(}\overline{)}𝒑\mathbf{}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{\right)}\mathbf{+}\overline{)}𝒑𝝎_{\mathbf{}}𝚺_{𝒊𝒋}^{𝒇\mathbf{,}𝑳}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{)}\mathbf{+}\overline{)}𝒑𝝎_\mathbf{+}𝚺_{𝒊𝒋}^{𝒇\mathbf{,}𝑹}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{)}\mathbf{+}\mathbf{\left(}𝒎_{𝒇\mathbf{,}𝒊}𝝎_{\mathbf{}}\mathbf{+}𝒎_{𝒇\mathbf{,}𝒋}𝝎_\mathbf{+}\mathbf{\right)}𝚺_{𝒊𝒋}^{𝒇\mathbf{,}𝑺}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{)}\mathbf{\right]}\mathbf{.}$$
(A.1)
At one loop in $`𝑹_𝝃`$ gauge, we find for the coefficient functions herein
$`𝚺_{𝒊𝒋}^{𝒇\mathbf{,}𝑳}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{}{\displaystyle \frac{𝜶}{\mathrm{𝟒}𝝅}}\mathbf{\{}𝜹_{𝒊𝒋}𝑸_𝒇^\mathrm{𝟐}\mathbf{[}\mathrm{𝟐}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝝀\mathbf{)}\mathbf{+}\mathrm{𝟏}\mathbf{+}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝝀\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{}𝝃_𝑨𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}\sqrt{𝝃_𝑨}𝝀\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}`$
$`\mathbf{+}{\displaystyle \frac{𝒑^\mathrm{𝟐}\mathbf{}𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}{𝝀^\mathrm{𝟐}}}\mathbf{(}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝝀\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}\sqrt{𝝃_𝑨}𝝀\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{)}\mathbf{]}`$
$`\mathbf{+}𝜹_{𝒊𝒋}\mathbf{\left(}𝒈_𝒇^{\mathbf{}}\mathbf{\right)}^\mathrm{𝟐}\mathbf{[}\mathrm{𝟐}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝑴_𝒁\mathbf{)}\mathbf{+}\mathrm{𝟏}`$
$`\mathbf{+}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝑴_𝒁\mathbf{)}\mathbf{}𝝃_𝒁𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}\sqrt{𝝃_𝒁}𝑴_𝒁\mathbf{)}`$
$`\mathbf{+}{\displaystyle \frac{𝒑^\mathrm{𝟐}\mathbf{}𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}{𝑴_𝒁^\mathrm{𝟐}}}\mathbf{(}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝑴_𝒁\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}\sqrt{𝝃_𝒁}𝑴_𝒁\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{)}\mathbf{]}`$
$`\mathbf{+}𝜹_{𝒊𝒋}{\displaystyle \frac{𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}{\mathrm{𝟒}𝒔_𝒘^\mathrm{𝟐}𝑴_𝑾^\mathrm{𝟐}}}\mathbf{\left[}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}\sqrt{𝝃_𝒁}𝑴_𝒁\mathbf{)}\mathbf{+}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝑴_𝑯\mathbf{)}\mathbf{\right]}`$
$`\mathbf{+}{\displaystyle \frac{\mathrm{𝟏}}{\mathrm{𝟐}𝒔_𝒘^\mathrm{𝟐}}}{\displaystyle \underset{𝒌}{\mathbf{}}}𝑽_{𝒊𝒌}𝑽_{𝒌𝒋}^{\mathbf{}}\mathbf{[}\mathrm{𝟐}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}\mathbf{,}𝑴_𝑾\mathbf{)}\mathbf{+}{\displaystyle \frac{𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}^\mathrm{𝟐}}{𝑴_𝑾^\mathrm{𝟐}}}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}\mathbf{,}\sqrt{𝝃_𝑾}𝑴_𝑾\mathbf{)}`$
$`\mathbf{+}\mathrm{𝟏}\mathbf{+}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}\mathbf{,}𝑴_𝑾\mathbf{)}\mathbf{}𝝃_𝑾𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}\mathbf{,}\sqrt{𝝃_𝑾}𝑴_𝑾\mathbf{)}`$
$`\mathbf{+}{\displaystyle \frac{𝒑^\mathrm{𝟐}\mathbf{}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}^\mathrm{𝟐}}{𝑴_𝑾^\mathrm{𝟐}}}\mathbf{[}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝑴_𝑾\mathbf{,}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}\mathbf{)}\mathbf{}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}\sqrt{𝝃_𝑾}𝑴_𝑾\mathbf{,}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}\mathbf{)}\mathbf{]}\mathbf{\}}\mathbf{,}`$
$`𝚺_{𝒊𝒋}^{𝒇\mathbf{,}𝑹}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{}{\displaystyle \frac{𝜶}{\mathrm{𝟒}𝝅}}\mathbf{\{}𝜹_{𝒊𝒋}𝑸_𝒇^\mathrm{𝟐}\mathbf{[}\mathrm{𝟐}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝝀\mathbf{)}\mathbf{+}\mathrm{𝟏}\mathbf{+}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝝀\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{}𝝃_𝑨𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}\sqrt{𝝃_𝑨}𝝀\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}`$
$`\mathbf{+}{\displaystyle \frac{𝒑^\mathrm{𝟐}\mathbf{}𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}{𝝀^\mathrm{𝟐}}}\mathbf{(}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝝀\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}\sqrt{𝝃_𝑨}𝝀\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{)}\mathbf{]}`$
$`\mathbf{+}𝜹_{𝒊𝒋}\mathbf{\left(}𝒈_𝒇^\mathbf{+}\mathbf{\right)}^\mathrm{𝟐}\mathbf{[}\mathrm{𝟐}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝑴_𝒁\mathbf{)}\mathbf{+}\mathrm{𝟏}`$
$`\mathbf{+}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝑴_𝒁\mathbf{)}\mathbf{}𝝃_𝒁𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}\sqrt{𝝃_𝒁}𝑴_𝒁\mathbf{)}`$
$`\mathbf{+}{\displaystyle \frac{𝒑^\mathrm{𝟐}\mathbf{}𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}{𝑴_𝒁^\mathrm{𝟐}}}\mathbf{(}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝑴_𝒁\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}\sqrt{𝝃_𝒁}𝑴_𝒁\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{)}\mathbf{)}\mathbf{]}`$
$`\mathbf{+}𝜹_{𝒊𝒋}{\displaystyle \frac{𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}{\mathrm{𝟒}𝒔_𝒘^\mathrm{𝟐}𝑴_𝑾^\mathrm{𝟐}}}\mathbf{\left[}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}\sqrt{𝝃_𝒁}𝑴_𝒁\mathbf{)}\mathbf{+}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝑴_𝑯\mathbf{)}\mathbf{\right]}`$
$`\mathbf{+}{\displaystyle \frac{𝒎_{𝒇\mathbf{,}𝒊}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒋}}{\mathrm{𝟐}𝒔_𝒘^\mathrm{𝟐}𝑴_𝑾^\mathrm{𝟐}}}{\displaystyle \underset{𝒌}{\mathbf{}}}𝑽_{𝒊𝒌}𝑽_{𝒌𝒋}^{\mathbf{}}𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}\mathbf{,}\sqrt{𝝃_𝑾}𝑴_𝑾\mathbf{)}\mathbf{\}}\mathbf{,}`$
$`𝚺_{𝒊𝒋}^{𝒇\mathbf{,}𝑺}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{)}`$ $`\mathbf{=}`$ $`\mathbf{}{\displaystyle \frac{𝜶}{\mathrm{𝟒}𝝅}}\mathbf{\{}𝜹_{𝒊𝒋}𝑸_𝒇^\mathrm{𝟐}\mathbf{[}\mathrm{𝟒}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝝀\mathbf{)}\mathbf{}\mathrm{𝟐}\mathbf{}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝝀\mathbf{)}\mathbf{+}𝝃_𝑨𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}\sqrt{𝝃_𝑨}𝝀\mathbf{)}\mathbf{]}`$ (A.2)
$`\mathbf{+}𝜹_{𝒊𝒋}𝒈_𝒇^\mathbf{+}𝒈_𝒇^{\mathbf{}}\mathbf{[}\mathrm{𝟒}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝑴_𝒁\mathbf{)}\mathbf{}\mathrm{𝟐}`$
$`\mathbf{}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝑴_𝒁\mathbf{)}\mathbf{+}𝝃_𝒁𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}\sqrt{𝝃_𝒁}𝑴_𝒁\mathbf{)}\mathbf{]}`$
$`\mathbf{+}𝜹_{𝒊𝒋}{\displaystyle \frac{𝒎_{𝒇\mathbf{,}𝒊}^\mathrm{𝟐}}{\mathrm{𝟒}𝒔_𝒘^\mathrm{𝟐}𝑴_𝑾^\mathrm{𝟐}}}\mathbf{\left[}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}\sqrt{𝝃_𝒁}𝑴_𝒁\mathbf{)}\mathbf{}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇\mathbf{,}𝒊}\mathbf{,}𝑴_𝑯\mathbf{)}\mathbf{\right]}`$
$`\mathbf{+}{\displaystyle \frac{\mathrm{𝟏}}{\mathrm{𝟐}𝒔_𝒘^\mathrm{𝟐}𝑴_𝑾^\mathrm{𝟐}}}{\displaystyle \underset{𝒌}{\mathbf{}}}𝑽_{𝒊𝒌}𝑽_{𝒌𝒋}^{\mathbf{}}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}^\mathrm{𝟐}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_{𝒇^{\mathbf{}}\mathbf{,}𝒌}\mathbf{,}\sqrt{𝝃_𝑾}𝑴_𝑾\mathbf{)}\mathbf{\}}\mathbf{,}`$
where
$`𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_\mathrm{𝟎}\mathbf{,}𝒎_\mathrm{𝟏}\mathbf{)}`$ $`\mathbf{=}`$ $`{\displaystyle \frac{\mathbf{(}\mathrm{𝟐}𝝅𝝁\mathbf{)}^{\mathrm{𝟒}\mathbf{}𝑫}}{𝒊𝝅^\mathrm{𝟐}}}{\displaystyle \mathbf{}\frac{𝒅^𝑫𝒒}{\mathbf{\left(}𝒒^\mathrm{𝟐}\mathbf{}𝒎_\mathrm{𝟎}^\mathrm{𝟐}\mathbf{\right)}\mathbf{\left[}\mathbf{(}𝒒\mathbf{+}𝒑\mathbf{)}^\mathrm{𝟐}\mathbf{}𝒎_\mathrm{𝟏}^\mathrm{𝟐}\mathbf{\right]}}}\mathbf{,}`$ (A.3)
$`𝑩_\mathrm{𝟏}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_\mathrm{𝟎}\mathbf{,}𝒎_\mathrm{𝟏}\mathbf{)}`$ $`\mathbf{=}`$ $`{\displaystyle \frac{𝒎_\mathrm{𝟏}^\mathrm{𝟐}\mathbf{}𝒎_\mathrm{𝟎}^\mathrm{𝟐}}{\mathrm{𝟐}𝒑^\mathrm{𝟐}}}\mathbf{\left[}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_\mathrm{𝟎}\mathbf{,}𝒎_\mathrm{𝟏}\mathbf{)}\mathbf{}𝑩_\mathrm{𝟎}\mathbf{(}\mathrm{𝟎}\mathbf{,}𝒎_\mathrm{𝟎}\mathbf{,}𝒎_\mathrm{𝟏}\mathbf{)}\mathbf{\right]}\mathbf{}{\displaystyle \frac{\mathrm{𝟏}}{\mathrm{𝟐}}}𝑩_\mathrm{𝟎}\mathbf{(}𝒑^\mathrm{𝟐}\mathbf{,}𝒎_\mathrm{𝟎}\mathbf{,}𝒎_\mathrm{𝟏}\mathbf{)}\mathbf{.}`$
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# Hydrodynamics beyond local equilibrium: application to electron gas
## I Introduction
The idea to apply a macroscopic hydrodynamical description to the collective dynamics of the inhomogeneous many-electron systems has been first suggested in a famous work by Bloch in 1933. In this paper he introduced a hydrodynamic theory of a degenerate Fermi gas as a simplest phenomenological extension of the Thomas-Fermi model to dynamical regime. Since application of nonequilibrium many-body theory to inhomogeneous systems is extremely complex, the Bloch’s hydrodynamic theory (BHT) remains popular and is in use until now. Since 1933 the BHT has been applied to variety of kinetic problems, albeit due to the heuristic nature of the BHT, its relation to the kinetic theory and the range of applicability remained unclear.
A common sense of hydrodynamical theory presumes an existence of a local statistical equilibrium in every point of space - the condition which allows to express the kinetic properties of a many-body system in simple terms of pressure, density and temperature distributions. This condition is, however, rarely met in electron gas where the electron-electron collisions commonly play a minor role. This circumstance is ignored in the BHT. Although the applications of the Bloch’s hydrodynamical theory have shown its usefulness, its failures are also well known, in particular a prediction of a wrong plasmon dispersion.
Is it possible to derive a macroscopic hydrodynamics-like theory which would be valid beyond the condition of a local equilibrium? What is the role of collisions and how is the static (i.e. collision dominated) limit as well as the high frequency (collisionless) limit recovered? These are the questions we address in the present paper. In the short publication we showed that such hydrodynamics existed. In the paper Ref., for the sake of simplicity we carried out the derivation of the hydrodynamics equations only in one dimensional case. We also restricted the theory to the second order in the long-wave limit of a kinetic equation, which left out the thermal conductivity and consequently the thermal equilibration processes.
In the present paper we give a new general derivation and extend the theory to include the heat transport. We show that the generalized hydrodynamics can be based on Landau theory of Fermi liquid and relate macroscopic parameters in hydrodynamics equations (i.e. shear and bulk moduli), to microscopic Landau parameters. We also present transparent analitical examples of the calculation of collective modes in confined systems of different dimensionality.
Let us first consider the mathematical structure of BHT and the formal microscopic restrictions of its applicability. The Bloch’s original idea was to use the set of hydrodynamics equations of an ideal charged liquid to describe the dynamics of a Fermi gas. Only the macroscopic variables - electron density $`n(𝐫,t)`$, velocity $`𝐯(𝐫,t)`$, pressure $`P`$, and electrostatic potential $`\phi (𝐫,t),`$ enter the equations:
$`D_tn`$ $`+`$ $`n𝐯=0,`$ (1)
$`mnD_t𝐯`$ $`+`$ $`Pen\phi =0,`$ (2)
$`^2\phi `$ $`=`$ $`4\pi en4\pi \rho _{ext},`$ (3)
where $`\rho _{ext}`$ is an external charge and $`D_t=_t+𝐯`$ is a substantial derivative. This set of equations (continuity equation (1), Euler (2) and Poisson (3) equations) becomes complete when the equation of state is added. In the original paper Bloch identified $`P`$ with the kinetic pressure of a degenerate Fermi gas:
$$P=\frac{1}{5m}(3\pi ^2)^{2/3}n^{5/3}$$
(4)
Further improvements of BHT addressed only the equation of state (inclusion of exchange, correlation and quantum gradient corrections).
Sometimes BHT is viewed as an approximate extension of the density functional theory to the dynamic regime. However, as we have mentioned above, the hydrodynamics theory (1)-(4) exhibits an inconsistency which has been reflected in many papers and textbooks. The system of equations (1)-(4) gives the wrong velocity coefficient $`v_0^2`$ in the plasmon dispersion law $`\omega ^2=\omega _p^2+v_0^2q^2`$. Instead of the correct result $`\frac{3}{5}v_F^2`$ ($`v_F`$ is the Fermi velocity) of the linear response theory (RPA) the BHT leads to the value $`v_0^2=\frac{1}{3}v_F^2`$. This inconsistency does not depend on a degree of degeneracy of an electron gas and is the general property of the hydrodynamics theory which is based on the equilibrium equation of state. For example, at arbitrary degeneracy the result of hydrodynamics is $`v_0^2=v_s^2`$, where $`v_s`$ is the velocity of sound, whereas in the linear response theory $`v_0^2`$ equals to the mean square of the particle velocity $`<v_p^2>`$.
It has been realized that this discrepancy originates from the assumption of a local equilibrium, which underlies the common hydrodynamic theory. The assumption allows to reduce the kinetic equation for the distribution function $`f_𝐩(𝐫,t)`$
$$_tf_𝐩+\frac{𝐩}{m}f_𝐩+e\phi \frac{f_𝐩}{𝐩}=I_𝐩\left[f_𝐩\right]$$
(5)
to a set of equations for macroscopic variables $`n(𝐫,t)`$,$`𝐯(𝐫,t)`$ and, in general case, temperature $`T(𝐫,t)`$. The requirement of the local equilibrium is fulfilled if the characteristic time of the process $`t_{pr}1/\omega `$ is much longer than the inverse collision frequency $`1/\nu _c`$ and the typical length of the inhomogeneity $`L`$ is greater than the mean free path $`lu/\nu _c`$ ($`u`$ is the average particle velocity)
$$\omega /\nu _c1,u/L\nu _c1.$$
(6)
In a zero order with respect to parameters (6) Eq. (5) reduces to $`I_𝐩\left[f_𝐩\right]`$ $`=0`$, which means that the distribution function $`f_𝐩`$ and, consequently, the equation of state has a locally equilibrium form. As a result, the equations for the first three moments of the distribution function i.e. density $`n`$, current $`𝐣=n𝐯`$ and stress tensor $`P_{ij}=\delta _{ij}P`$ (which is diagonal in the local equilibrium) constitute the closed set of hydrodynamics equations for an ideal liquid.
It is clear that from the microscopic point of view BHT cannot be a consistent theory since it extends the collision dominated hydrodynamics ($`\nu _c\mathrm{}`$) to the electron gas where the collisionless (Vlasov) limit ($`\nu _c0`$) is most common. Due to the high frequency of the plasma waves at least the first of inequalities (6) is strongly violated and the tensor structure of $`P_{ij}`$ as well as that of higher moments becomes important. In a degenerate case these tensors describe deformation of the Fermi sphere - the effect completely ignored in the BHT. It has been shown recently that within a linear response theory this effect leads to an effective shear modulus of a liquid.
In this paper we present a generalized hydrodynamics which remains valid far beyond the local equilibrium, when both conditions (6) are strongly violated. The theory is restricted to the long-wave limit of the kinetic equation and requires inclusion of the tensors of higher moments. In Sec. II we derive a general hierarchy of equations for the moments of the distribution function and introduce a regular procedure of truncation of the infinite chain of equations, which does not require the assumption of a local equilibrium. We also discuss the physical meaning of restrictions which make possible the truncation. In Sec. III we construct a hydrodynamic theory of the second order in the long wavelength expansion and show a relation of the generalized hydrodynamics to the theory of elasticity and to the theory of highly viscous fluids. Sec. IV is devoted to the generalized hydrodynamics of a charged Fermi-liquid on the basis of Landau theory. We find the relationship between macroscopic hydrodynamical parameters and microscopic Fermi-liquid parameters. Sec. V contains examples of application of the generalized hydrodynamics to collective oscillations of confined systems. The purpose of this section is purely illustrative. To show clearly the nonequilibrium effects we select an analytically solvable model of a charged Fermi liquid in a harmonic potential of different dimensionality. We show that deviations from the local equilibrium strongly shift the excitation frequencies and substantially change the structure of collective modes. In Sec. VI we extend the theory to the fourth order in the long wavelength expansion. It is shown that in the low frequency limit the fourth-order theory transforms into the set of equations of a classical hydrodynamics, which includes viscosity and a heat transport contribution. In Sec. VII we summarize our results.
## II The hierarchy of equations for the moments
The kinetic equation (5) can be transformed in an infinite chain of equations for the moments of the distribution function. The zeroth and the first moment are the particle density $`n(𝐫,t)`$ and the current density $`𝐣(𝐫,t)`$:
$$n(𝐫,t)=\underset{𝐩}{}f_𝐩(𝐫,t),𝐣(𝐫,t)=\underset{𝐩}{}\frac{𝐩}{m}f_𝐩(𝐫,t)$$
(7)
The velocity field can be defined in a usual way as $`𝐯=𝐣/n`$. In general, the moment of the $`k`$th order is a tensor of $`k`$th rank
$$𝐌^{(k)}=M_{i_1\mathrm{}i_k}^{(k)}=\frac{1}{m^{k1}}\underset{𝐩}{}p_{i_1}\mathrm{}p_{i_k}f_𝐩,k>1$$
(8)
It is a common practice to separate the macroscopic motion of the liquid from the relative motion of particles via transformation to the Lagrange comoving frame. In this frame the macroscopic velocity is equal to zero and the particle density $`n(𝐫,t)`$ is the same as in the laboratory frame. The transformed distribution function $`f_𝐩^L(𝐫,t)`$ is related to the distribution function in the laboratory frame $`f_𝐩(𝐫,t)`$ as follows
$$f_𝐩^L=f_{𝐩+m𝐯}$$
and apparently has the properties
$`{\displaystyle \underset{𝐩}{}}f_𝐩^L(𝐫,t)`$ $`=`$ $`{\displaystyle \underset{𝐩}{}}f_𝐩(𝐫,t)=n(𝐫,t),`$ (9)
$`{\displaystyle \underset{𝐩}{}}{\displaystyle \frac{𝐩}{m}}f_𝐩^L(𝐫,t)`$ $`=`$ $`{\displaystyle \underset{𝐩}{}}{\displaystyle \frac{𝐩}{m}}f_𝐩(𝐫,t)𝐯(𝐫,t)n(𝐫,t)=0.`$ (10)
Eqs. (9) and (10) can be considered as a definition of the comoving frame.
The higher moments of the transformed distribution function are defined similar to (8):
$$𝐋^{(k)}=L_{i_1\mathrm{}i_k}^{(k)}=\frac{1}{m^{k1}}\underset{𝐩}{}p_{i_1}\mathrm{}p_{i_k}f_𝐩^L,k>1$$
(11)
Equations for moments $`L_{i_1\mathrm{}i_k}^{(k)}`$ can be easily derived from the kinetic equation for the distribution function $`f_𝐩^L`$:
$`D_tf_𝐩^L`$ $`+`$ $`{\displaystyle \frac{𝐩}{m}}f_𝐩^L(𝐩)𝐯{\displaystyle \frac{f_𝐩^L}{𝐩}}`$ (12)
$``$ $`\left(mD_t𝐯e\phi \right){\displaystyle \frac{f_𝐩^L}{𝐩}}=I_𝐩^L\left[f_𝐩^L\right],`$ (13)
where $`D_t=_t+𝐯`$ is a substantial derivative. The zeroth and the first moments of this equation give respectively the continuity equation and the equation for the velocity $`𝐯`$:
$`D_tn`$ $`+`$ $`n_iv_i=0`$ (14)
$`mnD_tv_i`$ $`+`$ $`_jL_{ij}^{(2)}en_i\phi =0`$ (15)
In the derivation of Eqs. (14)-(15) the constraints due to conservation of particles, momentum and energy in a collision process
$$\underset{𝐩}{}I_𝐩=0,\underset{𝐩}{}p_iI_𝐩=0,\underset{𝐩}{}𝐩^2I_𝐩=0,$$
(16)
have been used. The second moment $`L_{ij}^{(2)}`$ is a stress tensor, which is related to the pressure $`P`$ of an electron gas as $`TrL_{ij}^{(2)}=3P`$. This tensor can be decomposed into a scalar and a traceless parts
$$L_{ij}^{(2)}=P\delta _{ij}+\pi _{ij},Tr\pi _{ij}=0.$$
(17)
The equation of motion (15) allows to eliminate the electrostatic potential $`\phi `$ from the kinetic equation (13)
$$D_tf_𝐩^L+\frac{p_i}{m}_if_𝐩^L\left(p_j_jv_i\frac{_jL_{ij}^{(2)}}{n}\right)\frac{f_𝐩^L}{p_i}=I_𝐩^L$$
(18)
and, consequently, from all equations for higher moments. The term in brackets in Eq. (18) has a clear physical meaning. This is a force which acts on a particle in the Lagrange frame.
Calculation of the $`k`$th moment of the kinetic equation (18) leads to the equation for the $`k`$th moment of the distribution function $`L_{i_1\mathrm{}i_k}^{(k)}`$. The $`k`$th moment of the first term in the left hand side of Eq. (18) gives the substantial derivative of the $`k`$th moment of the distribution function $`𝐋^{(k)}`$. Similarly the divergency of the $`k+1`$th moment $`𝐋^{(k+1)}`$ comes from the second term. Contributions from the last two terms on the left have more complicated tensor structure. For example, the $`k`$th moment of the first term in the brackets in Eq. (18) is
$`{\displaystyle \frac{1}{m^{k1}}}{\displaystyle \underset{𝐩}{}}p_{i_1}\mathrm{}p_{i_k}p_j{\displaystyle \frac{f_𝐩^L}{p_i}}_jv_i`$ $`=`$ $`{\displaystyle \frac{1}{m^{k1}}}{\displaystyle \underset{𝐩}{}}f_𝐩^L{\displaystyle \frac{}{p_i}}p_{i_1}\mathrm{}p_{i_k}p_j_jv_i`$ (19)
$`=`$ $`\left[\delta _{ii_1}L_{i_2\mathrm{}i_kj}^{(k)}+\delta _{ii_2}L_{i_1i_3\mathrm{}i_kj}^{(k)}+\mathrm{}+\delta _{ij}L_{i_1\mathrm{}i_k}^{(k)}\right]_jv_i`$ (20)
The contribution from the last term in the left-hand side of Eq. (18) has a similar structure but contains the $`k1`$th moment $`𝐋^{(k1)}`$ instead of $`𝐋^{(k)}`$ in Eq. (20). Since the first moment of the distribution function $`f_𝐩^L`$ equals to zero (see Eq.(10)) this contribution vanishes in equation for the second moment. Hence, in the case of $`k=2`$ one has the following equation for the moments:
$`D_tL_{ij}^{(2)}`$ $`+`$ $`L_{ij}^{(2)}_kv_k+L_{ik}^{(2)}_kv_j+L_{kj}^{(2)}_kv_i`$ (21)
$`+`$ $`_kL_{ijk}^{(3)}=I_{ij}^{(2)},`$ (22)
where $`I_{ij}^{(2)}`$ is the second moment of the collision integral.
By the definition (11), all moments $`L_{i_1\mathrm{}i_k}^{(k)}`$ are symmetric with respect to transmutation of all indexes. All equations for moments must preserve this symmetry. To explicitly incorporate this condition it is convenient to introduce special notations for two possible symmetric products of a symmetric $`k`$-th rank tensor $`T_{i_1\mathrm{}i_k}^{(k)}`$ and a vector $`a_j`$. The first one is the dot product, which decreases the rank of a tensor from $`k`$ to $`k1`$ and corresponds to the contraction of the vector index $`j`$ with one of tensor indexes
$$𝐓^{(k)}𝐚T_{i_1\mathrm{}i_{k1}j}^{(k)}a_j.$$
(23)
We also define a symmetric direct product of $`𝐓^{(k)}`$ and $`𝐚`$ which gives a tensor of the rank $`k+1`$ and is equal to the sum over all transmutations of the vector index $`j`$ and tensor indexes $`i_1\mathrm{}i_k`$
$`\{𝐓^{(k)}𝐚\}_S`$ $``$ $`T_{i_1\mathrm{}i_k}^{(k)}a_j+T_{ji_2\mathrm{}i_k}^{(k)}a_{i_1}`$ (24)
$`+`$ $`T_{i_1j\mathrm{}i_k}^{(k)}a_{i_2}+\mathrm{}+T_{i_1\mathrm{}i_{k1}j}^{(k)}a_{i_k}.`$ (25)
These notations allow us to rewrite the moments of different terms of the kinetic equation (18) in a compact way. For example, the $`k`$th moment of the first term in the round brackets in Eq. (18), which is given by Eq. (20) above, can be expressed as
$`{\displaystyle \frac{1}{m^{k1}}}{\displaystyle \underset{𝐩}{}}`$ $`p_{i_1}`$ $`\mathrm{}p_{i_k}p_j{\displaystyle \frac{f_𝐩^L}{p_i}}_jv_i`$ (26)
$`=`$ $`𝐋^{(k)}(𝐯)+\{(𝐋^{(k)})𝐯\}_S.`$ (27)
Similarly, the moment of the second term in the brackets in Eq. (18) is
$$\frac{1}{m^{k1}}\underset{𝐩}{}p_{i_1}\mathrm{}p_{i_k}\frac{f_𝐩^L}{p_i}_jL_{ij}^{(2)}=\frac{1}{m}\{𝐋^{(k1)}(𝐋^{(2)})\}_S$$
Finally the system of equations for the moments takes the form
$`D_tn`$ $`+`$ $`n𝐯=0,`$ (28)
$`mnD_t𝐯`$ $`+`$ $`𝐋^{(2)}en\phi =0,`$ (29)
$`D_t𝐋^{(2)}`$ $`+`$ $`𝐋^{(2)}(𝐯)+\{(𝐋^{(2)})𝐯\}_S`$ (30)
$`+`$ $`𝐋^{(3)}=𝐈^{(2)}`$ (31)
$`D_t𝐋^{(k)}`$ $`+`$ $`𝐋^{(k)}(𝐯)+\{(𝐋^{(k)})𝐯\}_S+𝐋^{(k+1)}`$ (32)
$``$ $`{\displaystyle \frac{1}{mn}}\{𝐋^{(k1)}(𝐋^{(2)})\}_S=𝐈^{(k)},`$ (33)
where the last equation (33) is valid for $`k>2`$.
In a collision dominated case when conditions (6) are fulfilled, the high moments are small and the infinite chain of equations (28)-(33) can be truncated (Chapman-Enskog or Grad methods, see Ref. ). As the traceless part $`\pi _{ij}`$ of the stress tensor (17) and the third moment tensor $`L_{ijk}^{(3)}`$ vanish for the locally equilibrium distribution function, they remain small under conditions (6). In this case tensors $`\pi _{ij}`$ and $`L_{ijk}^{(3)}`$ describe, respectively, viscosity and thermal conductivity in a collision dominated liquid. However, if conditions (6) are violated, $`\pi _{ij}`$ and $`L_{ijk}^{(3)}`$ as well as the higher-order moments are not small. Yet the infinite chain of equations can be decoupled if all physical quantities are slow varying functions of $`𝐫`$. It is seen from Eq. (22) that the third moment $`L_{ijk}^{(3)}`$ enters only under a spatial derivative, hence its contribution is proportional to $`1/L`$. This remains true for all higher moments in higher-order equations (33). Thus the truncation of the chain can be guaranteed by the smallness of the gradients of the moments instead of the smallness of the moments itself. The dimensionless parameter of this expansion is
$$\gamma \frac{u}{Lmax\{\omega ,\nu _c\}}1.$$
(34)
It is clear that both inequalities (6) can be violated while the condition (34) is fulfilled. For collisionless nondegenerate plasma the possibility to decouple the chain of equations for moments based on the condition similar to Eq. (34) was first pointed out in Ref. , where it was referred to as “the low temperature approximation”.
To justify the consistency of this truncation procedure we estimate the order of magnitude of the terms in Eq. (33). According to Eq. (28), the spatial derivative of the velocity field $`𝐯`$ is of the order of the inverse characteristic time $`1/t_{pr}\omega `$ hence the first three terms in Eq. (33) are proportional to $`\omega `$. The right-hand side of this equation is evidently proportional to the collision frequency $`\nu _c`$. The last two terms in the left-hand side contain only spatial derivatives and are of the order of $`u/L`$. Thus the contribution of $`𝐋^{(k)}`$ in the equation for $`𝐋^{(k1)}`$ contains an additional smallness $`\gamma `$ (34), which means that the correction originating from $`𝐋^{(k)}`$ to the equation for the density (28) is of the order $`\gamma ^k`$. Therefore to obtain a theory valid up to $`\gamma ^k`$ one should keep $`k+1`$ equations for $`𝐋^{(q)}`$ ($`0qk`$) with the contribution of the $`k+1`$-th moment and the spatial derivative of $`𝐋^{(2)}`$ being omitted in the last equation. The resulting theory is a generalization of hydrodynamics which is valid far from the equilibrium.
It should be mentioned that, in principle, the described truncation procedure may not provide the closed system of equations because the moments of the collision integral are, in general, the functionals of the distribution function. However, for a high frequency process ($`\omega /\nu _c1`$) the contribution from the collision term disappears and we obtain the closed set of equations of the ”collisionless hydrodynamics” which corresponds to the Vlasov limit of the kinetic equation. The collision terms become eventually important on the long time scale, when the system approaches the equilibrium. In this limit the moments of the collision integral are linear functions of the corresponding moments of the distribution function. Hence under the condition ($`\omega /\nu _c1`$) they can be presented in a close form. Matching these two limits gives the hydrodynamical theory which correctly describes the high- and the low-frequency limit and is approximately valid for all values of $`\omega /\nu _c`$.
There is an obvious relation of the generalized hydrodynamics to the standard linear response theory. In the high-frequency region the collisionless hydrodynamics is a regular expansion of Vlasov equation in terms of parameter (34). Consequently in the case of a weak perturbation the linearized hydrodynamics equations should give the same results as the linearized Vlasov equation. For example, the response function obtained from the linearized system (28)-(33) which is truncated at $`k`$th equation should coincide with the microscopic response function up to $`(u/L\omega )^k(uq/\omega )^k`$ (where $`q1/L`$ is a wave vector of a perturbation). However the hydrodynamics in its general form goes beyond the linear response theory (RPA) since it is valid for nonlinear regime.
In this section we introduced parameter $`\gamma `$, which governs the decoupling procedure. The smallness of this parameter restricts the region of applicability of the generalized hydrodynamics. The condition (34) physically means that a particle passes the characteristic length scale $`L`$ in a time much longer than the typical time of the process $`1/\omega `$ or a collisional time $`1/\nu _c`$. In this case there are no particles which move in resonance with a collective motion and may provide an energy exchange between single-particle and collective excitation. In other words, the contribution of the Landau damping to the evolution of the system is small. This is the physical restriction of the applicability of the generalized hydrodynamics which is (as any hydrodynamics) the theory of a collective motion.
## III The second-order theory. Hydrodynamics of a degenerate Fermi gas
The theory of the second order with respect to parameter $`\gamma `$ (34) corresponds to neglect the third moment $`𝐋^{(3)}`$ in Eq. (31). Let us decompose the stress tensor into the scalar and the traceless parts (17) and write the system of equations in terms of the fields $`P`$ and $`𝝅`$:
$`D_t`$ $`n+n𝐯=0,`$ (35)
$`m`$ $`D_t𝐯+P+𝝅en\phi =0,`$ (36)
$`D_t`$ $`P+{\displaystyle \frac{5}{3}}P𝐯+{\displaystyle \frac{2}{3}}(𝝅)𝐯=0,`$ (37)
$`D_t`$ $`𝝅+𝝅(𝐯)+\{(𝝅)𝐯\}_S\mathrm{𝟏}{\displaystyle \frac{2}{3}}(𝝅)𝐯`$ (38)
$`+`$ $`P\left(\{𝐯\}_S\mathrm{𝟏}{\displaystyle \frac{2}{3}}𝐯\right)=𝐈^{(2)},`$ (39)
where $`\mathrm{𝟏}`$ is the unit tensor. Eqs. (37) and (39) are the trace and the traceless part of Eq. (31) respectively. Due to the conservation of energy in a collision process the second moment of the collision integral $`𝐈^{(2)}`$ is a traceless tensor. Consequently it contributes only to Eq. (39). Eqs. (35)-(39) are the hydrodynamics equations which contain two scalars $`n`$ and $`P`$, one vector $`𝐯`$ and one traceless tensor $`𝝅`$. For a charged liquid one has to add the Poisson equation (3) for the scalar potential $`\phi `$.
Eigenmodes of a charge liquid (plasma waves) lie in a high frequency region where the collisionless limit is reached. In Sec. IV (see also Ref. ) it is shown that the system of equation (35)-(39) leads to the correct plasmon dispersion up to the second order of parameter $`\gamma `$. This is a reflection of the fact that the linearized Eqs. (35)-(39) correspond to a correct linear response theory up to the $`\gamma ^2`$.
In the low frequency limit an expression for collision integral is needed. To demonstrate how a long-time scale behavior is recovered, we take $`I_𝐩`$ in a Krook-Bhatnager-Gross (KBG) approximation (see, i.e. Ref. )
$$I_𝐩[f_𝐩]=\nu _c\left(f_𝐩(𝐫,t)f_𝐩^F(𝐫,t)\right),$$
(40)
where $`f_𝐩^F`$ is a local equilibrium Fermi function with position-dependent velocity, chemical potential $`\mu (𝐫,t)`$ and temperature $`T(𝐫,t)`$. By the definition the function $`f_𝐩^F`$ is chosen to give the same values of velocity $`𝐯(𝐫,t)`$, density $`n(𝐫,t)`$ and pressure $`P(𝐫,t)`$, as the exact distribution function $`f_𝐩(𝐫,t)`$:
$`n(𝐫,t)`$ $`=`$ $`{\displaystyle \underset{𝐩}{}}f_𝐩={\displaystyle \underset{𝐩}{}}f_𝐩^F,`$ (41)
$`𝐣(𝐫,t)`$ $`=`$ $`{\displaystyle \underset{𝐩}{}}{\displaystyle \frac{𝐩}{m}}f_𝐩={\displaystyle \underset{𝐩}{}}{\displaystyle \frac{𝐩}{m}}f_𝐩^F,`$ (42)
$`P(𝐫,t)`$ $`=`$ $`{\displaystyle \underset{𝐩}{}}{\displaystyle \frac{𝐩^2}{3m}}f_{𝐩+m𝐯}={\displaystyle \underset{𝐩}{}}{\displaystyle \frac{𝐩^2}{3m}}f_{𝐩+m𝐯}^F.`$ (43)
Eqs. (41-43) guarantee that the collision integral (40) satisfies the general properties (16).
The second moment of $`I_𝐩`$ (40) is equal to
$$𝐈^{(2)}=\nu _c\left(𝐋^{(2)}\mathrm{𝟏}P_F(n,T)\right),$$
(44)
where $`P_F(n,T)`$ is the pressure of a Fermi gas with the distribution function $`f_𝐩^F`$. By definition (43), this pressure equals to the exact pressure $`P_F(n,T)=P\frac{1}{3}Tr𝐋^{(2)}`$, hence the expression in the brackets in (44) equals to the traceless tensor $`𝝅`$. Consequently, Eq. (44) for $`𝐈^{(2)}`$ is equivalent to two equations
$`𝐈^{(2)}`$ $`=\nu _c𝝅,`$ (45)
$`P`$ $`=P_F(n,T).`$ (46)
Since $`P_F(n,T)`$ is an equilibrium pressure for given $`n`$ and $`T`$, the equation (46) may seem to be related to the assumption of the local equilibrium. In fact, it is not. Eq. (46) simply introduces the new independent scalar variable $`T(𝐫,t)`$ instead of $`P(𝐫,t)`$. The system can still be in an arbitrary (up to $`\gamma ^2`$) nonequilibrium state. In the second order theory such state is uniquely described by two scalar, one vector and one traceless tensor functions. The traceless second rank tensor $`𝝅`$ is responsible for deviations from the local equilibrium. In a degenerate Fermi system this tensor describes deviations of the shape of the Fermi surface from sphere.
The system of equations (35)-(39), with the collision term (45) and $`P`$ from Eq. (46) transforms into the common hydrodynamics theory in the limit $`\omega /\nu _c1`$. Indeed, as we already mentioned above the spatial derivative of $`𝐯`$ has the order of magnitude of $`\omega `$. Hence, the first four terms in the left hand side of (39) are proportional to $`\omega 𝝅`$, the last term in the left hand side is proportional to $`\omega P`$ and the right hand side is equal to $`\nu _c𝝅`$. Consequently, in a zero order of $`\omega /\nu _c`$ Eq. (39) leads to $`\pi _{ij}=0`$, and Eqs.(35)-(37) become identical to hydrodynamics equations for an ideal liquid:
$`D_tn`$ $`+`$ $`n𝐯=0,`$ (47)
$`mnD_t𝐯`$ $`+`$ $`Pen\phi =0,`$ (48)
$`D_tP`$ $`+`$ $`{\displaystyle \frac{5}{3}}P𝐯=0.`$ (49)
Here Eq. (48) is the Euler equation and Eq. (49) is the equation for the conservation of energy. For a degenerate Fermi gas this set of equations exactly corresponds to the BHT. We emphasize again, that it is valid only in an extremely collision dominated regime.
In the first order of $`\omega /\nu _c1`$ one has to neglect the first four terms in the left hand side of Eq. (39) in comparison with the last term which is proportional to the pressure $`P`$. Thus, the first order solution of Eq. (39) takes the form of the viscosity tensor
$$\pi _{ij}=\frac{P}{\nu _c}\left(_iv_j+_jv_i\frac{2}{3}\delta _{ij}_kv_k\right),$$
(50)
with the viscosity coefficient $`\eta =P/\nu _c`$. In this case Eqs. (36) and (37) are equivalent to Navier-Stokes equation and the equation for energy conservation in a viscous liquid respectively. In Eq. (50) we have turned back to common tensor notations to make the relation to the ordinary viscous hydrodynamics more transparent.
The formulas (45), (46) were obtained using the KBG collision integral (44). They are, however, more general than the KBG approximation itself. As it has been mentioned at the end of the previous section, the reason is that in the low frequency range, where collisions are important, the second moment $`I_{ij}^{(2)}`$ is always a linear function of $`\pi _{ij}`$. The coefficient $`\nu _c`$ can thus be considered as a phenomenological parameter, which is related to viscosity $`\eta `$ as $`\nu _c=P/\eta `$.
Eqs. (35)-(39) with $`𝐈^{(2)}`$ from (45) and $`P`$ from (46) constitute the closed set of equations of generalized hydrodynamics in the second order of $`\gamma `$. This set of equations gives a correct description of the high-frequency collisionless regime. It also leads to the classical hydrodynamics of a viscous fluid in a low-frequency limit. Consequently it should be valid, with a reasonable accuracy, for intermediate regime.
However, the second order theory has an inherent inconsistency. Though in a low frequency limit the viscous term is recovered correctly, the thermal conductivity contribution $`^2T`$ (Ref. ) is still missing in the energy conservation equation. Without this contribution the correct static limit $`T(𝐫)=0`$ cannot be recovered. In fact, the term $`^2T`$ corresponds to correction of the fourth order of $`\gamma `$ in the continuity equation. Physically this means an absence of a dissipative flow of energy with an accuracy up to $`\gamma ^2`$. To include the thermal conductivity one has to consider the third and the fourth moments and neglect the fifth moment (see Sec. VI bellow).
The theory of the second order in $`\gamma `$ is obviously applicable for a description of dynamics of a dense degenerate Fermi gas. In this case the condition $`T/E_F1`$ ($`E_F`$ is the local Fermi energy) is always satisfied and it is unnecessary to take into account the third and the fourth moments. At $`T=0`$ Eq. (46) reduces to $`P=\frac{1}{5m}(3\pi ^2)^{2/3}n^{5/3}`$ and Eq. (39) transforms into
$$(𝝅)𝐯=0.$$
(51)
The equation for pressure $`P`$ and the condition (51) together with the system of differential equations
$`D_t`$ $`n+n𝐯=0,`$ (52)
$`m`$ $`D_t𝐯+P+𝝅en\phi =0,`$ (53)
$`D_t`$ $`𝝅+𝝅(𝐯)+\{(𝝅)𝐯\}_S\mathrm{𝟏}{\displaystyle \frac{2}{3}}(𝝅)𝐯`$ (54)
$`+`$ $`P\left(\{𝐯\}_S\mathrm{𝟏}{\displaystyle \frac{2}{3}}𝐯\right)=\nu _c𝝅,`$ (55)
provide a complete system of equations of the generalized hydrodynamics for a degenerate Fermi gas.
Most applications of the Bloch’s hydrodynamics concerned just the case of the degenerate Fermi gas. The set of equations (51)-(55) of generalized hydrodynamics should replace the BHT, which is incorrect except the collision dominated limit or the case of spatially homogeneous system.
There is an interesting relation of the linearized version of the generalized hydrodynamics in a high-frequency (collisionless) limit to the theory of elasticity. In the limit $`\omega /\nu _c1`$ linearization of the system (35)-(39) gives
$`_t\delta n`$ $`+`$ $`n_0_kv_k=0`$ (56)
$`mn_0_tv_i`$ $`+`$ $`_i\delta P+_j\pi _{ij}en_0_i\phi =0`$ (57)
$`_t\delta P`$ $`+`$ $`{\displaystyle \frac{5}{3}}P_0_kv_k=0`$ (58)
$`_t\pi _{ij}`$ $`+`$ $`P_0\left(_iv_j+_jv_i{\displaystyle \frac{2}{3}}\delta _{ij}_kv_k\right)=0,`$ (59)
where $`\delta n`$ and $`\delta P`$ are deviations from the equilibrium density $`n_0`$ and pressure $`P_0`$. Let us introduce a displacement vector $`𝐮(𝐫,t)`$ as $`\delta n=n_0𝐮`$. The continuity equation (56) gives a usual relation between velocity and displacement $`_t𝐮=𝐯`$. Introducing a stress tensor $`\sigma _{ij}=\delta P\delta _{ij}\pi _{ij}`$ one can rewrite Eq. (57) as
$$mn_0_t^2u_i_j\sigma _{ij}en_0_i\phi =0.$$
(60)
The relationship of the tensor $`\sigma _{ij}`$ to the displacement $`u_i`$ follows from Eqs. (58, 59) and takes the same form as in the elasticity theory
$$\sigma _{ij}=K_ku_k\delta _{ij}+\mu \left(_iu_j+_ju_i\frac{2}{3}\delta _{ij}_ku_k\right),$$
(61)
where the bulk modulus $`K`$ and the shear modulus $`\mu `$ of an electron gas are $`K=\frac{5}{3}P_0`$ and $`\mu =P_0`$. Physically the bulk modulus is responsible for the increase of the energy which is caused by the local change of the occupied volume in a momentum space, whereas the shear stress describes the deviation of the shape of this volume from a sphere. In the next section we derive the hydrodynamics for a Fermi liquid and find the correlation contribution to bulk and shear modulus. The elastic description of Fermi systems has been recently discussed in Ref. within the linear response theory.
Inclusion of collisions violates the exact correspondence to the elasticity theory. In this case the shear stress tensor $`\sigma _{ij}\frac{1}{3}\delta _{ij}\sigma _{kk}=\pi _{ij}`$ should be determined from the equation
$$_t\pi _{ij}+\mu _t\left(_iu_j+_ju_i\frac{2}{3}\delta _{ij}_ku_k\right)=\nu _c\pi _{ij},$$
(62)
which follows from Eqs. (59) and (45). Eq. (62) exactly coincides with the equation for the shear stress tensor in the phenomenological theory of highly viscous fluids by Maxwell. These fluids (for instance, glycerin or resin) behave as solids at short intervals of time, but as viscous liquids on a large time scale (see Ref. ). The generalized hydrodynamics thus provide a surprising but clear analogy between an electron gas and Maxwellian highly viscous fluids.
## IV The generalized hydrodynamics of a Fermi liquid
In this section we develop the hydrodynamic description of a Fermi liquid based on the Landau theory. Since the Landau theory is valid only for small deviations of the quasiparticle distribution function $`n_𝐩`$ from the Fermi function, only linearized version of the generalized hydrodynamics can be constructed. Although such theory is formally equivalent to the linear response approach, the hydrodynamical formulation is generally more efficient for spatially inhomogeneous problems. In addition, the derivation of the generalized hydrodynamics of a Fermi liquid is of a general interest since it transparently shows how correlation effects contribute to the stress tensor.
We start from a kinetic equation for a charged Fermi liquid
$$_tn_𝐩+\frac{\epsilon _𝐩}{𝐩}n_𝐩\epsilon _𝐩\frac{n_𝐩}{𝐩}+e\phi \frac{n_𝐩}{𝐩}=I_𝐩,$$
(63)
where $`\epsilon _𝐩`$ is the local quasiparticle energy, which depends on the distribution function and, consequently, is a function of spatial coordinates. The gradient of this function gives an additional force which has its origin in the correlation effects.
Following the derivation of the hydrodynamic equations for a Fermi gas (Sec. II) we separate the macroscopic and relative motion of a liquid by transformation to the comoving frame. The kinetic equation for the distribution function $`n_𝐩^L`$ in this frame looks similar to Eq. (13), but with an additional correlation force:
$`D_tn_𝐩^L`$ $`+`$ $`{\displaystyle \frac{\epsilon _𝐩^L}{𝐩}}n_𝐩^L(𝐩)𝐯{\displaystyle \frac{n_𝐩^L}{𝐩}}`$ (64)
$``$ $`\left(mD_t𝐯e\phi +\epsilon _𝐩^L\right){\displaystyle \frac{n_𝐩^L}{𝐩}}=I_𝐩^L`$ (65)
According to the Landau assumption, the quasiparticle energy is a linear functional of the deviation of the distribution function from the Fermi step function:
$`\epsilon _𝐩^L`$ $`=`$ $`\epsilon _𝐩^{(0)}+{\displaystyle \underset{𝐩^{}}{}}f_{𝐩^{}𝐩}\delta n_𝐩^{}^L,`$ (66)
$`n_𝐩^L`$ $`=`$ $`n_𝐩^{(0)}+\delta n_𝐩^L,n_𝐩^{(0)}=\theta (pp_F),`$ (67)
where $`\epsilon _𝐩^{(0)}`$ is the quasiparticle energy in the undisturbed system and $`f_{𝐩^{}𝐩}`$ is the quasiparticle-quasiparticle interaction. Linearization of the kinetic equation (65) leads to the Landau-Silin equation, written in a comoving frame
$`_t`$ $`\delta n_𝐩^L+{\displaystyle \frac{p_i}{m^{}}}_i\delta n_𝐩^Lp_j_jv_i{\displaystyle \frac{n_𝐩^{(0)}}{p_i}}`$ (68)
$``$ $`\left(m_tv_ie_i\phi +{\displaystyle \underset{𝐩^{}}{}}f_{𝐩^{}𝐩}_i\delta n_𝐩^{}^L\right){\displaystyle \frac{n_𝐩^{(0)}}{p_i}}=I_𝐩^L.`$ (69)
The quasiparticle mass $`m^{}`$, which enters equation (69), is related to the bare fermion mass $`m`$ as follows
$$\left(\frac{^2}{p^2}\epsilon _𝐩^{(0)}\right)^1m^{}=m\left(1+\frac{1}{3}F_1\right).$$
(70)
where $`F_1`$ is the first of Landau Fermi-liquid parameters:
$$F_l=\frac{m^{}p_F}{\pi ^2}\frac{d\mathrm{\Omega }^{}}{4\pi }f_{𝐩^{}𝐩}P_l(\mathrm{cos}\theta ),$$
(71)
where $`\theta `$ is the angle between $`𝐩`$ and $`𝐩^{}`$, $`d\mathrm{\Omega }^{}`$ is an element of a solid angle around $`𝐩^{}`$ and $`P_l(\mathrm{cos}\theta )`$ are Legendre polynomials.
Zeroth and first moments of Eq. (69) correspond to the continuity equation and the equation for velocity of a Fermi liquid:
$`_t\delta n`$ $`+`$ $`n_0_iv_i=0`$ (72)
$`mn_0_tv_i`$ $`+`$ $`_j\delta P_{ij}en_0_i\phi =0`$ (73)
In Eq. (73) we introduced the stress tensor of a Fermi liquid
$$\delta P_{ij}=\underset{𝐩}{}\frac{p_ip_j}{m^{}}\delta n_𝐩^L\underset{𝐩,𝐩^{}}{}f_{𝐩^{}𝐩}\frac{n_𝐩^{(0)}}{p_i}p_j\delta n_𝐩^{}^L$$
(74)
The first term in Eq. (74) corresponds to the kinetic stress tensor of quasiparticles, whereas the second term describes the interaction of qusiparticles and is responsible for a correlation contribution to the stress.
According to the results of the previous section, the second order approximation (up to $`\gamma ^2`$) gives a consistent hydrodynamical theory of a degenerate Fermi system. To make the system of equations (72)-(73) complete one needs only an equation for the second moment or the stress tensor. To simplify the derivation of this equation we separate the stress tensor $`\delta P_{ij}`$ into the scalar $`\delta P`$ and the traceless $`\pi _{ij}`$ parts
$$\delta P_{ij}=\delta P\delta _{ij}+\pi _{ij}$$
(75)
First we rewrite (74) as follows
$`\delta P_{ij}`$ $`=`$ $`{\displaystyle \frac{1}{m^{}}}{\displaystyle \underset{𝐩}{}}T_{ij}(𝐩)\delta n_𝐩^L`$ (76)
$`T_{ij}(𝐩)`$ $`=`$ $`p_ip_j+{\displaystyle \frac{m^{}p_F}{\pi ^2}}{\displaystyle \frac{d\mathrm{\Omega }^{}}{4\pi }f_{𝐩^{}𝐩}p_i^{}p_j^{}}`$ (77)
In the Landau theory all momentums reside on the Fermi surface. By symmetry tensor $`T_{ij}(𝐩)`$ is uniquely representable in the form
$$T_{ij}(𝐩)=\alpha p^2\delta _{ij}+\beta (p_ip_j\frac{1}{3}p^2\delta _{ij}),$$
(78)
where the constants $`\alpha `$ and $`\beta `$ are determined by the following two equations
$`{\displaystyle \frac{1}{p_F^2}}T_{ii}=3\alpha =1`$ $`+`$ $`{\displaystyle \frac{m^{}p_F}{\pi ^2}}{\displaystyle \frac{d\mathrm{\Omega }^{}}{4\pi }f_{𝐩^{}𝐩}},`$ (79)
$`{\displaystyle \frac{1}{p_F^4}}p_iT_{ij}p_j=\alpha +{\displaystyle \frac{2}{3}}\beta `$ $`=`$ (80)
$`1`$ $`+`$ $`{\displaystyle \frac{m^{}p_F}{\pi ^2}}{\displaystyle \frac{d\mathrm{\Omega }^{}}{4\pi }f_{𝐩^{}𝐩}\frac{(\mathrm{𝐩𝐩}^{})^2}{p_F^4}}.`$ (81)
Solutions to these equations are
$$\alpha =1+F_0,\beta =1+\frac{1}{5}F_2$$
(82)
Substituting $`\alpha `$ and $`\beta `$ (82) into Eqs. (78),(76) we obtain the following microscopic expressions for $`\delta P`$ and $`\pi _{ij}`$
$`\delta P`$ $`=`$ $`(1+F_0){\displaystyle \underset{𝐩}{}}{\displaystyle \frac{p^2}{3m^{}}}\delta n_𝐩^Lv_s^2\delta \rho ,`$ (83)
$`\pi _{ij}`$ $`=`$ $`(1+{\displaystyle \frac{1}{5}}F_2){\displaystyle \underset{𝐩}{}}\left\{{\displaystyle \frac{p_ip_j}{m^{}}}{\displaystyle \frac{p^2}{3m^{}}}\delta _{ij}\right\}\delta n_𝐩^L,`$ (84)
where $`v_s^2=(1+F_0)p_F^2/3mm^{}`$ is the square of velocity of sound and $`\delta \rho =m\delta n`$ is a variation of the mass density. The equivalence in Eq. (83) shows that the scalar part of the stress tensor $`\delta P`$ exactly equals to the variation of the pressure of the Fermi liquid.
Expressions for $`\delta P`$ and $`\pi _{ij}`$ (83) and (84) have the same form as for the gas case except the factors in front of the sums and the change of the bare mass to the mass of a quasiparticle. Thus the derivation of an equation for the stress tensor is straightforward. The final set of the generalized hydrodynamics equations of a Fermi liquid in the second order in parameter $`\gamma `$ (34) takes the form
$`_t\delta n`$ $`+`$ $`n_0_kv_k=0`$ (85)
$`mn_0_tv_i`$ $`+`$ $`_i\delta P+_j\pi _{ij}en_0_i\phi =0`$ (86)
$`_t\delta P`$ $`+`$ $`K_kv_k=0`$ (87)
$`_t\pi _{ij}`$ $`+`$ $`\mu \left(_iv_j+_jv_i{\displaystyle \frac{2}{3}}\delta _{ij}_kv_k\right)=\nu _c\pi _{ij}.`$ (88)
Hence, in a Fermi liquid we have again the system of equations which is similar to the Maxwell’s theory of highly viscous fluids, but with the bulk $`K`$ and the shear $`\mu `$ moduli directly related to microscopic Landau parameters
$`K`$ $`=`$ $`mn_0v_s^2={\displaystyle \frac{5}{3}}P_0{\displaystyle \frac{1+F_0}{1+\frac{1}{3}F_1}},`$ (89)
$`\mu `$ $`=`$ $`P_0{\displaystyle \frac{1+\frac{1}{5}F_2}{1+\frac{1}{3}F_1}},`$ (90)
where $`P_0`$ is a pressure of an ideal Fermi gas.
The bulk modulus (89) and the shear modulus (90) exactly coincide with those of Ref. , where they have been obtained by analyzing the linear response function. In ideal Fermi gas ($`F_l=0`$) Eqs. (89) and (90) transform to the ideal gas moduli (see Sec. III).
## V Hydrodynamic theory of collective modes
### A Plasma waves in a homogeneous system
The frequencies of collective plasma modes in a charged liquid have an order of magnitude of $`\omega _p`$ which commonly resides in the high frequency (collisionless) region $`\omega /\nu _c1`$. The plasma oscillations are solutions of the linearized system of the hydrodynamics equations. To obtain the dispersion of plasmons we have to solve Eqs. (85)-(88) with the collision term being omitted in Eq. (88). In the collisionless limit it is convenient to introduce a displacement vector $`𝐮`$ and rewrite this system in the form of the elasticity theory (60), (61), where $`K`$ and $`\mu `$ are determined, in general, by Eqs. (89) and (90). Substituting Eq. (61) to Eq.(60) we have an equation for the displacement vector of a ”charged elastic medium”
$$mn_0_t^2𝐮(K+\frac{1}{3}\mu )(𝐮)\mu ^2𝐮en_0\phi =0,$$
(91)
which should be solved together with the Poisson equation
$$^2\phi =4\pi en_0𝐮$$
(92)
Considering a plane-wave solution $`e^{i(\omega t\mathrm{𝐪𝐫})}`$ we get the dispersion of the plasma waves
$$\omega ^2(q)=\omega _p^2+v_0^2q^2$$
(93)
where
$$v_0^2=\frac{1}{mn_0}(K+\frac{4}{3}\mu )=v_s^2+\frac{4\mu }{3mn_0}.$$
(94)
The first contribution in $`v_0^2`$ equals to the square of the sound velocity $`v_s^2`$ (the result of the ordinary hydrodynamics) and comes from the fluctuations of pressure $`\delta P`$ (see formulas (83) and (89)). The second contribution arises from the traceless part of the stress tensor $`\pi _{ij}`$. The sum gives a correct coefficient which can be obtained from the longitudinal response function of a charged Fermi liquid. In the case of a noninteracting gas at arbitrary degeneracy the expression for $`v_0^2`$ reduces to the formula
$$v_0^2=\frac{5P_0}{3mn_0}+\frac{4P_0}{3mn_0}=\frac{3P_0}{mn_0}<v_p^2>.$$
which coincides with the result of the RPA (up to the second order of $`qu/\omega \gamma `$) and in a degenerate Fermi gas gives the well known result $`\frac{3}{5}v_F^2`$. It is straightforward to show that taking into account the third and fourth moments in the linearized system of equations leads to the correct plasmon dispersion up to $`q^4`$.
### B Plasma oscillations of a confined electron liquid
In the case of one dimensional motion (e.g. for plasma waves in an infinite medium which are considered in the previous subsection) the displacement vector $`𝐮`$ depends only on one coordinate. Hence, the equation of motion takes the following simple form:
$$_t^2𝐮v_0^2^2𝐮\frac{e}{m}\phi =0.$$
(95)
Eq. (95) coincides with that of the linearized Bloch’s theory. The only difference is that the dispersion coefficient equals to the correct value $`v_0^2`$ instead of a square of a velocity of sound $`v_s^2`$. It seems easy to phenomenologically improve the BHT. One could try to replace the dispersion coefficient by the correct value $`v_0^2`$ and hence obtain a correct dispersion of plasmon. However, from the general point of view such an improvement cannot be consistent since the correct static limit requires $`v_s^2`$ as a coefficient in Eq. (95). The frequency-dependent coefficient was commonly used in hydrodynamics calculations to recover both the high- and the low-frequency regimes (see Ref. and references therein). It has been then recognized that the theory with a frequency dependent dispersion coefficient violates the harmonic potential theorem (HPT) in a spatially inhomogeneous case. Nonetheless, it was believed that since the theory with the replaced coefficient gives the correct plasmon dispersion in the infinite medium it should hopefully give the proper description of the high-frequency plasma oscillations for any geometries of a confined electron gas.
As we show in the previous sections, the consistent theory which is valid both in a high- and low-frequency regimes must inevitably include the tensor fields of higher rank (e.g. $`\pi _{ij}`$ in the second order theory). At low frequency $`\pi _{ij}`$ goes to zero that gives a correct static limit. However, in the high-frequency regime it does not contain any small parameter and gives the contribution of the same order of magnitude as the variation of the scalar pressure. Within the linearized equations the contribution of the tensor $`\pi _{ij}`$ can be interpreted as a shear stress of a Fermi system. Formally this leads to more complicated structure of the differential operator in Eq. (91) than one of the operator expected from the BHT Eq. (95).
In the present subsection we consider two simple examples of a collective motion of a Fermi liquid which is confined in three and one dimensions. These examples show that the collective eigenfrequencies strongly differ from predictions of the standard hydrodynamics.
Consider a homogeneous (in equilibrium state) electron liquid which is confined by an external potential. The homogeneity of the equilibrium state can be reached for example with the help of a positively charged background of the density equal to the density of electrons $`n_0`$ or with the parabolic potential of a proper frequency $`\omega _0`$. For example, in a three-dimensional spherically symmetric case (parabolic quantum dot) the frequency $`\omega _0=\omega _p/\sqrt{3}`$. The two-dimensional (parabolic quantum wire) and one-dimensional (quantum well) potentials have frequencies $`\omega _0=\omega _p/\sqrt{2}`$ and $`\omega _0=\omega _p`$ respectively.
Eigenmodes are solutions to equations (60),(61),(92) with the proper boundary conditions. Integration of Eqs. (60) and (92) across the surface gives the set of boundary conditions
$`n_j\sigma _{ij}(𝐫_s)`$ $`=`$ $`0,`$ (96)
$`𝐧\phi ^>(𝐫_s)𝐧\phi ^<(𝐫_s)`$ $`=`$ $`4\pi en_0\mathrm{𝐧𝐮}(𝐫_s),`$ (97)
$`\phi ^>(𝐫_s)`$ $`=`$ $`\phi ^<(𝐫_s),`$ (98)
where $`𝐫_s`$ is a coordinate of a surface point, $`\phi ^>`$ and $`\phi ^<`$ are respectively the electrostatic potential $`\phi `$ outside and inside the surface and $`𝐧`$ is the unit vector normal to the surface. The boundary condition (96) corresponds to the ”free surface” and allows an electron liquid to cross the boundary and oscillate with respect to the equilibrium position. For example, the rigid oscillations of an electron liquid $`𝐮=𝐮_0e^{i\omega t}`$ ($`𝐮_0`$ is a constant vector) satisfy Eqs. (60),(61),(92) and conditions (96)-(98). Such oscillations have the frequency $`\omega =\omega _0`$. This fact apparently reflects the HPT.
The above system of equations and boundary conditions also satisfy the HPT (or generalized Kohn theorem) in the most general formulation, which was suggested by Dobson. We demonstrate this for a Fermi liquid confined by a spherically symmetric harmonic potential with $`\omega _0=\omega _p/\sqrt{3}`$ and subjected to external homogeneous electric $`𝐄`$ and magnetic $`𝐇`$ fields. The HPT motion corresponds to rigid displacement $`𝐮(𝐫,𝐭)=𝐮_0(t)`$ of the liquid. Vector $`𝐮_0(t)`$ is obviously the center-of-mass coordinate. Since the stress tensor is equal to zero, the boundary condition (96) is trivially satisfied and equation of motion (60) has the form
$$m_t^2𝐮_0e𝐇\times _t𝐮_0+e𝐄e\phi ^<=0.$$
(99)
The scalar potential $`\phi `$ satisfies the Laplace equation and electrostatic boundary conditions (97), (98). For rigid motion these conditions coincide with the conditions on the boundary of a homogeneously polarized sphere with a polarization vector $`𝐏=en_0𝐮_0`$. The soliton of this problem inside the sphere is a depolarization field
$$\phi ^<=\frac{4\pi }{3}𝐏=\frac{4\pi en_0}{3}𝐮_0.$$
Substituting this solution into the Eq.(99) we find that the center-of-mass moves according to equation
$$m_t^2𝐮_0e𝐇\times _t𝐮_0+e𝐄+\omega _0^2𝐮_0=0,$$
(100)
in exact correspondence with the HPT.
To avoid a confusion we would like to outline that the boundary condition (96) assumes the parabolic potential outside the region occupied by an electron liquid. Physically it can be realized, for example, at the boundary of a depletion region in p-n junctions or in artificial quantum structures. The condition (96) is not directly applicable to a model surface of metals with abrupt change of the positive background at the boundary. The reason is a jump of a potential at the surface which forbids electrons to penetrate the boundary.
Let us separate longitudinal and transverse components of the displacement $`𝐮`$
$$𝐮=\psi +𝐮_t,𝐮_t=0$$
(101)
($`\psi `$ is a potential of the longitudinal component) and rewrite the system of equations (60),(61),(92) and boundary conditions (96)-(98) in terms of variables $`\psi `$ and $`𝐮_t`$. We are looking for harmonic in time solutions $`e^{i\omega t}`$.
In the region occupied by an electron liquid we have the system of differential equations:
$`\omega ^2\psi +c_l^2^2\psi +{\displaystyle \frac{e}{m}}\phi `$ $`=`$ $`0,`$ (102)
$`\omega ^2𝐮_t+c_t^2^2𝐮_t`$ $`=`$ $`0,𝐮_t=0,`$ (103)
$`^2(\phi +4\pi en_0\psi )`$ $`=`$ $`0.`$ (104)
Outside the surface the Laplace equation for the scalar potential $`\phi `$ must be satisfied
$$^2\phi =0.$$
(105)
All solutions of Eqs. (102)-(105) must fulfill the boundary conditions
$`(c_l^22c_t^2)`$ $`n_i`$ $`^2\psi +2c_t^2n_j_j_i\psi `$ (106)
$`+`$ $`2c_t^2n_j(_iu_t^j+_ju_t^i)=0,`$ (107)
$`𝐧(\phi ^<`$ $`+`$ $`4\pi en_0\psi )=𝐧\phi ^>,`$ (108)
$`\phi ^<`$ $`=`$ $`\phi ^>,`$ (109)
where all functions are taken at the surface. In Eqs. (102)-(109) we introduced the notations:
$`c_l^2`$ $`=`$ $`v_0^2={\displaystyle \frac{1}{mn_0}}(K+{\displaystyle \frac{4}{3}}\mu ),`$ (110)
$`c_t^2`$ $`=`$ $`{\displaystyle \frac{\mu }{mn_0}}.`$ (111)
The boundary condition (107) demonstrates a peculiarity of the generalized hydrodynamics. Longitudinal and transverse oscillations are, in general, mixed at the surface. This effect is well known in the elasticity theory. If the shear modulus is nonzero, the incident purely longitudinal (transverse) wave transforms into the mixed wave containing both longitudinal and transverse components. As a result the structure of eigenmodes is changed in comparison with the common hydrodynamics ($`\mu =0`$). The qualitative structure of the wave is not changed only in the case of a normal incidence (for example radial modes in the spherical symmetric system which are purely longitudinal). However, even in this case the shear modulus contributes to the boundary condition. The boundary condition for a purely longitudinal motion takes the form
$$(c_l^22c_t^2)n_i^2\psi +2c_t^2n_j_j_i\psi =0$$
which strongly differs from the common hydrodynamical condition $`^2\psi =0`$ for a parabolic potential.
Bellow we solve equations (102)-(109) for two specific cases and show that nonzero shear modulus strongly influences the dispersion of the eigenmodes even for purely longitudinal waves.
#### 1 Radial plasma oscillations in a parabolic quantum dot.
In a spherical symmetric quantum dot the normal vector $`𝐧=\widehat{𝐫}`$. Thus all normal derivatives transform to the derivative with respect to the radial coordinate $`r`$. The direction of the displacement $`𝐮`$ for radial modes coincides with $`\widehat{𝐫}`$. Hence the transverse component $`𝐮_t`$ equals to zero and functions $`\psi `$ and $`\phi `$ depend only on $`r`$.
Radial solutions of Eqs. (102), (104), (105) which are regular at zero and at the infinity take the following general form
$$\psi (r)=A_1\frac{\mathrm{sin}qr}{r}+B_1,\phi _<(r)=A_2\frac{\mathrm{sin}qr}{r}+B_2,$$
(113)
$$\phi _>(r)=\frac{C}{r},$$
(114)
Substituting (V B 1) into Eqs. (102), (104) and into boundary conditions for the scalar potential (108), (109) one gets the result
$`\psi (r)`$ $`=`$ $`A_1\left\{{\displaystyle \frac{\mathrm{sin}qr}{r}}{\displaystyle \frac{\omega _p^2}{\omega _p^2+c_l^2q^2}}{\displaystyle \frac{\mathrm{sin}qR}{R}}\right\}`$ (115)
$`\phi _<(r)`$ $`=`$ $`A_14\pi en_0\left\{{\displaystyle \frac{\mathrm{sin}qr}{r}}{\displaystyle \frac{\mathrm{sin}qR}{R}}\right\}`$ (116)
Here $`R`$ is the radius of the quantum dot. The wave vector $`q`$ is related to the frequency $`\omega `$
$$\omega ^2=\omega _p^2+c_l^2q^2.$$
(117)
The last boundary condition (107) in the spherical symmetric case reads
$$(c_l^22c_t^2)^2\psi (R)+2c_t^2_r^2\psi (R)=0.$$
This equation together with the solution (115) gives the final dispersion equation
$$\mathrm{tan}qR=\frac{qR}{1(c_l/2c_t)^2q^2R^2}$$
(118)
which determine the allowed values of $`q`$ and, consequently, the frequencies of the eigenmodes (117). The corresponding result of BHT can be obtained from Eq. (118) in the limit $`c_t0`$ and reads
$$\mathrm{tan}qR=0.$$
The quantity $`(c_l/2c_t)^2`$ which governs the difference of the last two dispersion equations can be expressed in terms of Landau parameters:
$$(c_l/2c_t)^2=\frac{1}{3}\left(1+\frac{5}{4}\frac{1+F_0}{1+F_2/5}\right).$$
This value is obviously far from infinity (as assumed in the BHT). For example, in the case of a Fermi gas ($`F_l=0`$) we have $`(c_l/2c_t)^2=3/4`$ which is less then one.
#### 2 Surface plasma modes at the edge of a parabolic potential well.
Let us consider the surface plasma oscillations at the edge of an electron system confined by a one dimensional parabolic potential with $`\omega _0=\omega _p`$. The situation is common in p-n junctions where the potential in the depletion region is parabolic and density of electrons (holes) in n (p) regions is approximately constant up to the corresponding boundary of the depletion region. Possible applications of the results to the surface modes of the parabolic potential well is obvious. Within BHT this problem has been considered in Ref. .
We suppose that the electron gas occupies the lower half space ($`z<0`$). Let the surface wave propagate along $`x`$-axis $`\psi ,𝐮_te^{iqx}`$. Due to the symmetry, the $`y`$-component of the vector $`𝐮_t`$ equals to zero. Hence the solution of Eq. (103) can be taken as
$$𝐮_t=𝐚e^{\kappa z}e^{iqx}$$
(119)
where $`𝐚=(a_x,0,a_z)`$ is a two dimensional constant vector and
$$\kappa ^2=q^2\omega ^2/c_t^2.$$
(120)
Due to the condition $`𝐮_t=0`$ (see Eq. (103)) the constants $`a_x`$ and $`a_z`$ are not independent:
$$a_x=\frac{\kappa }{iq}a_z.$$
The general solution of Es.(102, 104) takes the form
$`\psi `$ $`=`$ $`\left(Ae^{qz}+Be^{pz}\right)e^{iqx},`$ (121)
$`\phi _<`$ $`=`$ $`\left({\displaystyle \frac{\omega ^2m}{e}}Ae^{qz}+4\pi en_0Be^{pz}\right)e^{iqx},`$ (122)
with
$$p^2=q^2+(\omega _p^2\omega ^2)/c_l^2.$$
In the upper half-space we have the solution of the Poisson equation (105):
$$\phi _>=Ce^{qz}e^{iqx}.$$
(123)
The relationship of the constants $`B`$ and $`C`$ to the constant $`A`$ can be easily obtained from the electrostatic boundary conditions (108), (109). The result for the potential of the irrotational part of displacement $`\psi `$ is
$$\psi =A\left[e^{qz}+\left(1\frac{\omega ^2}{\omega _s^2}\right)e^{pz}\right]e^{iqx}$$
(124)
where $`\omega _s=\omega _p/\sqrt{2}`$ is the classical surface plasmon frequency in the infinite wave length limit. Let us rewrite the boundary condition (107) explicitly in the case of the plane boundary :
$`(`$ $`c_l^2`$ $`2c_t^2)n_i^2\psi +2c_t^2_z^2\psi +4c_t^2_zu_t^z=0,`$ (125)
$`c_t^2`$ $`(_x_z\psi +_xu_t^z+_zu_t^x)=0.`$ (126)
Substitution of the solutions (119) and (124) to the boundary conditions gives the system of equations for the frequencies of the surface plasma modes:
$`\left(1{\displaystyle \frac{\omega ^2}{\omega _p^2}}\right)\left(\omega ^2\omega _s^22c_t^2q^2\right)A2c_t^2\kappa a_z=0,`$ (127)
$`c_t^2q^2\left[q+\left(1{\displaystyle \frac{\omega ^2}{\omega _s^2}}\right)p\right]A+c_t^2\left(\kappa ^2+q^2\right)a_z=0.`$ (128)
At $`q=0`$ we have
$$\omega ^2=\omega _s^2\frac{\omega _p^2}{2};a_z=0,$$
which corresponds to the usual infinite wave length surface plasma oscillations. At small $`q`$ the dispersion of the surface plasmon takes the form
$$\omega ^2=\omega _s^2+2c_t^2q^2+i\frac{4c_t^3q^3}{\omega _s}.$$
(129)
The damping of the surface plasmon is not surprising since at $`q0`$ there is mixing with the transverse modes which can not decay from the surface at high frequency. There are decaying transverse solutions only for $`\omega ^2<c_t^2q^2`$. Hence only the propagating transverse waves are allowed at the frequency $`\omega _s`$. The coupling of the longitudinal and transverse components leads to the energy transfer to propagating waves and, consequently, provides damping of the surface plasmon. It should be mentioned that the hydrodynamic theory does not take into account the Landau damping due to the single particle excitations. Inclusion of the Landau damping will bring the usual contribution to the lifetime of the surface plasmon. Yet in a weakly coupled electron gas the Landau damping destroys the transverse modes. More precisely, if $`c_t<v_F`$ the transverse modes are damped. Hence the decay channel of the surface plasmon into the transverse waves will be really open if the inverse inequality $`c_t>v_F`$ is fulfilled and transverse modes are well defined. Using Eq. (111) for $`c_t`$ and Eq. (90) for $`\mu `$ we can express the condition $`c_t>v_F`$ in terms of Landau Fermi liquid parameters:
$$F_1(1+F_2/5)+3F_2/5>6.$$
Consequently, the transverse modes and the decay channel described above could exist for strongly correlated charged Fermi liquid. It was mentioned in Ref. that the condition should be satisfied for a electron liquid of low density.
With the accuracy up to $`q^2`$ we can neglect the imaginary term in the dispersion law (129). In this long wavelength approximation the surface plasmon is a longitudinal wave with the dispersion
$$\omega ^2=\omega _s^2+2c_t^2q^2.$$
(130)
The last equation shows that the dispersion of the almost longitudinal surface plasmon is totally determined by the shear modulus ($`c_t^2\mu `$). Within the BHT the shear modulus equals to zero. In the case $`c_t=0`$ equations (127, 128) lead to the dispersionless surface plasmon with the frequency $`\omega =\omega _s`$. This agrees with the results obtained in for a wide parabolic quantum well. The correct dispersion law (130) obtained from the generalized hydrodynamics has the dispersion coefficient $`2c_t^2`$ which is of the same order of magnitude as the coefficient $`v_0^2=c_l^2`$ in the volume plasmon dispersion (93). For example, in the case of a high density electron gas ($`F_l0`$) we have the ratio
$$2c_t^2/c_l^2=\frac{2}{3}.$$
## VI The fourth-order theory. Thermal conductivity
As we have seen in Sec. III, the theory of the second order in parameter $`\gamma `$ (34) does not contain effects of thermal conductivity which are responsible for the thermal equilibration processes. Hence the second-order approximation cannot give correct static solution which should correspond to a constant temperature ($`T(𝐫)=0`$). To account for the heat conduction one has to consider at least fourth-order approximation which contains two additional equations for the third $`𝐋^{(3)}`$ and fourth $`𝐋^{(4)}`$ moments. The system of equations follows Eqs. (28)-(33). The first two equations - the continuity equation and the equation for the velocity field, are the same as in the second-order theory (Eqs. (35) and (36) respectively). In the equations for the scalar $`P`$ and traceless $`𝝅`$ parts of the second moment we get additional terms which are proportional to the gradient of the $`𝐋^{(3)}`$ (see Eq. 31):
$`D_tP`$ $`+`$ $`{\displaystyle \frac{5}{3}}P𝐯+{\displaystyle \frac{2}{3}}(𝝅)𝐯+Tr(𝐋^{(3)})=0,`$ (131)
$`D_t𝝅`$ $`+`$ $`𝝅(𝐯)+\{(𝝅)𝐯\}_S\mathrm{𝟏}{\displaystyle \frac{2}{3}}(𝝅)𝐯+P\left(\{𝐯\}_S\mathrm{𝟏}{\displaystyle \frac{2}{3}}𝐯\right)+𝐋^{(3)}Tr(𝐋^{(3)})=𝐈^{(2)}.`$ (132)
To close the set of equations with the accuracy of $`\gamma ^4`$ we should take an equation for $`𝐋^{(3)}`$ in the general form (33) and omit the spatial derivatives of the fifth and the second moments in the equation for $`𝐋^{(4)}`$:
$`D_t𝐋^{(3)}`$ $`+`$ $`𝐋^{(3)}(𝐯)+\{(𝐋^{(3)})𝐯\}_S+𝐋^{(4)}{\displaystyle \frac{1}{mn}}\{𝐋^{(2)}(𝐋^{(2)})\}_S=𝐈^{(3)},`$ (133)
$`D_t𝐋^{(4)}`$ $`+`$ $`𝐋^{(4)}(𝐯)+\{(𝐋^{(4)})𝐯\}_S=𝐈^{(4)}.`$ (134)
Eqs. (35), (36) and (131)-(134) constitute the system of the generalized hydrodynamics equations in the fourth-order approximation.
The last term in the left-hand side of Eq. (131) can be rewritten as a divergency of a vector
$$Tr(𝐋^{(3)})=_jL_{iij}^{(3)}=_jQ_j,$$
(135)
where the vector $`Q_j=L_{iij}^{(3)}`$ describes the flow of energy in a comoving frame.
To demonstrate the recovering of the thermal conductivity we take the moments of the collision integral $`𝐈^{(3)}`$ and $`𝐈^{(4)}`$ in a linear approximation, similar to (45):
$`𝐈^{(3)}`$ $`=`$ $`\nu _{c1}𝐋^{(3)},`$ (136)
$`𝐈^{(4)}`$ $`=`$ $`\nu _{c2}\left(𝐋^{(4)}𝐋_F\right),`$ (137)
where we introduced the notation $`𝐋_F`$ for the fourth moment of the distribution function $`f_𝐩^L(𝐫,t)`$ (see Eqs. (40)-(43)):
$$𝐋_F=\frac{1}{m^3}\underset{𝐩}{}p_ip_jp_kp_lf_𝐩^L$$
(138)
Due to the spherical symmetry of the function $`f_𝐩^L`$ the fourth rank tensor $`𝐋_F`$ (138) takes the form:
$`𝐋_F`$ $`=`$ $`(\delta _{ij}\delta _{kl}+\delta _{kj}\delta _{il}+\delta _{lj}\delta _{ik})L_F,`$ (139)
$`L_F`$ $`=`$ $`{\displaystyle \frac{1}{15}}{\displaystyle \underset{𝐩}{}}{\displaystyle \frac{p^4}{m^3}}f_𝐩^L.`$ (140)
Parameters $`\nu _{c1}`$ and $`\nu _{c2}`$ in Eqs. (136) and (137) can be considered as phenomenological coefficients which are, in general, not equal to the parameter $`\nu _c`$ (45), but have the same order of magnitude.
In the low frequency limit $`\omega /\nu _c1`$ we solve Eqs. (132)-(134) using a perturbation theory. Since the collision terms in all these equations have an order of magnitude of the collision frequency the solution in zero order of $`1/\nu _c`$ follows from the equations
$$𝐈^{(𝐣)}=0,j=2,3,4$$
and takes the local equilibrium form
$$𝐋^{(4)}=𝐋_F,𝐋^{(3)}=0,𝝅=0.$$
The stress tensor becomes diagonal and proportional to the local pressure
$$𝐋^{(2)}=\mathrm{𝟏}P_F$$
The first order solution of equation (39) for the tensor $`𝝅`$ has been already obtained in Sec. III. It has the form of the viscosity tensor (50). Eq. (132), which determines $`𝝅`$ with the accuracy of $`\gamma ^4`$, differs from Eq. (39) since the third rank tensor $`𝐋^{(3)}`$ contributes in the left-hand side of (132). However, the contribution of the third moment in the Eq. (132), being of the higher order in $`\omega /\nu _c`$, does not change the first order result of Sec. III. Hence the viscous solution (50) is still valid.
Let us consider solution of Eqs. (133), (134) for the third moment $`𝐋^{(3)}`$. In the first-order of $`\omega /\nu _c`$ only the last two terms in the left-hand side in Eq. (133) contribute to the equation
$$\nu _{c1}𝐋^{(3)}𝐋^{(4)}\frac{1}{mn}\{𝐋^{(2)}(𝐋^{(2)})\}_S,$$
(141)
where zero-order expressions for $`𝐋^{(4)}`$ and $`𝐋^{(2)}`$ should be used in the right-hand side.
Substituting zero-order (local equilibrium) solutions ($`𝐋^{(2)}=\mathrm{𝟏}P_F`$ and $`𝐋^{(4)}=𝐋_F`$) to Eq. (141) one gets the the first-order expression for the third moment
$$𝐋^{(3)}=\frac{1}{\nu _{c1}}\left\{\mathrm{𝟏}\left(L_F\frac{1}{mn}P_FP_F\right)\right\}_S.$$
(142)
Contraction of the tensor $`𝐋^{(3)}`$ (142) over the couple of indexes gives the energy flow vector $`Q_j=L_{iij}^{(3)}`$, which enters the equation of the energy conservation (131)
$$𝐐=\frac{5}{\nu _{c1}}\left(L_F\frac{1}{mn}P_FP_F\right).$$
(143)
Let us show that vector $`𝐐`$ is proportional to the gradient of temperature $`T(𝐫,t)`$ at arbitrary degeneracy of an electron gas. The functions $`L_F`$ and $`P_F`$ depend on spatial coordinates via the spatial dependence of the local chemical potential $`\mu (𝐫,t)`$ and the local temperature $`T(𝐫,t)`$ entering the distribution function $`f_𝐩^L(𝐫,t)`$
$$f_𝐩^L(𝐫,t)=f^L\left(\frac{\epsilon _𝐩\mu }{T}\right),\epsilon _𝐩=\frac{p^2}{2m}.$$
Hence gradients of $`P_F`$ and $`L_F`$ can be rewritten as follows
$`P_F`$ $`=`$ $`{\displaystyle \frac{1}{3m}}{\displaystyle \underset{𝐩}{}}p^2f_𝐩^L=\left(\mu +{\displaystyle \frac{\mu }{T}}T\right){\displaystyle \frac{1}{3m}}{\displaystyle \underset{𝐩}{}}p^2{\displaystyle \frac{f_𝐩^L}{\epsilon _𝐩}}{\displaystyle \frac{T}{T}}{\displaystyle \frac{1}{6m^2}}{\displaystyle \underset{𝐩}{}}p^4{\displaystyle \frac{f_𝐩^L}{\epsilon _𝐩}}`$ (144)
$`L_F`$ $`=`$ $`{\displaystyle \frac{1}{15m^3}}{\displaystyle \underset{𝐩}{}}p^4f_𝐩^L=\left(\mu +{\displaystyle \frac{\mu }{T}}T\right){\displaystyle \frac{1}{15m^3}}{\displaystyle \underset{𝐩}{}}p^2{\displaystyle \frac{f_𝐩^L}{\epsilon _𝐩}}{\displaystyle \frac{T}{T}}{\displaystyle \frac{1}{30m^4}}{\displaystyle \underset{𝐩}{}}p^6{\displaystyle \frac{f_𝐩^L}{\epsilon _𝐩}}`$ (145)
The partial integration allows one to express the momentum integrals in Eqs. (144), (145) in terms of macroscopic variables $`n`$, $`P_F`$ and $`L_F`$:
$`{\displaystyle \frac{1}{3m}}{\displaystyle \underset{𝐩}{}}p^2{\displaystyle \frac{f_𝐩^L}{\epsilon _𝐩}}`$ $`=`$ $`{\displaystyle \underset{𝐩}{}}f_𝐩^L=n,`$ (146)
$`{\displaystyle \underset{𝐩}{}}p^4{\displaystyle \frac{f_𝐩^L}{\epsilon _𝐩}}`$ $`=`$ $`5m{\displaystyle \underset{𝐩}{}}p^2f_𝐩^L=15m^2P_F,`$ (147)
$`{\displaystyle \frac{1}{30m^4}}{\displaystyle \underset{𝐩}{}}p^6{\displaystyle \frac{f_𝐩^L}{\epsilon _𝐩}}`$ $`=`$ $`{\displaystyle \frac{7}{30m^3}}{\displaystyle \underset{𝐩}{}}p^4f_𝐩^L={\displaystyle \frac{7}{2}}L_F.`$ (148)
Using the last expressions we can represent Eqs. (144) and (145) in the following compact form
$`P_F=\left(\mu {\displaystyle \frac{\mu }{T}}T\right)n+{\displaystyle \frac{5P_F}{2T}}T,`$ (149)
$`L_F=\left(\mu {\displaystyle \frac{\mu }{T}}T\right){\displaystyle \frac{P_F}{m}}+{\displaystyle \frac{7L_F}{2T}}T.`$ (150)
Comparison of Eqs. (149) and (150) with Eq. (143) shows that the contribution of the gradient of the chemical potential exactly cancels in the equation for the energy flow vector. Thus the vector of the energy flow takes the usual form
$$𝐐=\lambda T$$
(151)
with the coefficient of heat conduction
$$\lambda =\frac{5}{2\nu _{c1}T}\left(7L_F\frac{5P_F^2}{mn}\right).$$
(152)
Finally Eq. (131) with $`𝝅`$ (39) and $`𝐐`$ (152) transforms to the common hydrodynamical equation of the energy conservation
$$D_tP+\frac{5}{3}P𝐯+\frac{2}{3}(𝝅)𝐯+𝐐=0,$$
(153)
which contain both the viscous and the heat conduction terms.
The formula (152) for the coefficient of heat conduction is valid at arbitrary degree of degeneracy of an electron gas. In the nondegenerate case $`P_F=nT`$ and $`L_F=nT^2/m`$. Hence Eq. (152) transforms to the result:
$$\lambda =\frac{5nT}{\nu _{c1}m}.$$
In a degenerate gas at $`T=0`$ we have $`L_F=5P_F^2/7mn`$ hence the expression in the brackets in (152) cancels. At $`T/\epsilon _F1`$ the correction has the order of magnitude $`(T/\epsilon _F)^2`$ and the coefficient of heat conduction (152) is proportional to $`T/\nu _{c1}`$. Since the collision frequency in a degenerate Fermi system is a quadratic function of $`T`$ the heat conduction diverges as $`1/T`$ (see, for example, Ref. ).
Thus, in the low frequency limit the generalized hydrodynamics equations (35), (36, (131)-(134) coincides with the correct set of the common hydrodynamics which consist of the continuity equation, Navier-Stokes equation and the equation of the energy conservation in the form (153).
In the high frequency (collisionless) limit the generalized hydrodynamics leads to the following plasmon dispersion
$$\omega ^2=\omega _p^2+<v_𝐩^2>q^2+\frac{<v_𝐩^4><v_𝐩^2>^2}{\omega _p^2}q^4$$
that is exactly the result of the linear response theory.
The system of equations (35), (36), (131)-(134) with the collision terms (45), (136), (137) constitute the closed set of equations and provides the smooth interpolation between high and low frequency regimes.
## VII Conclusion
A description of dynamics of many-electron systems in terms of macroscopic collective variables, which is usually referred to as a hydrodynamical approach, provides a simple, physically transparent and powerful tool for studying spatially inhomogeneous problems. The common Bloch’s hydrodynamics of an electron gas is based on the assumption of a local equilibrium and consequently shows a number of inconsistencies. We have shown that it is possible to construct an inherently consistent generalized hydrodynamics which correctly describes both the collisionless high-frequency limit and the collision dominated low-frequency regime, where the theory coincides with the standard Navier-Stokes hydrodynamics. The theory follows from the long-wavelength expansion of a kinetic equation and requires inclusion of new collective variables with a nontrivial tensor structure. We remind that only scalar (pressure and density) and vector (velocity) variables enter the common hydrodynamics.
The appearance of higher rank tensors is physically evident since the occupied region in a momentum space loses its spherical symmetry under a general (not locally equilibrium) evolution of the system. A need to describe a non-spherical isoenergetic surface inevitably requires tensor variables. As long as this surface is smooth it is possible to approximately describe it as a surface of a finite order and thus to take into account the tensor fields of the finite rank. We have actually shown above that the smoothness of the surface, which bounds the occupied region in the momentum space, is governed by the basic parameter (34).
In extension of the previous publication, we presented a new general and more transparent derivation of the generalized hydrodynamics and showed how the closed set of equations of the standard Navier-Stokes hydrodynamics (including the heat transport) is recovered in the low frequency regime. We have shown that the generalized hydrodynamics can be built on the basis of Landau theory of Fermi liquid which allows to determine the correlation contribution to the stress tensor. Although this hydrodynamics (as the Landau theory of Fermi liquid itself) is meaningful only in the linear approximation, we belive that it could be of a practical importance since it allows to express the results of hydrodynamical calculations, which are relatively simple, in terms of the microscopic Landau parameters. For example, the theory offers a possibility to determine the contribution of Fermi-liquid correlation effects to the eigenmodes of spatially inhomogeneous systems.
In a linear approximation the tensor variables, which describe the absence of a local equilibrium, can be interpreted as an effective shear stress of a liquid. On a time scale much longer than the collision time the contribution of the shear stress vanishes and the dynamics of a liquid is governed only by the usual bulk stress. However, as the collective modes of a charged liquid normally belong to the high frequency range, the shear contribution cannot be neglected. It is important to realize that this contribution cannot be modeled, even qualitatively, merely by a change of the bulk modulus - the procedure, which has been frequently used in literature to obtain the correct dispersion of plasmon. The reason is that the relative contribution of the shear and bulk stresses is different for different modes and different geometries. As a result the eigenfrequencies and the structure of plasma modes in confined systems strongly deviate from predictions of BHT even with “improved” bulk modulus, which is adjusted to provide a correct plasmon dispersion in a homogeneous situation. To demonstrate this qualitative non-local-equilibrium effect we calculated eigenmodes of a Fermi liquid confined by a harmonic potential of different dimensionality. The results show a nontrivial contribution of the shear modulus to eigenfrequencies of the longitudinal plasma oscillation. The most transparent result concerns surface waves at the edge of 1D harmonically trapped system. In the absence of the shear modulus these modes are absolutely dispersionless in agreement with the BHT results. They acquire a dispersion of the same order of magnitude as the bulk plasmon due to deviation from the local equilibrium. We have also found that existence of a nonzero shear modulus may open a new channel of decay of these surface modes, in addition to the usual Landau damping.
The work of one of the authors (I.T.) was supported by the Alexander von Humboldt Foundation.
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# Algebraic Surfaces Holomorphically Dominable by ℂ²
## 1 Introduction
An $`n`$-dimensional complex manifold $`M`$ is said to be (holomorphically) dominable by $`^n`$ if there is a map $`F:^nM`$ which is holomorphic such that the Jacobian determinant $`det(DF)`$ is not identically zero. Such a map $`F`$ is called a dominating map. In this paper, we attempt to classify algebraic surfaces $`X`$ which are dominable by $`^2`$ using a combination of techniques from algebraic topology, complex geometry and analysis. One of the key tools in the study of algebraic surfaces is the notion of Kodaira dimension (defined in section 2). By Kodaira’s pioneering work \[Ko1\] and its extensions (see, for example, \[CG\] and \[KO\]), an algebraic surface which is dominable by $`^2`$ must have Kodaira dimension less than two. Using the Kodaira dimension and the fundamental group of $`X`$, we succeed in classifying algebraic surfaces which are dominable by $`^2`$ except for certain cases in which $`X`$ is an algebraic surface of Kodaira dimension zero and the case when $`X`$ is rational without any logarithmic $`1`$-form. More specifically, in the case when $`X`$ is compact (namely projective), we need to exclude only the case when $`X`$ is birationally equivalent to a K3 surface (a simply connected compact complex surface which admits a globally non-vanishing holomorphic 2-form) that is neither elliptic nor Kummer (see sections 3 and 4 for the definition of these types of surfaces).
With the exceptions noted above, we show that for any algebraic surface of Kodaira dimension less than 2, dominability by $`^2`$ is equivalent to the apparently weaker requirement of the existence of a holomorphic image of $``$ which is Zariski dense in the surface. With the same exceptions, we will also show the very interesting and revealing fact that dominability by $`^2`$ is preserved even if a sufficiently small neighborhood of any finite set of points is removed from the surface. In fact, we will provide a complete classification in the more general category of (not necessarily algebraic) compact complex surfaces before tackling the problem in the case of non-compact algebraic surfaces.
We remark that both elliptic K3 and Kummer K3 surfaces are dense in the moduli space of K3 surfaces; the former is dense of codimension-one while the latter is dense of codimension sixteen in this moduli space (see \[PS, LP\]) and intersects the former transversally (these density results hold also in any universal family). Dominability by $`^2`$ holds for both types of K3 surfaces. This suggests that it might hold for all K3 surfaces so that our statements above would be valid without exception for projective (and, more generally, for compact Kähler) surfaces. Indeed, their density plus Brody’s lemma (\[Br\]) tells us that every K3 surface contains a non-trivial holomorphic image of $``$ and that the generic K3 surface, which is non-projective but remains Kähler, even contains such an image that is Zariski dense. We mention here that dominability by $`^2`$ can be shown for some non-elliptic K3 surfaces which are close to Kummer surfaces using an argument similar to that of section 6; for length considerations, we omit this non-elliptic case from this paper. However, we note that the statement equating dominability to the weaker condition of having a Zariski dense image of $``$ is quite false in the non-Kähler category, as is amply demonstrated by Inoue surfaces (see \[In0\] or \[BPV, V.19\]).
Observe that if there is a dominating map $`F:^2X`$, then there is also a holomorphic image of $``$ which is Zariski dense: First we may assume that the Jacobian of $`F`$ is non-zero at the origin. Defining $`h:^2`$ by $`h(z)=(\mathrm{sin}(2\pi z),\mathrm{sin}(2\pi z^2))`$, we see that $`h(n)=(0,0)`$ with corresponding tangent direction $`(2\pi ,4\pi n)`$ for each $`n`$. Taking $`Fh`$, we obtain a holomorphic image of $``$ with an infinite number of tangent directions at one point, which implies that the image is Zariski dense.
We say that an algebraic variety $`X`$ satisfies property C if every holomorphic image of $``$ in $`X`$ is algebraically degenerate; i.e., is not Zariski dense. Our first main result is that for algebraic surfaces of Kodaira dimension less than 2 and with the exceptions mentioned above, dominability by $`^2`$ is equivalent to the failure of property C. We will state only the main results in the projective category in this introduction for simplicity but will discuss fully the compact non-projective case and much of the quasi-projective case in this paper.
###### THEOREM 1.1
Let $`X`$ be a projective surface of Kodaira dimension less than $`2`$ and suppose that $`X`$ is not birational to a K3 surface which is either elliptic or Kummer. Then $`X`$ is dominable by $`^2`$ if and only if it does not satisfy property C. Equivalently, there is a dominating holomorphic map $`F:^2X`$ if and only if there is a holomorphic image of $``$ in $`X`$ which is Zariski dense.
By a recent result of the second named author, this theorem is also true for a projective surface of Kodaira dimension 2, which is the maximum Kodaira dimension for surfaces. As previously mentioned, a surface of Kodaira dimension 2 is not dominable by $`^2`$ \[Ko1\]; indeed, a surface of Kodaira dimension 2 is precisely a surface which admits a possibly degenerate hyperbolic volume form. Thus in the case of Kodaira dimension 2, theorem 1.1 can be established by showing that such a surface satisfies property C. The question of whether a variety of maximum Kodaira dimension satisfies property C was first raised explicitly by Serge Lang \[Lang\].
In the following theorem we give, again modulo the above mentioned exceptions, a classification of projective surfaces which are dominable by $`^2`$ and hence a classification of projective surfaces of Kodaira dimension less than 2 which fail to satisfy property C. We will do this in terms of the Kodaira dimension and the fundamental group, both of which are invariant under birational maps.
###### THEOREM 1.2
A projective surface $`X`$ not birationally equivalent to a K3 surface is dominable by $`^2`$ if and only if it has Kodaira dimension less than two and its fundamental group is a finite extension of an abelian group (of even rank four or less). If $`\kappa (X)=\mathrm{}`$, then the fundamental group condition can be replaced by the simpler condition of non-existence of more than one linearly independent holomorphic one-form. If $`\kappa (X)=0`$ and $`X`$ is not birationally equivalent to a K3 surface, then $`X`$ is dominable by $`^2`$. If $`X`$ is birationally equivalent to an elliptic K3 surface or to a Kummer K3 surface, then $`X`$ is dominable by $`^2`$.
As with theorem 1.1, this theorem fails if we include compact non-Kähler surfaces (even after simple minded modification of this theorem). For instance, the Kodaira surfaces are dominable by $`^2`$ but their fundamental groups are not finite extensions of abelian groups (\[Ko4\]). But this theorem remains valid in the Kähler category, thanks, for example, to Kodaira’s result that all Kähler surfaces are deformations of projective surfaces (\[Ko2\], \[Ko3\]).
More general versions of theorem 1.1 and theorem 1.2 for compact complex surfaces will be given at the end of section 4.
In the quasi-projective category, we also prove the analogue of theorem 1.1 modulo the same exceptions mentioned in the beginning, following mainly the work of Kawamata \[K1\] and M. Miyanishi \[M\]. In this setting, the analogue of the fundamental group characterization requires the study of a new but very natural class of objects of complex dimensional one that are related to orbifolds. As for explicit examples, we will work out theorems 1.1 and an analogue of 1.2 for the complement of a reduced curve $`C`$ in $`^2`$ in the case when $`C`$ is normal crossing, where we show that dominability is characterized by $`\mathrm{deg}C3`$, and for the overlapping case in which $`C`$ is either not a rational curve of high degree or has at most one singular point. Here, the most fascinating and revealing example is the case in which $`C`$ is a non-singular cubic curve, whose complement is a noncompact analogue of a K3 surface. The question of the dominability of the complement of a non-singular cubic was discussed by Bernard Shiffman at MSRI in 1996, and the positive resolution of this problem served as the first result in and inspiration for this paper.
The key tools we introduce here for constructing dominating maps are the mapping theorems we establish via a combination of complex geometry and analysis. One of these theorems utilizes Kodaira’s theory of Jacobian fibrations to deal with general elliptic fibrations (see section 3). Other such theorems construct the required self-maps of $`^2`$ directly via complex analysis to deal with $`^{}`$-fibrations, abelian and Kummer surfaces.
In particular, the constructions in sections 4 and 6 show that given any complex 2-torus and any finite set of points in this torus, there is an open set containing this finite set and a dominating map from $`^2`$ into the complement of the open set. This should be compared with \[Gr\] in which it was claimed that the complement of any open set in a simple complex torus is Kobayashi hyperbolic (a complex torus is simple if it has no nontrivial complex subtori). There is no contradiction because it was later realized that the proof given in \[Gr\] is incorrect since the topological closure of a complex one-parameter group need not be a complex torus. Despite this, the validity of this claim appears to have been an open question until the current paper, which shows the claim to be false in dimension 2. The $`n`$-dimensional analogue of our result is given in \[Bu\].
Many of the tools and results we develop may be of interest to other areas of mathematics besides complex analysis and holomorphic geometry, especially to Diophantine (arithmetic) geometry in view of the connection between the transcendental holomorphic properties and arithmetic properties of algebraic varieties. For example, the important technique of constructing sections of elliptic fibrations, which is very difficult to achieve in the algebro-geometric category but certainly useful in arithmetic and algebraic geometry, turns out to be quite natural and relatively easy to do in the holomorphic category. Also, we undertake a global study, from the viewpoint of holomorphic geometry, of the monodromy action on the fundamental group of an elliptic fibration. Needless to say, without the deep and beautiful contributions of Kodaira on complex analytic surfaces, we would not be able to go much beyond dealing with some special examples, as is the case with much of the scarce literature on the subject. However, we have not avoided, due to the nature of this joint paper, giving elementary lemmas and proofs while avoiding the unnecessary full force of Kodaira’s theory on elliptic fibrations, especially as we deal with fibrations over curves that are not necessarily quasiprojective.
The paper is organized as follows. Section 2 introduces some basic birational invariants and general notation and provides a list of the classification of projective surfaces. Section 3 deals with projective surfaces not of zero Kodaira dimension and solves the dominability problem completely for elliptic fibrations, including the non-algebraic ones. Section 4 deals with the remaining projective and compact complex cases while section 5 deal with the non-compact algebraic surfaces. Section 6 goes beyond these theorems to deal with algebraic surfaces minus small open balls.
We are very grateful to Bernard Shiffman for posing the question which motivated and inspired this paper and for his constant encouragement during its preparation.
## 2 Classification of algebraic surfaces
In this section we will first introduce some basic invariants in the (logarithmic) classification theory of algebraic varieties (see \[Ii\] for more details, also compare with \[Ue\]). Then we will provide a list of the birational classification of projective surfaces and discuss briefly the dominability problem in the quasi-projective category. Finally, we will introduce the more general category of compactifiable complex manifolds and a basic invariant which distinguishes the algebraic case in dimension two.
Let $`\overline{X}`$ be a complex manifold with a normal crossing divisor $`D`$. This means that around any point $`q`$ of $`\overline{X}`$, there exist a local coordinate $`(z_1,\mathrm{},z_n)`$ centered at $`q`$ such that, for some $`rn`$, $`D`$ is defined by $`z_1z_2\mathrm{}z_r=0`$ in this coordinate neighborhood. If all the components of $`D`$ are smooth, then $`D`$ is called a simple normal crossing divisor. Following Iitaka (\[Ii\]), we define the logarithmic cotangent sheaf $`\mathrm{\Omega }_{\overline{X}}(\mathrm{log}D)`$ as the locally free subsheaf of the sheaf of meromorphic 1-forms, whose restriction to $`X=\overline{X}D`$ is identical to $`\mathrm{\Omega }_X`$ and whose localization at any point $`qD`$ is given by
$$\mathrm{\Omega }_{\overline{X}}(\mathrm{log}D)=\underset{i=1}{\overset{r}{}}𝒪_{\overline{X},q}\frac{dz_i}{z_i}+\underset{j=r+1}{\overset{n}{}}𝒪_{\overline{X},q}dz_j,$$
where the local coordinates $`z_1,\mathrm{},z_n`$ around $`q`$ are chosen as before. Its dual, the logarithmic tangent sheaf $`T_{\overline{X}}(\mathrm{log}D)`$, is a locally free subsheaf of $`T_{\overline{X}}`$. We will follow a general abuse of notation and use the same notation to denote both a locally free sheaf and a vector bundle.
By an algebraic variety in this paper, we mean a complex analytic space $`X_0`$ such that $`X_0`$ has an algebraic structure in the following sense: $`X_0`$ is covered by a finite number of neighborhoods, each of which is isomorphic to a closed analytic subspace of a complex vector space defined by polynomial equations and which piece together with rational coordinate transformations. A proper birational map from $`X_0`$ to another variety $`X_1`$ is, by the graph definition, an algebraic subvariety of $`X_0\times X_1`$ which projects generically one-to-one onto each factor as a proper morphism. If such a map exists, we say that the two varieties are properly birational. This notion corresponds to that of a bimeromorphic map in the holomorphic context. Two algebraic varieties are said to be birationally equivalent if they have isomorphic rational function fields; or equivalently, if they have birational compactifications. Hironaka’s resolution of singularity theorem \[Hi\] (an elementary proof of which can be found in \[BM\]) implies that given any algebraic variety $`X_0`$, there is a smooth projective variety $`\overline{X}`$ with a simple normal crossing divisor $`D`$ such that $`X=\overline{X}D`$ is properly birational to $`X_0`$. If $`X_0`$ is smooth, then we can even take $`X`$ to be $`X_0`$ so that $`X_0`$ can be compactified by adding a simple normal crossing boundary divisor. In this paper, a surface will mean a complex two dimensional manifold while a curve that is not explicitly a subvariety (or a subscheme) will mean a (not necessarily quasi-projective) complex one-dimensional manifold. All surfaces and curves are assumed to be connected. In particular, every algebraic surface is isomorphic to the complement of a finite set of transversely intersecting smooth curves without triple intersection in some projective surface. We will use the Enriques-Kodaira classification of compact surfaces to simplify our problem for surfaces.
One of the most important invariants under proper birational maps is the (logarithmic) Kodaira dimension. Let $`X_0`$, $`X`$, $`\overline{X}`$, and $`D`$ be as above, and let $`K_{\overline{X}}=det_{}(T_{\overline{X}}^{})`$ where $`T_{\overline{X}}^{}`$ is the complex cotangent bundle of $`\overline{X}`$. The (holomorphic) line bundle $`K_{\overline{X}}`$ is called the canonical bundle of $`\overline{X}`$. Identifying a line bundle and its sheaf of holomorphic sections, we define a new line bundle $`K=K_{\overline{X}}(D)=K_{\overline{X}}𝒪(D)`$ corresponding to the sheaf of meromorphic sections of $`K_{\overline{X}}`$ which are holomorphic except for simple poles along $`D`$ (see Griffiths and Harris \[GH\] among many other standard references). In fact,
$$K=det\mathrm{\Omega }_{\overline{X}}(\mathrm{log}D).$$
This line bundle on $`\overline{X}`$ is called the logarithmic canonical bundle of $`X=\overline{X}D`$, or more specifically, of $`(\overline{X},D)`$. We will write tensor products of line bundles additively by a standard abuse of notation; for example, $`mK=K^m`$. Given a projective manifold $`\overline{Y}`$ and a birational morphism $`f:\overline{Y}\overline{X}`$ such that $`f^1(D)`$ is the same as a normal crossing divisor $`E`$ in $`\overline{Y}`$, then any section of $`mK`$ as a tensor power of rational 2-form on $`X`$ pulls back via $`f`$ to a section of $`mK_{\overline{Y}}(E)`$. Conversely, any section of $`mK_{\overline{Y}}(E)`$ pulls back (via $`f^1`$) to a section of $`mK`$ outside a codimension-two subset (the indeterminacy set of $`f^1`$ ), which therefore extends to a section of $`mK`$ by the classical extension theorem of Riemann. It follows that, for every positive integer $`m`$, $`h^0(mK):=dimH^0(mK)`$ is independent of the choice of $`\overline{X}`$ for $`X_0`$ and is a proper birational invariant of $`X_0`$. This allows us to introduce the following birational invariant of $`X_0`$.
###### DEFINITION 2.1
The Kodaira dimension of $`X_0`$ is defined as
$$\overline{\kappa }(X_0)=\underset{m\mathrm{}}{lim\; sup}\frac{\mathrm{log}h^0(mK)}{\mathrm{log}m}.$$
The simpler notation $`\kappa (X_0)`$ is used when $`X_0`$ is projective. The Riemann-Roch formula shows that $`\overline{\kappa }(X_0)`$ takes values in the set
$$\{\mathrm{},0,1,\mathrm{},dimX_0\}.$$
By the same argument as that for $`h^0(mK)`$, we see that another proper birational invariant is given by the (logarithmic) irregularity of $`X_0`$ defined by
$$\overline{q}(X_0)=h^0(\mathrm{\Omega }_{\overline{X}}(\mathrm{log}D)).$$
If $`D=0`$, then $`\overline{q}(X_0)`$ is just the dimension of the space of global holomorphic one-forms $`q(X)=h^0(\mathrm{\Omega }_X)`$ on $`X`$.
If $`\overline{\kappa }(X_0)=dim(X_0)`$, then $`X_0`$ is called a variety of general type. A theorem of Carlson and Griffiths \[CG\] (see also Kodaira \[Ko1\]) says that $`X_0`$ cannot be dominated (even meromorphically) by $`^n`$ in this case. Hence for both theorem 1.1 and theorem 1.2, we need consider only those surfaces with Kodaira dimension less than 2.
A projective surface $`X`$ whose canonical bundle has non-negative intersection with (or, equivalently, non-negative degree when restricted to) any curve in $`X`$ is called minimal. We say that $`K_X`$ is nef (short for numerically effective) in this case. In general, we say that a line bundle $`L`$ on $`X`$ is nef if $`LC0`$ for any curve $`C`$ in $`X`$.
Every algebraic surface is either projective or admits a projective compactification by adding a set of smooth curves with at most normal crossing singularities. Moreover, the Enriques-Kodaira classification \[BPV, Ch. VI\] says that a projective surface admits a birational morphism (as a composition of blowing up smooth points) to one of the following.
* A surface of general type: $`\kappa =2`$.
* $`^2`$ or a ruled surface over a curve $`C`$ of genus $`g=h^0(\mathrm{\Omega }_C)`$ (that is, a holomorphic $`^1`$ bundle over $`C`$). The latter is birationally equivalent to $`^1\times C`$. Here, $`\kappa =\mathrm{}`$.
* An abelian surface (a projective torus given by $`^2`$/a lattice). Here, $`\kappa =0`$.
* A K3 surface (a simply connected surface with trivial canonical bundle). $`\kappa =0`$.
* A minimal surface with the structure of an elliptic fibration (see section 3.2).
Here $`\kappa `$ can be $`0`$, $`1`$, or $`\mathrm{}`$.
The characteristic property of the surfaces listed above is the absence of $`(1)`$-curves. A $`(1)`$-curve is a smooth rational curve (image of $`^1`$) in a surface with self-intersection $`1`$, i.e. whose normal bundle has degree $`1`$. From Castelnuovo’s criterion \[BPV, III4.1\], a $`(1)`$ curve is always the blow-up of a (smooth) point on a surface. A simple argument (via the linear independence of the total transform of blown up $`(1)`$-curves in $`H_1`$) shows that, given any projective surface with $`\kappa <2`$, one can always reach one of the surfaces listed above by blowing down $`(1)`$-curves a finite number of times. It is a standard fact that a projective surface with $`\kappa 0`$ is minimal if and only if it does not have any $`(1)`$-curve, and that such a surface is the unique one in its birational class having this property.
Let $`X_0`$ be an algebraic surface having a compactification $`\overline{X}_0`$ which is birational to one of the model surfaces listed above, say $`\overline{X}`$. There is a maximum Zariski open subset $`U`$ of $`X_0`$ that is properly birational to the complement of a reduced divisor $`C`$ and a finite set $`T^{}`$ of points in $`\overline{X}`$. Now the indeterminacy set of this proper birational map from $`X=\overline{X}\{CT\}`$ to $`UX_0`$ must consist of a finite set of points. So to produce a dominating map from $`^2`$ to $`X_0`$, it suffices to produce, for each finite set of points $`T`$ in $`\overline{X}`$, a dominating map from $`^2`$ into the complement of $`T`$ in $`\overline{X}C`$. Nevertheless, $`\overline{X}C`$ may not be dominable by $`^2`$ when $`X_0`$ is dominable by $`^2`$; for example, a point on $`X_0`$ may correspond to an infinitely near point on $`\overline{X}`$ over a point of $`C`$. However, if we think of $`X_0`$ as an open subset of the space of infinitely near points of $`\overline{X}`$, then we can recover the equivalence in dominability through the above procedure (see section 5).
Although we have chosen to introduce and state our results so far in the algebraic category for simplicity, we will in fact deal with a more general class of surfaces in the next two sections: the class of compactifiable surfaces. These are Zariski open subsets of compact complex surfaces and the invariants $`\overline{\kappa }`$ and $`\overline{q}`$ carry over to them verbatim as they are defined by compactifications with normal crossing divisors, which exist by complex surface theory. If a surface $`X`$ is compact, the transcendency degree $`a(X)`$ of the field of meromorphic functions on $`X`$ is, by definition, a bimeromorphic invariant and is called the algebraic dimension.
## 3 Compact surfaces with $`\kappa 0`$ and $`a0`$
In this section we solve the $`^2`$ dominability problem for compact surfaces whose Kodaira dimension and algebraic dimension are both non-zero. The bulk of this section is devoted to the case of elliptic fibrations, which we treat completely, including all the noncompact cases. In particular, we solve our problem for every projective surface that is birational to a minimal one listed in (1) and (4) above. Cases (2) and (3) will be discussed in section 5.
### 3.1 Projective surfaces with Kodaira dimension $`\mathrm{}`$
Since any $`^1`$-bundle over a curve $`C`$ is birational to the trivial $`^1`$-bundle over $`C`$ and since $`^2`$ is birational to $`^1\times ^1`$, any projective surface $`X`$ with $`\overline{\kappa }(X)=\mathrm{}`$ is birational to a surface $`Y`$ which is a trivial $`^1`$-bundle over a curve $`C`$ of genus $`g:=h^0(\mathrm{\Omega }_C)`$. In the case where $`C`$ is of genus $`g>1`$, any holomorphic image of $``$ in $`Y`$ must lie in a fiber of the bundle since $`C`$ is hyperbolic. Hence $`X`$ satisfies property C and so cannot be dominated by $`^2`$. In the case $`Y`$ is a $`^1`$ bundle over an elliptic curve or over $`^1`$, one can easily construct a dominating map from $`^1\times ^1`$ and hence from $`^2`$ which respects the bundle structure (even algebraically in the latter case). In fact, by composing with the map
$$(\pi ^1,h\pi ^2):^2^2$$
(3.1)
where $`h:`$ is holomorphic with prescribed zeros (which we can do by Weierstrass’ theorem) and $`\pi ^1,\pi ^2`$ are the respective projections, we can arrange to have the dominating map miss any finite subset in $`Y`$. Choosing this finite subset to be the set of indeterminacies of the birational map from $`Y`$ to $`X`$, this dominating map lifts to give a dominating map into $`X`$. Since $`^1`$ admits no holomorphic differentials and is simply connected, we obtain, respectively,
$$q(X)=q(Y)=q(C)=g\text{and}\pi _1(X)=\pi _1(Y)=\pi _1(C).$$
Coupling this with the fact that the fundamental group of a curve of genus greater than 1 is not a finite extension of an abelian group gives us the following.
###### THEOREM 3.1
If $`X`$ is a projective surface with $`\kappa (X)=\mathrm{}`$, then the following are equivalent.
* $`X`$ is dominable by $`^2`$.
* $`q(X):=h^0(\mathrm{\Omega }_X)<2`$.
* $`X`$ admits a Zariski dense holomorphic image of $``$.
* $`\pi _1(X)`$ is a finite extension of an abelian group.
### 3.2 Elliptic fibrations
If $`X`$ is any compact non-projective surface with $`a(X)0`$, then $`X`$ is an elliptic surface by \[Ko2\]. Also, if $`X`$ is projective and $`\kappa (X)=1`$, then $`X`$ is again an elliptic surface by classification. Hence the only remaining cases of $`\kappa 0`$ and $`a0`$ are elliptic surfaces. In this section we resolve completely the case of elliptic surfaces.
###### DEFINITION 3.2
An elliptic fibration is a proper holomorphic map from a surface to a curve whose general fiber is an elliptic curve, i.e., a curve of genus one. Such a surface is called an elliptic surface. An elliptic fibration is called relatively minimal if there are no $`(1)`$-curves on any fiber.
Note that an elliptic fibration structure on a minimal surface must be relatively minimal.
Let $`f:XC`$ be a fibration (i.e. a proper holomorphic map with connected fibers) between complex manifolds $`X`$ and $`C`$. If $`f^{}:X^{}C^{}`$ is another map where $`C^{}C`$, then a map $`h:X^{}X`$ is called fiber-preserving if $`fh=f^{}`$. If rank$`(df)=\text{dim}C`$ at every point on a fiber $`X_s=f^1(s)`$, then $`X_s`$ must be smooth by the implicit function theorem. If rank $`(df)<\text{dim}C`$ somewhere on $`X_s`$, then $`X_s`$ is called a singular fiber. Outside the singular fibers, all fibers are diffeomorphic by Ehresmann’s theorem.
In the case $`f`$ is a fibration of a surface $`X`$ over a curve $`C`$, then each fiber, as a subscheme via the structure sheaf from $`C`$, is naturally an effective divisor on $`X`$ as follows. We write $`X_s=n_iC_i`$ where each $`C_i`$ is the $`i`$-th component of the fiber $`(X_s)_{red}`$ (without the scheme structure) and where $`n_i1`$ is the vanishing order of $`df`$ for a generic point on $`C_i`$. The positive integer coefficient $`n_i`$ is called the multiplicity of the $`i`$-th component. The multiplicity of a fiber $`X_s=n_iC_i`$ is defined as the greatest common divisor $`n_s`$ of $`\{n_i\}`$. A fiber $`X_s`$ with $`n_s>1`$ is called a multiple fiber. A smooth fiber is then a fiber of multiplicity one having only one component. The singular fibers form a discrete set in $`X`$ by analyticity. We will assume this setup for $`X`$ and $`C`$ from now on.
Let $`\alpha :\stackrel{~}{C}C`$ be a finite proper morphism. The ramification index at a point $`\stackrel{~}{s}\stackrel{~}{C}`$ is defined as the vanishing order of $`d\alpha `$ at $`\stackrel{~}{s}`$ plus one. Suppose $`\alpha `$ has ramification index $`n_s`$ at every point above $`sC`$ and suppose that this is true for every $`sC`$. Then, according to \[BPV, III, Theorem 9.1\], pulling back the fibration via this ramified cover yields an unramified covering $`\stackrel{~}{X}`$ over $`X`$. Also, the resulting fibration $`\stackrel{~}{X}\stackrel{~}{C}`$ no longer has any multiple fibers. Such a ramified covering $`\stackrel{~}{C}`$ is called an orbifold covering of $`C`$ with the given branched (orbifold) structure on $`C`$. More generally we have:
###### DEFINITION 3.3
Given a curve $`C`$ with an assignment of a positive integer $`n_s`$ for each $`sC`$ such that the set $`S=\{sC|n_s>1\}`$ is discrete in $`C`$, define $`D={\displaystyle \underset{n_s>1}{}}\left(1{\displaystyle \frac{1}{n_s}}\right)s`$. Suppose $`\alpha :\stackrel{~}{C}C`$ is a holomorphic map such that $`\alpha :\stackrel{~}{C}\alpha ^1(S)CS`$ is an unramified covering and such that, for each point $`sS`$, every point on $`\stackrel{~}{C}`$ above $`s`$ has ramification index $`n_s`$. Then $`\stackrel{~}{C}`$ is called an orbifold covering of the orbifold $`(C,D)`$. If also $`\stackrel{~}{C}`$ is simply connected, then $`\stackrel{~}{C}`$ is called a uniformizing orbifold covering. A fibration over $`C`$ defines a natural (branched) orbifold structure $`D`$ on $`C`$ by assigning $`n_s`$ to be the multiplicity of the fiber at $`s`$ of the fibration.
Therefore, we have the following:
###### PROPOSITION 3.4
Let $`X`$ be a fibration over $`C`$. Let $`n_s`$ denote the multiplicity of the fiber $`X_s`$ for every point $`sC`$, thus endowing $`C`$ with an orbifold structure $`D`$ as above. Let $`\stackrel{~}{C}`$ be an orbifold covering of $`(C,D)`$. Then the pull back fibration $`\stackrel{~}{X}\stackrel{~}{C}`$ has no multiple fibers and $`\stackrel{~}{X}X`$ is an unramified holomorphic covering map.
#### 3.2.1 The Jacobian Fibration
We first begin with a preliminary discussion in the absolute case, the case where the base is just one point.
Let $`Z`$ be a one dimensional subscheme (or a curve) in a complex projective surface. The arithmetic genus of $`Z`$, defined by $`p_a(Z)=h^1(𝒪_Z):=\text{dim}_{}H^1(𝒪_Z)`$, is equal to the geometric genus when $`Z`$ is smooth. Assume now that $`Z`$ is an arbitrary fiber in an elliptic fibration. Since $`p_a`$ is an invariant in any algebraic family of curves (\[Ha, III, cor. 9.13\]), we have $`p_a(Z)=1`$ and so $`H^1(𝒪_Z)=.`$ From the exponential exact sequence $`0𝒪𝒪^{}0`$, we construct the cohomology long exact sequence over $`Z`$ to deduce:
$$\begin{array}{ccccccccc}0& \hfill & H^1(Z,)& \stackrel{i}{}\hfill & H^1(𝒪_Z)& \hfill & H^1(𝒪_Z^{})& \stackrel{\delta }{}\hfill & H^2(Z,)0\\ & & ||& & ||& & ||& & ||\\ & & \text{}\text{-module}& & & & \text{Pic}(Z)& & \end{array}$$
Fact: (Let $`Z`$ be non-singular.) Pic$`(Z)`$ is naturally identified with the space of holomorphic line bundles over $`Z`$, which, in our case of $`p_a=1`$, is a 1-dimensional complex Lie group under tensor product. Every line bundle $`L`$ can be written as $`𝒪(E)`$ for some divisor $`E=a_is_i`$ ($`a_i,s_iZ`$) and $`\delta (L)=\mathrm{deg}E:=a_i`$.
###### DEFINITION 3.5
$`\text{Pic}^0(Z):=\mathrm{ker}\delta `$ is the subgroup of Pic$`(Z)`$ of line bundles $`L`$ with trivial first Chern class $`c_1(L):=\mathrm{deg}(L)`$.
If $`Z`$ is a smooth elliptic curve with a base point $`\sigma `$, we can construct a group homomorphism from $`Z`$ to Pic$`{}_{}{}^{0}(Z)`$ by the map
$$xZ\stackrel{f}{}𝒪(x\sigma )\text{Pic}^0(Z).$$
###### LEMMA 3.6
The map $`f`$ is holomorphic, one-to-one and hence onto.
Proof: As $`f`$ is holomorphic by construction, we need to prove only that it is one-to-one. Assume not, so that $`𝒪(x\sigma )=𝒪(x^{}\sigma )`$ where $`xx^{}`$. Then $`𝒪(xx^{})`$ corresponds to the trivial line bundle over $`Z`$ and so $`Z`$ has a rational function with a simple pole at $`x^{}`$ and a simple zero at $`x`$. This gives a 1-1 and hence surjective holomorphic map from $`Z`$, which has genus $`1`$, to $`^1`$, which has genus $`0`$. This is a contradiction.
Note: $`\text{Pic}^0(Z)=H^1(𝒪_Z)/i(H^1(Z,)).`$
We now return to the case in which the base is a curve.
Given an elliptic fibration $`f:XC`$ without multiple fibers, one can construct a relative version of $`\text{Pic}^0`$ as follows (see \[BPV, p. 153\]). We first form the $`𝒪_C`$ module
$$𝒥ac(f)=f_1(𝒪_X)/f_1$$
over $`C`$. Since $`p_a(X_s)=1`$ for every fiber, it follows that $`f_1(𝒪_X)`$ is locally free of rank $`1`$ (by a well known theorem of Grauert) and hence is the sheaf of sections of a line bundle $`L`$ over $`C`$. Hence $`𝒥ac(f)`$ corresponds to the sheaf of sections of
$$\text{Jac}(f):=L/f_1,$$
which is a holomorphic fibration of complex Lie groups with a zero section (see \[Ko2\], compare also \[BPV, V.9\]). Note that when $`X_s`$ is smooth elliptic, $`(f_1)_s=H^1(X_s,)`$ which embeds in $`L_s=H^1(𝒪_{X_s})=`$. So Jac$`(f)_s=\text{Pic}^0(X_s)`$. Note also that Jac$`(f)`$ is a holomorphic quotient of a line bundle $`L`$ over $`C`$.
We have the following theorem from Kodaira \[Ko2\] (see \[BPV, V9.1\]).
###### PROPOSITION 3.7
Let $`f:XC`$ be a relatively minimal elliptic fibration over a curve $`C`$ with a holomorphic section $`\sigma :CX`$. Let $`X_\sigma ^{}`$ consist of all irreducible components of fibers $`X_s`$ not meeting $`\sigma (C)`$, and let $`X^\sigma =XX_\sigma ^{}`$. Then there is a canonical fiber-preserving isomorphism $`h`$ from Jac$`(f)`$ onto $`X^\sigma `$ mapping the zero-section in Jac$`(f)`$ onto $`\sigma (C)`$.
Hence it is useful to construct holomorphic sections of elliptic fibrations for which we develop the following key lemma.
###### LEMMA 3.8
Given a relatively minimal elliptic fibration $`f:XC`$ without multiple fibers, assume $`C`$ is non-compact. Then $`f`$ has a holomorphic section. Furthermore, given a countable subset $`T`$ of $`X`$ whose image $`f(T)`$ is discrete in $`C`$, the section can be chosen to avoid $`T`$.
Proof: From Kodaira’s table of non-multiple singular fibers (\[Ko2\] or \[BPV, Table 3 p. 150\]), we see that every fiber which is not multiple in a relatively minimal elliptic fibration has a component of multiplicity one. So, every point on $`C`$ admits a neighborhood with a section. We now choose a locally finite good covering of $`C`$ by open sets $`U_1,U_2,\mathrm{},`$ with sections $`\tau _1,\tau _2,\mathrm{}`$ of $`f|_{U_1},f|_{U_2},\mathrm{},`$ respectively. We may further stipulate that there are no singular fibers on the intersection of any two $`U_j`$’s.
Let $`L=f_1𝒪_X`$, which is a holomorphically trivial line bundle over $`C`$ since $`C`$ is Stein. Let $`UC`$ be open and $`\tau ^{}H^0(U,L)`$ a section. If $`\tau `$ is a section of $`f|_U`$, then we can form the section $`\tau +\tau ^{}`$ of $`f|_U`$ by proposition 3.7. By the same proposition and the fact that all fibers are elliptic curves over $`U_iU_j`$, there is a section $`\tau _{ij}^{}H^0(U_iU_j,L)`$ such that $`\tau _i+\tau _{ij}^{}=\tau _j`$ on $`U_iU_j`$.
As $`\{\tau _{ij}^{}\}`$ satisfies the cocycle condition, so does $`\{\tau _{ij}^{}\}`$. By the solution to the classical additive Cousin problem (or from the fact that $`H^1(\{U_i\},L)=H^1(C,L)=H^1(C,𝒪)=0`$ by Leray’s theorem, Dolbeault’s isomorphism, and the fact that $`C`$ is Stein), one can find holomorphic sections $`\tau _i^{}H^0(U_i,L)`$ such that $`\tau _i^{}\tau _{ij}^{}=\tau _j^{}`$. Then $`\tau _i+\tau _i^{}=\tau _j+\tau _j^{}`$ on $`U_iU_j`$ for all $`i,j`$. This gives rise to a global section of $`f:XC`$.
Given such a global section, proposition 3.7 gives a fiber-preserving dominating map $`F:LX`$ where $`F^1(x)L`$ is at most a countable discrete set for all $`x`$ in $`X`$. Hence $`F^1(T)`$ is also a countable set and is supported on the fibers of $`L`$ over $`f(T)`$. For each $`sf(T)`$, therefore, we may choose a point $`q_s`$ in $`LT`$. As $`L`$ is isomorphic to the trivial line bundle ($`C`$ being non-compact), the classical interpolation theorems of Mittag-Leffler and of Weierstrass give us a holomorphic section $`\sigma `$ of $`L`$ with the prescribed value $`q_s`$ for all $`sT`$. But then $`F\sigma `$ is a section of $`f`$ which avoids $`T`$. This completes the proof.
#### 3.2.2 Theorem 1.1 in the case of elliptic fibrations
###### THEOREM 3.9
Let $`f:XC`$ be a relatively minimal elliptic fibration with a finite number of multiple fibers. Assume that $`C`$ is a Zariski open subset of a projective curve $`\overline{C}`$. Let $`n_s`$ be the multiplicity of the fiber $`X_s`$. Then the following are equivalent.
* $`X`$ is dominable by $`^2`$.
* $`\chi :=22g(\overline{C})\mathrm{\#}(\overline{C}C){\displaystyle \underset{n_s2}{}}(1{\displaystyle \frac{1}{n_s}})0`$.
* There exists a holomorphic map of $``$ to $`X`$ whose image is Zariski dense.
Remark 1: $`\chi =\chi (C,D)`$ is the orbifold Euler characteristic of $`(C,D)`$. It can be written as $`\chi (C,D)=22g(\overline{C}){\displaystyle \underset{s\overline{C}}{}}(1{\displaystyle \frac{1}{n_s}})`$ if we set $`n_s=\mathrm{}`$ for $`s\overline{C}C`$ (where $`\frac{1}{\mathrm{}}=0`$). Hence, if we complete the $``$-divisor $`D={\displaystyle \underset{sC}{}}(1{\displaystyle \frac{1}{n_s}})s`$ to $`\overline{D}={\displaystyle \underset{s\overline{C}}{}}(1{\displaystyle \frac{1}{n_s}})s`$ on $`\overline{C}`$, then
$$\chi (C,D)=22g(\overline{C})\mathrm{deg}\overline{D}.$$
Proof of theorem: The pair $`(C,D)`$ defines an orbifold as given in definition 3.3. We will show that (a) holds if $`\chi (C,D)0`$ while property C holds for $`X`$ (that is, (c) fails to hold) if $`\chi (C,D)<0`$. This will conclude the proof.
From the classical uniformization theorem for orbifold Riemann surfaces (see, for example, \[FK, IV 9.12\]), $`(C,D)`$ has a uniformizing orbifold covering $`\stackrel{~}{C}`$ which is $`^1`$, $``$ or $`𝔻`$ according to $`\chi (C,D)>0`$, $`\chi (C,D)=0`$ or $`\chi (C,D)<0`$ respectively, unless $`\overline{C}=^1`$ and $`\overline{D}`$ has one or two components. In the latter (“unless”) case, we simply redefine $`C`$ to be the complement of the components of $`\overline{D}`$ in $`^1`$ and reset $`D`$ to be $`0`$, shrinking $`X`$ as a result. We can do this because it does not change the fact that $`\chi (C,D)0`$ and because once we show that the resulting $`X`$ is dominable by $`^2`$, the original $`X`$ is also.
By pulling back the fibration to $`\stackrel{~}{C}`$, we obtain a relatively minimal elliptic fibration $`Y`$ over $`\stackrel{~}{C}`$. Now, proposition 3.4 implies that the natural map from $`Y`$ to $`X`$ is an unramified covering. Hence any holomorphic map from $``$ to $`X`$ must lift to a holomorphic map to $`Y`$. It follows that if $`\stackrel{~}{C}=𝔻`$, then any such map must lift to a fiber and hence its image in $`X`$ must lie in a fiber. So, property C holds and $`X`$ cannot be dominated by $`^2`$ in this case.
It remains to show that $`X`$ is dominated by $`^2`$ in the case $`\stackrel{~}{C}=`$ or $`^1`$ to complete the proof of this theorem. Note that the latter case can be reduced to the former by simply removing a point from $`\stackrel{~}{C}`$. Hence, we may take $`\stackrel{~}{C}`$ to be $``$ which is non-compact. Lemma 3.8 now applies to give a section of the pullback fibration $`\stackrel{~}{f}:Y`$. By proposition 3.7, $`Y`$ is dominated by Jac$`(\stackrel{~}{f})`$ which in turn is dominated by a line bundle $`L`$ over $``$ by construction. Hence $`X`$ is dominated by $`L=^2`$ (since any line bundle over $``$ is holomorphically trivial) as required.
Now, let $`f^{}:X^{}C`$ be an arbitrary elliptic fibration. By contracting the $`(1)`$-curves on the fiber, we get a bimeromorphic map $`\alpha `$ from $`X^{}`$ to a surface $`X`$ having a relatively minimal elliptic fibration structure over $`C`$. As before, $`X`$ defines an orbifold structure $`D`$ on $`C`$. If $`X`$ has an infinite number of multiple fibers or if $`C`$ is not quasi-projective, then $`𝔻`$ is the universal covering of $`(C,D)`$ and conditions (a) and (c) of this theorem both fail for $`X`$. Otherwise the above theorem can be applied to conclude that conditions (a) and (c) are still equivalent for $`X`$. Let $`T`$ be the indeterminacy set of $`\alpha ^1`$. By examining the last paragraph of the above proof, we see that lemma 3.8 actually applies to give us a dominating map from the trivial line bundle $`L`$ over $`\stackrel{~}{C}`$ to $`X`$, and the zero-section of $`L`$ maps to a section of $`f`$ that avoids $`T`$. Composing with a self-map of $`L`$ given by a section of $`L`$ with prescribed zeros (just as in equation 3.1) then gives us a dominating map from $`L`$ to $`X`$ which avoids $`T`$. Hence, if $`X`$ is dominable by $`^2`$, then $`X^{}`$ is also. It is clear that $`X^{}`$ satisfies property C if $`X`$ does. Hence, we obtain the following, which covers theorem 1.1 in the case of elliptic surfaces.
###### THEOREM 3.10
Let $`f:XC`$ be an elliptic fibration. Then conditions (a) and (c) of theorem 3.9 above are equivalent for $`X`$; that is, dominability by $`^2`$ is equivalent to having a Zariski-dense holomorphic image of $``$.
Note that we do not require $`C`$ to be quasi-projective in this theorem.
#### 3.2.3 An algebro-geometric characterization
In this section, we will give, without proof, a characterization of dominability by $`^2`$ for a projective elliptic fibration in terms of familiar quantities in algebraic geometry and not involving the fundamental group. Unfortunately, the condition given is not straightforward nor does it seem very tractable. Hence, we will leave the proof (which is based on the simple fact that the saturation of the cotangent sheaf of the base, pulled back by the fibration map, includes the orbifold cotangent sheaf as a $``$ subsheaf) to the reader. We will deal only with the case of $`\kappa =1`$ since the other possibility of $`\kappa =0`$ contains the, so far, problematic K3 surfaces. However, all surfaces with $`\kappa =0`$ other than the K3’s are dominable by $`^2`$. We note that from the classification list in section 2, a surface with $`\kappa =1`$ is necessarily elliptic. Before the statement of the following proposition, recall that a vector sheaf is called big if it contains an ample subsheaf. Recall also that a divisor in a surface is nef if its intersection with any effective divisor is non-negative.
###### PROPOSITION 3.11
Let $`X`$ be a projective surface with $`\kappa (X)=1`$. Then $`X`$ is dominable by $`^2`$ if and only if there exists a nef and big divisor $`H`$ such that, for every nef divisor $`N`$ with $`K_XN=0`$, there exists a positive integer $`m`$ with $`\mathrm{S}^m\mathrm{\Omega }_X(HN)`$ big.
It is not difficult to extract a birational invariant out of this from $``$-subsheaves of the cotangent sheaf of such an elliptic surface; again we leave this to the interested reader.
In the remainder of this section, we give a more satisfactory and elementary characterization of dominability, now in terms of the fundamental group.
### 3.3 The fundamental group of an elliptic fibration
We begin with the remark that, except for our narrow focus on holomorphic geometry, most of the results we obtain in this section are not presumed to be new.
Let $`f:XC`$ be an elliptic fibration. Then the fibration determines a branched orbifold structure $`D`$ on $`C`$ as given in definition 3.3. Let $`C^{}`$ be the complement of the set of branch points in $`C`$. Then $`X^{}=f^1(C)`$ is an elliptic fibration defined by $`f^{}=f|_X^{}`$, which has no multiple fiber. Let $`X^{}`$ be the complement of the singular fibers in $`\overline{X}`$. Then $`f^{}=f|_X^{}`$ defines a smooth fibration over a curve $`C^{}C`$, and is therefore differentiably locally trivial by Ehresmann’s theorem. We have the following commutative diagram.
$$\begin{array}{ccccc}\hfill X^{}& & \hfill X^{}& & X\hfill \\ \hfill f^{}& & \hfill f^{}& & f\hfill \\ \hfill C^{}& & \hfill C^{}& & C\hfill \end{array}$$
(3.2)
We first observe the following trivial lemma for our consideration of $`\pi _1(X)`$. Throughout this section, all paths are assumed to be continuous.
###### LEMMA 3.12
Assume that we are given a real codimension two subset $`W`$ of $`X`$ and a path $`\nu :[0,1]X`$ such that $`\nu (0)`$ and $`\nu (1)`$ lies outside $`W`$. Then $`\nu `$ is homotopic to a path that avoids $`W`$ keeping the end points fixed.
Proof: We first impose a metric on $`X`$. Since $`[0,1]`$ is compact, there is an integer $`n`$ such that $`\nu ([(i1)/n,i/n])`$ is contained in a geodesically convex open ball $`B_i`$ for all $`i\{1,2,\mathrm{},n\}`$. Then the intersection of these balls are also geodesically convex and, in particular, connected. Now replace $`\nu (i/n)`$ by a point in $`B_iB_{i+1}W`$ for each integer $`i[1,n1]`$. Then replace $`\nu |_{[(i1)/n,i/n]}`$ by a path in $`B_iW`$ connecting $`\nu ((i1)/n)`$ with $`\nu (i/n)`$, for each integer $`i[1,n]`$. This is possible because the complement of $`W`$ in each of the open balls is connected as $`W`$ is of real codimension two in them. Since the balls are contractible and intersect in connected open sets, we see that the new path is homotopic to the original one fixing the end points but now avoids $`W`$.
If the path $`\nu `$ given above has the same end points, that is $`\nu (0)=\nu (1)`$, then we call $`\nu `$ a loop. We will often identify $`\nu `$ with its image.
For the next two propositions, we observe from Kodaira’s table of singular fibers (see \[BPV, V.7\]) that, for a fiber $`X_s`$ of an elliptic fibration (as a topological space or a simplicial complex), $`\pi _1(X_s)`$ is either $``$ (corresponding to a nonsingular elliptic curve), $``$ (corresponding to the (semi-)stable singular fibers), or the trivial group (corresponding to the other singular fibers).
###### PROPOSITION 3.13
Let $`f:XC`$ be an elliptic fibration. In the case $`C=^1`$, let $`X_{\mathrm{}}`$ be a multiple fiber if one exists. Assume $`f`$ has no multiple fibers except possibly for $`X_{\mathrm{}}`$ and that $`C`$ is simply connected. Then $`\pi _1(X)`$ is a quotient of $`\pi _1(X_s)`$ for every fiber outside $`X_{\mathrm{}}`$. In particular, $`\pi _1(X)`$ is abelian.
Proof: Since contracting $`(1)`$-curves does not change the fundamental group, we may assume without loss of generality that $`f`$ is relatively minimal. Let $`X_s`$ be an arbitrary fiber. Being a CW-subcomplex of $`X`$, it is a deformation retract of a small neighborhood $`U`$ which we may assume to contain a smooth fiber $`X_s^{}`$ nearby. Since $`X`$ is path connected, we can choose any base point in considering its fundamental group. Fix then a base point $`qX_s^{}`$ and a loop $`Q`$ with this base point. We will show that $`Q`$ is pointed homotopy equivalent in $`X`$ to a loop in $`X_s^{}U`$. The theorem then follows as $`X_s`$ is a deformation retract of $`U`$.
Since the singular fibers form a real codimension two subset, we can modify $`Q`$ to avoid them up to pointed homotopy equivalence by lemma 3.12 above. In the case $`C=^1`$ but $`X_{\mathrm{}}`$ is not already given, let $`X_{\mathrm{}}`$ be a fiber outside $`U`$ and $`Q`$. Since every homotopy (of $`Q`$) in $`XX_{\mathrm{}}`$ is also one in $`X`$, we may safely replace $`X`$ by $`XX_{\mathrm{}}`$ so that $`C`$ becomes contractible in this case. Hence, we may assume in all cases that $`C`$ is contractible and that $`Q`$ is a loop in $`X^{}`$, the complement of the singular fibers in $`X`$. So lemma 3.8 and proposition 3.7 apply to give an isomorphism from Jac$`(f)`$ to $`X`$ with parts of the singular fibers complemented. Hence, we get, by construction of Jac$`(f)`$, a map $`\theta `$ from a holomorphically trivial line bundle $`L`$ over $`C`$ to $`X`$ which is an unramified covering above $`X^{}X`$. Fixing a point $`q_0\theta ^1(q)L_s^{}`$, we see that $`Q`$ can be lifted to a path $`\stackrel{~}{Q}`$ in $`L`$ from $`q_0`$ to a point $`q_1L_s^{}`$ by the theory of covering spaces. As $`C`$ is contractible, there is a homotopy retraction of $`L`$ to $`L_s^{}`$ which provides a pointed homotopy of $`\stackrel{~}{Q}`$ to a path in $`L_s^{}`$. Pushing down this homotopy (via $`\theta `$) to $`X`$ gives a pointed homotopy from $`Q`$ to a loop in $`X_s^{}`$ as required.
Looking back at the above proof, we see that we can reach the same conclusion by allowing $`X_s`$ to be a multiple fiber as long as $`C`$ is contractible and $`X`$ is free of other multiple fibers. This can be done by contracting the loop $`Q`$, as given in the proof, but only to the neighborhood $`U`$ of $`X_s`$ before homotoping to $`X_s`$ via the deformation retraction of $`U`$ to $`X_s`$. Of course, $`X_s`$ as stated in the theorem is no longer arbitrary in this case as it is a multiple fiber. If $`C=^1`$ and $`D`$ has two components (corresponding to $`X`$ having two multiple fibers), we can remove one of the components (corresponding to removing one multiple fiber from $`X`$) for the same conclusion. We recall that
$$D=\underset{sC}{}\left(1\frac{1}{n_s}\right)s$$
defines the orbifold structure on $`X`$ where $`n_s`$ is the multiplicity of the fiber $`X_s`$. Hence, we get a complement to the above proposition.
###### PROPOSITION 3.14
Let $`f:XC`$ be an elliptic fibration defining the orbifold structure $`D`$ on $`C`$. If $`C=^1`$ and $`D`$ has one or two components, or if $`C`$ is contractible and $`D`$ has one component, then $`\pi _1(X)`$ is a quotient of $`\pi _1(X_s)`$ for every component $`s`$ of $`D`$. Hence, $`\pi _1(X)`$ is abelian in these cases.
#### 3.3.1 Monodromy action as conjugation in the fundamental group
Although it is not absolutely necessary, some familiarity with the notion of monodromy and vanishing cycles used in geometry may be useful for reading this section.
Let the setup be as in diagram 3.2 and let $`X_r`$ be a non-singular fiber. Fix a base point $`q`$ in $`X_r`$ for all fundamental group considerations from now on. There is an action of $`\pi _1(X^{},q)`$ on $`\pi _1(X_r,q)`$ via the monodromy action which, in the case $`C^{}`$ is not $`^1`$, is just the conjugation action in $`\pi _1(X^{},q)`$. Indeed, in this case, we have the following exact sequence from the theory of fiber bundles (or from elementary covering space theory)
$$0\pi _1(X_r,q)\pi _1(X^{},q)\pi _1(C^{},r)0,$$
(3.3)
from which we deduce that the monodromy action is really an action of $`\pi _1(C^{},r)`$ on $`\pi _1(X_r,q)`$ since the latter is abelian.
In general, we will let $`H`$ denote the image of $`\pi _1(X_r,q)`$ in $`\pi _1(X,q)`$ under the inclusion of $`X_r`$ in $`X`$. It is easy to see that $`H`$ is a normal subgroup in $`\pi _1(X)`$ (by the definition of the monodromy action). In this paper, we will be mainly interested in the monodromy action on $`H`$. As opposed to the usual case of the monodromy action on the homology level, this action need not be trivial, unless we know, for example, that $`\pi _1(X,q)`$ is abelian. Hence, it is of interest for us to know how far $`\pi _1(X,q)`$ is from abelian.
With the same setup, suppose $`(C,D)`$ has a uniformizing orbifold cover $`\stackrel{~}{C}`$. This is the case unless $`\overline{C}=^1`$ and $`\overline{D}`$ has one or two components, again by the uniformization theorem (\[FK, IV 9.12\]). In the latter cases, proposition 3.14 and proposition 3.13 tell us that $`\pi _1(X)`$ is abelian so that the monodromy action on $`H`$ is trivial. In all other cases, let $`\stackrel{~}{f}:\stackrel{~}{X}\stackrel{~}{C}`$ be the pullback fibration. Proposition 3.4 implies that $`\stackrel{~}{X}`$ is an unramified cover over $`X`$. Let $`R`$ be the covering group and $`G=\pi _1(X)`$. From the theory of covering spaces, we know that $`G`$ is an extension of $`\pi _1(\stackrel{~}{X})`$ by $`R`$. Since $`\pi _1(X_r)`$ surjects to $`\pi _1(\stackrel{~}{X})\pi _1(X)`$ by proposition 3.4, we see that
$$H=\pi _1(\stackrel{~}{X}).$$
Note that $`R`$ is a quotient of $`\pi _1(C^{})`$ and hence also of $`\pi _1(C^{})`$, allowing us to identify the conjugation action of $`R`$ on $`H`$ with the monodromy action. Hence, we have the following exact sequence (which we can regard as a quotient of the exact sequence 3.3)
$$0HGR0.$$
(3.4)
The following proposition tells us that this monodromy action on $`H`$ via loops in $`X^{}`$, which induces the conjugation action of $`R`$ on $`H`$, depends only on the pointed homotopy class of the image of these loops in $`C`$. Hence, the monodromy action on $`H`$ is really an action by the group $`\pi _1(C)`$, which is a quotient of $`R`$. In particular, it tells us that the action is trivial when $`C`$ is simply connected. This is the closest analogue, on the level of $`\pi _1`$, of the fact that vanishing cycles are vanishing on the level of homology.
###### PROPOSITION 3.15
Let $`f:XC`$ be an elliptic fibration. Let $`X_r`$ be a non-singular fiber with a base point $`q`$. If $`\alpha `$, $`\beta `$ and $`\gamma `$ are loops based at $`q`$ with $`\alpha `$ in $`X_r`$, and $`f\beta `$ is pointed homotopic to $`f\gamma `$ in $`C`$, then $`\beta ^1\alpha \beta `$ is pointed homotopic in $`X`$ to $`\gamma ^1\alpha \gamma `$.
Proof: We may assume, via lemma 3.12, that $`\beta `$ and $`\gamma `$ lie in $`X^{}`$. Let $`h:[0,1]\times [0,1]C`$ be a pointed homotopy between $`f\beta `$ and $`f\gamma `$, which exists by assumption. Note that
$$(f\beta )(f\gamma )^1=h\left(([0,1]\times [0,1])\right)$$
as loops up to pointed homotopy equivalence, where $``$ means the oriented boundary. Our conclusion would follow if we show that the monodromy action of this latter loop, call it $`\mu `$, on $`\alpha `$ is trivial in $`\pi _1(X)`$.
By compactness of $`h([0,1]\times [0,1])`$, there is a partition $`\{0=a_0<a_1<\mathrm{}<a_n=1\}`$ of $`[0,1]`$ such that $`h([a_{i1},a_i]\times [a_{j1},a_j])`$ is contained in an open disk $`D_{ij}`$ containing at most one branch point and such that the loop
$$\mu _{ij}:=h\left(([a_{i1},a_i]\times [a_{j1},a_j])\right)$$
lies in $`X^{}`$, for all $`i,j\{1,2,\mathrm{},n\}`$. Since $`\pi _1\left(f^1(D_{ij})\right)`$ is abelian by proposition 3.14, the monodromy action of $`\mu _{ij}`$ on any pointed loop in the fiber is trivial in $`\pi _1(f^1(D_{ij}))`$, and hence in $`\pi _1(X)`$ as well, for all $`i,j\{1,2,\mathrm{},n\}`$. Our result now follows from the fact that the monodromy action of $`\mu `$ is just the sum of the monodromy action of the $`\mu _{ij}`$’s.
We can do a bit better when $`X`$ is compact.
###### LEMMA 3.16
With the setup as in the above proposition, assume further that either $`X`$ is compact or $`X`$ is a holomorphic fiber bundle over $`C`$. Then an integer $`m`$ exists, independent of $`\beta `$, such that $`\beta ^m`$ commutes with $`\alpha `$ in $`\pi _1(X)`$.
Proof: Let $`H`$ be the image of $`\pi _1(X_r)`$ in $`\pi _1(X)`$. We may assume, as before, that $`f`$ is relatively minimal.
If $`X`$ has a singular fiber, then $`H`$ is cyclic and hence the result follows from the fact that the automorphism group of a cyclic group is finite.
If $`X`$ is a holomorphic fiber bundle over $`C`$, then the monodromy actions can be realized as holomorphic automorphisms of the fiber. The group of such automorphisms is a finite cyclic extension of the group of lattice translations (this can be deduced easily or determined from the table in V.5 of \[BPV\] listing such groups). Hence every monodromy action up to a power is a translation on the fiber, which therefore leaves every element of $`\pi _1(X_r)`$ invariant.
If $`X`$ is compact and has no singular fibers, then it is a holomorphic fiber bundle by Kodaira’s theory of Jacobian fibrations. So the result follows by the last paragraph.
If $`X`$ is non-compact and $`f`$ is algebraic without singular fibers, then the conclusion of this lemma may no longer hold. Neverthless, we can embed $`X`$ in a projective surface $`\overline{X}`$, which is again elliptic. Deligne’s Invariant Subspace Theorem \[Del\] implies that elements in $`\pi _1(X_r)`$ which vanish in $`\pi _1(X)`$ are generated over $``$ by commutators of the form given by this lemma. But we can deduce this directly from the fact that the abelianization of $`\pi _1(\overline{X})`$ must have even rank so that either $`H`$ lies in the center of $`\pi _1(\overline{X})`$ (in the case when $`\overline{X}`$ is birational to an elliptic fiber bundle) or the commutator subgroup of $`\pi _1(\overline{X})`$ generates $`H`$ over $``$. In fact, Kodaira’s theory allows us to deduce a strong version of the Invariant Subspace Theorem (in the case of elliptic fibrations) which is valid even outside the algebraic category:
###### PROPOSITION 3.17
Let $`f:XC`$ be an elliptic fibration without singular fibers and such that $`C`$ is the complement of a discrete set in a quasi-projective curve. Let $`X_r`$ be a fiber. Then either $`f`$ is holomorphically locally trivial or $`\pi _1(X_r)=H_1(X_r)`$ is generated by the vanishing cycles — that is, by loops of the form $`\alpha ^1\beta ^1\alpha \beta `$ (naturally identified as elements of $`\pi _1(X_r)`$ via monodromy) in the notation of proposition 3.15, where $`\alpha `$ is a loop in $`X_r`$.
We remark that a weaker form of this proposition is in fact due to Kodaira and is disguised in the proof of theorem 11.7 in \[Ko2\]. We will follow his method, almost verbatim, in our proof.
Proof: We begin with some preliminaries concerning the period function $`z(s)`$, which takes values in the upper half plane. Recall that, as far as monodromy actions are concerned, we can identify $`\beta \pi _1(X)`$ with an element of $`\pi _1(C)`$, which we will denote again as $`\beta `$ by abuse of notation.
By theorem 7.1 and theorem 7.2 of \[Ko2\] (neither of which requires the additional assumption of that section concerning the compactification), we have a multivalued holomorphic period function $`z(s)`$ on $`C`$ with positive imaginary part such that, under the monodromy representation $`(\beta )`$SL$`(2,)`$ of $`\beta \pi _1(C)`$ as an automorphism of the lattice $`H_1(X_r)`$ with a fixed choice of basis, $`z(r)`$ transforms as
$$\beta _{}:z(r)\frac{az(r)+b}{cz(r)+d},\text{where}(\beta )=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\text{SL}(2,)$$
under our choice of basis. By definition, $`(1,z(s))`$ is the period defining the elliptic curve $`X_s`$ via analytic continuation of $`(1,z(r))`$, which is fixed by our choice of basis on $`H_1(X_r)`$ (see equation 7.3 in \[Ko2\]).
With a choice of basis over the point $`r`$ fixed, we can regard the period function $`z`$ as a single valued holomorphic function on the universal cover $`\stackrel{~}{C}`$. Also, we can naturally identify $`\pi _1(C)`$ with the covering transformation group of $`\stackrel{~}{C}`$ over $`C`$. Then we have (see equation 8.2 in \[Ko2\])
$$z(\beta (\xi ))=\beta _{}z(\xi )=\frac{az(\xi )+b}{cz(\xi )+d},\text{where}(\beta )=\left(\begin{array}{cc}a& b\\ c& d\end{array}\right)\text{and}\xi \stackrel{~}{C}.$$
Let $`M`$ denote the submodule of $`\pi _1(X_r)=H_1(X_r)=`$ generated by the vanishing cycles. After a suitable change of basis, we may assume that $`M=nm`$, where $`m`$ and $`n`$ are integers. If $`M`$ does not generate $`H_1(X_r)`$ over $``$, then either $`m`$ or $`n`$ must vanish. If $`m`$ vanishes, then we must have
$$(\beta )=\left(\begin{array}{cc}1& b_\beta \\ 0& 1\end{array}\right)(\text{for some}b_\beta ),$$
and therefore $`z(\beta (\xi ))=z(\xi )+b_\beta `$ for all $`\beta \pi _1(C)`$. Since the imaginary part of $`z(s)`$ is positive, exp$`[2\pi iz(s)]`$ defines a single valued holomorphic function on $`C`$ with modulus less than $`1`$. Hence, it must extend to a bounded holomorphic function on the compactification $`\overline{C}`$ of $`C`$ and therefore must be constant. It follows that $`z(s)`$ is constant and so the fibration is locally holomorphically trivial. If $`n`$ vanishes, then
$$(\beta )=\left(\begin{array}{cc}1& 0\\ c_\beta & 1\end{array}\right)(\text{for some}c_\beta ),$$
and therefore
$$1/z(\beta (\xi ))=1/z(\xi )+c_\beta .$$
Hence, considering exp$`[2\pi i/z(s)]`$ instead of exp$`[2\pi iz(s)]`$ gives us the same conclusion. This ends our proof.
In order to study $`G=\pi _1(X)`$, we need some information about its quotient $`R`$. This is fortunately a classical subject that we now turn to.
#### 3.3.2 Fuchsian groups versus elementary groups
In this section, we will collect some basic definitions and facts that we will need about Kleinian groups. We refer the reader to \[FK, IV.5-IV.9\] and \[Mas, I-V\] for more details.
Let $`(C,D)`$ be an orbifold and let $`\stackrel{~}{C}`$ be its uniformizing orbifold covering with covering group $`R`$ which acts holomorphically on $`\stackrel{~}{C}`$. Since $`\stackrel{~}{C}=^1,`$ or $`𝔻`$, all of which have natural embeddings into $`^1`$, $`R`$ can be identified as a subgroup of the group $`𝕄`$ of holomorphic automorphisms of $`^1`$, the group of Mobius transformations. So identified, $`R`$ becomes a Kleinian group; that is, a subgroup of $`𝕄`$ with a properly discontinuous action at some point, and hence in some maximum open subset $`\mathrm{\Omega }`$, of $`^1`$. The set of points $`\mathrm{\Lambda }=^1\mathrm{\Omega }`$ where $`R`$ does not act properly discontinuously is called the limit set of $`R`$.
An elementary group is a Kleinian group $`R`$ with no more than two points in its limit set. Such a group acts properly discontinuously on $`\mathrm{\Delta }^1`$, where $`\mathrm{\Delta }`$ is $`^1`$, $``$ or $`^{}`$.
By a Fuchsian group, we mean a Kleinian group $`R`$ with with a properly discontinuous action on some disk $`𝔻^1`$ such that $`𝔻/R`$ is quasi-projective; that is, $`𝔻`$ is the uniformizing orbifold covering of an orbifold $`(C,D)`$ where $`C=𝔻/R`$ is quasi-projective. If $`D`$ has finitely many components, then $`(X,D)`$ is known as a finite marked Riemann surface and $`R`$ is called basic. The limit set of a Fuchsian group necessarily contains the boundary of $`𝔻`$ (which characterizes Fuchsian groups of the first kind in the literature). It follows that a Fuchsian group cannot be an elementary group. We can also see this directly as follows.
###### LEMMA 3.18
An elementary Kleinian group is not a Fuchsian group.
Proof: Let $`R`$ be an elementary Kleinian group, then $`R`$ acts properly discontinuously on $`\mathrm{\Delta }=^1,`$ or $`^{}`$ as a subset of $`^1`$. If $`R`$ also acts on a disk $`𝔻^1`$, then the boundary of this disk with at most two points removed is contained in $`\mathrm{\Delta }`$. Since $`R`$ is properly discontinuous on $`\mathrm{\Delta }`$, and hence on this punctured boundary, $`𝔻/R`$ is not quasi-projective. Hence $`R`$ is not Fuchsian.
The following is a direct consequence of the uniformization theorem.
###### PROPOSITION 3.19
Let $`(C,D)`$ be a uniformizable orbifold where $`D`$ has a finite number of components. Let $`R`$ be the uniformizing orbifold covering group of $`(C,D)`$ properly regarded as a Kleinian group. Then $`\chi (C,D)<0`$ if and only if $`R`$ is a Fuchsian group while $`\chi (C,D)0`$ if and only if $`R`$ is an elementary group.
The reader is cautioned that lemma 3.18 is not a corollary of this since the definition of an elementary group is more general than that given in this proposition.
Concerning $`R`$ as an abstract group, proposition 3.19 and the basic theory of elementary Kleinian groups (see \[Mas, V.C and V.D\] or \[FK, IV 9.5\]) gives:
###### PROPOSITION 3.20
With the same setup as proposition 3.19, assume that $`\chi (C,D)0`$. Then there is a finite orbifold covering $`\stackrel{~}{C}`$ of $`(C,D)`$ such that $`\stackrel{~}{C}=^1`$, $`^{}`$ or an elliptic curve. In particular, $`R`$ is a finite extension of a free abelian group of rank at most two.
Quoting \[Mas, V.G.6\], using lemma 3.18 and proposition 3.20, we have:
###### PROPOSITION 3.21
Let $`R`$ be a Fuchsian group as defined above. Then $`R`$ is not a finite extension of an abelian group. Hence, $`R`$ is not isomorphic to an elementary group as an abstract group.
#### 3.3.3 The fundamental group characterization in theorem 1.2
Before stating the main theorem of this section, we need the following proposition from \[BPV, V.5\]. We first note from the same source that an elliptic fiber bundle over an elliptic curve is called a primary Kodaira surface if it is not Kähler. A non-trivial free quotient of such a surface by a finite group is called a secondary Kodaira surface. The fundamental group of such a surface is unfortunately not a finite extension of an abelian group, even though the surface is $`^2`$-dominable.
###### PROPOSITION 3.22
An elliptic fiber bundle over an elliptic curve is either a primary Kodaira surface, or a free and finite quotient of a compact complex 2-dimensional torus.
Armed with this, we are ready to tackle our second main theorem, theorem 1.2, in the case of elliptic fibrations. We will state a more general theorem:
###### THEOREM 3.23
Let $`f:XC`$ be an elliptic fibration with $`C`$ quasi-projective. Assume that $`X`$ is not bimeromorphic to a free and finite quotient of a primary Kodaira surface. Then $`X`$ is dominable by $`^2`$ if and only if $`\pi _1(X)`$ is a finite extension of an abelian group (of rank at most 4).
Proof: With the assumptions as in the theorem, we let $`G=\pi _1(X)`$ as before. By the same argument as that for theorem 3.10, we may assume, without loss of generality, that $`X`$ is relatively minimal by contracting the $`(1)`$-curves (as $`G`$ is unchanged in this process). If $`X`$ has an infinite number of multiple fibers, then the orbifold $`(C,D)`$ is uniformized by $`𝔻`$ and so $`X`$ is not dominable by $`^2`$. Proposition 3.21 tells us that $`R`$ is not isomorphic to a finite extension of an abelian group in this case. Hence, we may also assume that $`D`$ has only a finite number of components for the rest of the proof. Theorem 3.9 then applies and so it is sufficient to show that $`\chi (C,D)0`$ if and only if $`G`$ is a finite extension of an abelian group.
Assume that $`\chi (C,D)<0`$. If $`(C,D)`$ projectivizes to $`(^1,\overline{D})`$ (see the first remark after theorem 3.9 for the definition of $`\overline{D}`$), then $`\overline{D}`$ must have more than two components by the definition of $`\chi `$. Hence $`(C,D)`$ is uniformizable and we may apply proposition 3.19 and proposition 3.21 to conclude that the orbifold uniformizing group $`R`$ of $`(C,D)`$ is not a finite extension of an abelian group. But then neither is $`G`$ as $`R`$ is a quotient of $`G`$.
Conversely, assume $`\chi (C,D)0`$. If $`(C,D)`$ projectivizes to $`(^1,\overline{D})`$ and $`\overline{D}`$ has no more than two components, then $`G`$ is abelian by proposition 3.13 and proposition 3.14. Otherwise, $`(C,D)`$ is uniformizable and, with the notation as in section 3.3.1, the exact sequence 3.4 implies that $`G`$ is an extension of $`H`$ by $`R`$. Proposition 3.20 now applies to give a pull back elliptic fibration $`\widehat{f}:\widehat{X}\widehat{C}`$ without multiple fibers such that $`\widehat{X}`$ is a finite unramified covering of $`X`$ and such that $`\widehat{C}=^1`$, $`^{}`$ or an elliptic curve. We will consider each of these cases for $`\stackrel{~}{C}`$ separately. Note first that $`G=\pi _1(X)`$ (respectively $`R`$) is a finite extension of $`\widehat{G}=\pi _1(\widehat{X})`$ (respectively $`\widehat{R}`$) and that $`\widehat{H}=H`$. Replacing $`\widehat{C}`$ by a finite unramified covering of $`\widehat{C}`$, we may assume, thanks to lemma 3.16 and proposition 3.17, that $`H`$ lies in the center of $`\widehat{G}`$ (that is, the conjugation action of $`\widehat{G}`$ on $`H`$ is trivial).
In the case when $`\widehat{C}=^1`$, proposition 3.13 implies that $`\widehat{G}`$ is a quotient of a free abelian group of rank two. Hence $`\widehat{G}`$ is abelian of rank no greater than two. Since $`X`$ is Kähler if and only if $`\widehat{X}`$ is, this rank is even if $`X`$ is Kähler and odd if not.
In the case when $`\widehat{C}=^{}`$, the triviality of the conjugation action of $`\widehat{R}=`$ implies immediately that $`\widehat{G}`$ is abelian, of rank one greater than that of $`H`$.
In the case when $`\widehat{C}`$ is an elliptic curve, proposition 3.17 implies that $`\widehat{X}`$ must either be a holomorphically locally trivial fibration over $`\widehat{C}`$, or $`H`$ is finite cyclic. In the former case, proposition 3.22 tells us that $`\widehat{G}`$ is a finite extension of a free abelian group of rank four. In the latter case, let $`m/2`$ be the order of $`H`$. Since $`\widehat{R}`$ is abelian, the commutator of two elements $`a`$ and $`b`$ in $`\widehat{G}`$ must lie in $`H`$. Hence, $`ab=bac`$ for some $`cH`$. Since $`c`$ commutes with both $`a`$ and $`b`$, we have $`a^mb=ba^m`$ and $`a^mb^m=(ab)^m`$. This shows that $`\widehat{G}^m=\{a^m|a\widehat{G}\}`$ is an abelian subgroup of $`\widehat{G}`$ intersecting $`H`$ at $`1`$. Hence, we can form the internal direct sum $`G^mH`$ in $`G`$ which we can easily identify with the inverse image of $`\widehat{R}^m`$ in $`G`$, where $`\widehat{R}^m`$ is a normal subgroup of index $`m^2`$ in $`\widehat{R}`$. (We note as an aside that $`G^m`$ is canonically isomorphic to $`R^m`$.) It follows that $`\widehat{G}`$ becomes abelian if we replace $`\widehat{X}`$ by a finite covering of itself and so our theorem is proved.
## 4 Other compact complex surfaces
We deal with the remaining cases of compact complex surfaces in this section. These are the case of zero Kodaira dimension and the case of zero algebraic dimension. In fact, by Kodaira’s classification, all surfaces with Kodaira dimension zero are elliptic fibrations except for those bimeromorphic to compact complex 2-dimensional tori and K3 surfaces, where the elliptic ones form a dense codimension one family in their respective moduli space. As we have already resolved the case of elliptic fibrations in the previous section, we need to consider only the tori and the K3 surface cases. We first resolve the case of tori, and indeed prove a much stronger result of independent interest, before considering the other cases.
### 4.1 Compact complex tori
A 2-dimensional compact complex torus is the quotient of $`^2`$ by a lattice $`\mathrm{\Lambda }`$ of real rank 4. Let $`X`$ be such a surface, which we call a torus surface. Any compact surface $`Y`$ bimeromorphic to $`X`$ admits a dominating holomorphic map from the complement of finitely many points in $`X`$. We show in this section that the complement of finitely many points in $`X`$ is dominable by $`^2`$. This will follow immediately from proposition 4.1 below. Hence, $`Y`$ is also dominable by $`^2`$ as a result.
Following Rosay and Rudin \[RR1\], we say that a discrete set $`\mathrm{\Lambda }`$ in $`^2`$ is tame if there is a holomorphic automorphism, $`F`$, of $`^2`$ such that $`F(\mathrm{\Lambda })`$ is contained in a complex line. Using techniques of \[RR1\] or \[BF\], the complement of a tame set is dominable by $`^2`$, and in fact, there exists an injective holomorphic map from $`^2`$ to $`^2\mathrm{\Lambda }`$.
By a lattice, we mean a discrete $``$-module. For the following proposition, let $`\mathrm{\Lambda }`$ be a lattice in $`^2`$, let $`q_1,\mathrm{},q_m^2`$, and let $`\mathrm{\Lambda }_0=_{j=1}^m\mathrm{\Lambda }+q_j`$, where $`\mathrm{\Lambda }+q_j`$ represents translation by $`q_j`$.
###### PROPOSITION 4.1
The set $`\mathrm{\Lambda }_0`$ is tame. In particular, $`^2\mathrm{\Lambda }_0`$ is dominable by $`^2`$ using an injective holomorphic map.
This result will be strengthened considerably in section 6. Before proving this proposition, we need a lemma.
###### LEMMA 4.2
There exists an invertible, complex linear transformation $`A:^2^2`$ such that $`\mathrm{Im}\pi ^1A(\mathrm{\Lambda }_0)`$ is a discrete set in $``$. Moreover, we may assume that if $`p,qA(\mathrm{\Lambda }_0)`$ with $`pq`$, then $`|pq|1`$ and either $`\mathrm{Im}\pi ^1p=\mathrm{Im}\pi ^1q`$ or $`|\mathrm{Im}\pi ^1p\mathrm{Im}\pi ^1q|1`$.
Proof: Let $`v_1,v_2,v_3,v_4`$ be a $``$-basis for $`\mathrm{\Lambda }`$, and let $`E`$ be the span over $``$ of $`v_1,v_2,v_3`$. Using the real inner product, let $`u_00`$ be orthogonal to $`E`$. Using the complex inner product, let $`u_10`$ be orthogonal to $`u_0`$. Then $`u_1`$ and $`iu_1`$ are both real orthogonal to $`u_0`$, so $`u_1E`$. Choose $`A_1`$ complex linear such that $`A_1(u_0)=(1,0)`$ and $`A_1(u_1)=(0,1)`$. Then $`\pi ^1A_1(E)`$ is a one (real) dimensional subspace of $``$, so by rotating in the first coordinate, we may assume that $`\pi ^1A_1(E)`$ is the real line in $``$.
Let $`\mu _0=\mathrm{Im}\pi ^1A_1(v_4)`$, and $`\mu _j=\mathrm{Im}\pi ^1A_1(p_j)`$ for $`j=1,\mathrm{},m`$. Then for each $`j=1,\mathrm{},m`$ and $`k`$, we have
$$\pi ^1A_1(E+kv_4+p_j)+i(k\mu _0+\mu _j),$$
so that $`\mathrm{Im}\pi ^1A_1(\mathrm{\Lambda }_0)`$ is discrete in $``$. Applying an appropriate dilation to $`A_1`$ gives $`A`$ as desired.
Note that this lemma implies that given a finite set of points in a complex 2-torus, there is an open set, $`U`$, containing this finite set and a nonconstant image of $``$ avoiding $`U`$. In particular, the complement of $`U`$ in this torus is not Kobayashi hyperbolic. As mentioned in the introduction, this result will be strengthened in section 6 to show that there is a dominating map into the complement of such an open set $`U`$.
Proof of proposition 4.1: Lemma 4.2 implies that there is a complex line $`L=(z_0,w_0)`$ with orthogonal projection $`\pi _L:^2L`$ and real numbers $`\mu _0,\mathrm{},\mu _m`$ such that
$$\pi _L(\mathrm{\Lambda }_0)_{j=1}^m(\mu _0+\mu _j+i)(z_0,w_0).$$
(4.1)
I.e., identifying $`L`$ with $``$ in the natural way, the image of $`\mathrm{\Lambda }_0`$ under $`\pi _L`$ is contained in a union of lines parallel to the imaginary axis, and this union of lines intersects the real axis in a discrete set.
Making a linear change of coordinates, we may assume that $`L=(0,1)`$, in which case we may identify $`\pi _L`$ with projection to the second coordinate, $`\pi ^2`$. Let $`\pi ^1`$ denote projection to the first coordinate, and let $`E=_{j=1}^m(\mu _0+\mu _j+i)`$.
We next show that there is a continuous, positive function $`f_0`$ on $`E`$ such that if $`(z,w)\mathrm{\Lambda }_0`$ with $`z0`$, then $`f_0(w)|z|2|w|`$. First, define
$$r_1(w)=\{\begin{array}{cc}\frac{|w|}{\mathrm{min}\{|z|:(z,w)\mathrm{\Lambda }_0,z0\}}\hfill & \mathrm{if}w\pi ^2(\mathrm{\Lambda }_0);\hfill \\ 0\hfill & \mathrm{if}wE\pi ^2(\mathrm{\Lambda }_0).\hfill \end{array}$$
Then $`r_1(w)0`$, and since $`\mathrm{\Lambda }_0`$ is discrete, $`r`$ is upper-semicontinuous.
Let $`r_2(w)=2(r_1(w)+1)`$ for $`wE`$. Since $`r_2`$ is also upper-semicontinuous, it is bounded above on compacta, so a standard construction gives a function $`f_0`$ which is continuous on $`E`$ with $`f_0(w)r_2(w)>0`$. Then for $`(z,w)\mathrm{\Lambda }_0`$ with $`z0`$, we have $`f_0(w)|z|2r_1(w)|z|2|w|`$ by definition of $`r_1`$.
We next find a non-vanishing entire function $`f`$ so that $`|f(w)z||w|`$ if $`(z,w)\mathrm{\Lambda }_0`$ with $`z0`$. Since $`f_0`$ is positive on $`E`$, $`\mathrm{log}f_0(w)`$ is continuous and real-valued on $`E`$, and $`\mathrm{log}f_0(w)\mathrm{log}2+\mathrm{log}(r_1(w)+1)`$. By Arakelian’s theorem (e.g. \[RR2\]), there exists an entire $`g(w)`$ with $`|\mathrm{log}f_0(w)g(w)|<\mathrm{log}2`$ for $`wE`$. Then $`f(w)=\mathrm{exp}(g(w))`$ is entire and non-vanishing, and if $`(z,w)\mathrm{\Lambda }_0`$ with $`z0`$, then
$$|f(w)z|=\mathrm{exp}(\mathrm{Re}g(w))|z|r_1(w)|z||w|.$$
Finally, define $`F(z,w)=(f(w)z,w)`$. Then $`F`$ is a biholomorphic map of $`^2`$ onto itself, and for $`(z,w)\mathrm{\Lambda }_0`$ with $`z0`$, we have $`|\pi ^1F(z,w)||\pi ^2F(z,w)|`$. Since $`F(\mathrm{\Lambda }_0)`$ is discrete, we see that $`\pi ^1F(\mathrm{\Lambda }_0)`$ is discrete. Hence $`F(\mathrm{\Lambda }_0)`$ is tame by \[RR1, theorem 3.9\]. By definition of tame, $`\mathrm{\Lambda }_0`$ is also tame, so as mentioned earlier, $`^2\mathrm{\Lambda }_0`$ is dominable by $`^2`$.
###### COROLLARY 4.3
The complement of a finite set of points in a two dimensional compact complex torus is dominable by $`^2`$. Hence any surface bimeromorphic to such a torus is dominable by $`^2`$..
We remark that not all tori are elliptic. The elliptic torus surfaces form a $`3`$ dimensional family in the $`4`$ dimensional family of torus surfaces and the generic torus contains no curves. All compact complex tori are Kähler. Also a compact surface bimeromorphic to a torus can be characterized by $`\kappa =0`$ and $`q=2`$.
### 4.2 K3 surfaces
A compact complex surface $`X`$ is called a K3 surface if its fundamental group and canonical bundle are trivial. A useful fact in the compact complex category, due to Siu (\[Siu\]), is that all K3 surfaces are Kähler. One can show that $`H^2(X,)`$ is isometric to a fixed lattice $`L`$ of rank $`22`$. If $`\varphi `$ is such an isometry, then $`(X,\varphi )`$ is called a marked K3 surface. The set of such surfaces is parametrized by a $`20`$ dimensional non-Hausdorff manifold $``$ \[BPV, VIII\] (The fact that $``$ is smooth follows from S.T. Yau’s resolution of the Calabi conjecture in \[Yau\] (see e.g., \[T\]) and the fact that $``$ is not Hausdorff is due to Atiyah (\[At\]).)
We first observe a few facts from the classical work of Piatetsky-Shapiro and Shafarevich in \[PS\] (see also \[LP\],\[Shi\],\[BPV, VIII\]), where they obtained a global version of the Torelli theorem for K3 surfaces. Given a marked K3 surface and a point $`o`$ corresponding to it, there is a smooth Hausdorff neighborhood $`𝒰`$ of $`o`$, a smooth complex manifold $`Z`$, and a proper holomorphic map $`Z\stackrel{p}{}𝒰`$ whose fibers are exactly the marked K3 surfaces parametrized by $`𝒰`$. Within this local family, the subset of projective K3 surfaces is parametrized by a topologically dense subset of $`𝒰`$ which is a countable union of codimension one subvarieties. The elliptic K3 surfaces (that is, K3 surfaces admitting an elliptic fibration) also form a topologically dense codimension one family in $`𝒰`$.
The following proposition follows directly from theorem 3.23 and the fact that the fundamental group of a K3 surface is trivial.
###### PROPOSITION 4.4
A compact complex surface bimeromorphic to an elliptic K3 surface is holomorphically dominable by $`^2`$.
The previous section on complex tori allows us to deal with another class of K3 surfaces — the Kummer surfaces, which form a 4 dimensional family in the 20 dimensional family of K3 surfaces. Such a surface $`X`$ is, by definition, obtained by taking the quotient of a torus surface $`A`$ (given as a complex Lie group $`^2`$/lattice) by the natural involution $`g(x)=x`$, then blowing up the 16 orbifold singular points (resulting in 16 ($`2`$) curves). Alternatively, one can describe $`X`$ as a $`_2`$ quotient of $`\widehat{A}`$, where $`\widehat{A}`$ is the blowing up of $`A`$ at the 16 points of order $`2`$ and where the quotient map is branched along the exceptional (-1)-curves of the blowing up. Since the inverse image of any finite set of points in $`X`$ is finite in $`\widehat{A}`$ and hence also finite in $`A`$, any surface bimeromorphic to a Kummer surface is dominable by $`^2`$ according to corollary 4.3.
###### PROPOSITION 4.5
A compact surface bimeromorphic to a Kummer surface is dominable by $`^2`$.
Before we leave the subject of K3 surfaces, it is worth mentioning that projective K3 surfaces are dominable by $`𝔻\times `$ by the work of \[GG\] and \[MM\]. Clearly, elliptic K3 surfaces and Kummer surfaces are so dominable as well. Such a surface cannot be measure hyperbolic as defined by Kobayashi (\[Kob\]). However, it is still an unsolved problem whether all K3 surfaces are so dominable. The only other compact complex surfaces for which this problem remains open are the non-elliptic and non-Hopf surfaces of class VII<sub>0</sub> outside the Inoue-Hirzebruch construction.
### 4.3 Other compact surfaces and our two main theorems
Besides those bimeromorphic to K3 and torus surfaces, the remaining compact complex surfaces with zero Kodaira dimension are all elliptic, and are all dominable by $`^2`$. Such a surface must be bimeromorphic to either a Kodaira surface (defined and characterized in section 3.3.3), a hyperelliptic surface (which is a finite free quotient of a product of elliptic curves, and hence projective), or an Enriques surface (which is a surface admitting an unramified double covering by an elliptic K3 surface). Except for the first among these three types, the fundamental group is always a finite extension of an abelian group.
Finally, the only remaining compact complex surfaces are those with algebraic dimension 0 and $`\kappa =\mathrm{}`$. This category includes the non-elliptic Hopf surfaces, which are dominable by $`^2`$ by construction (see \[Ko4\]). This category also includes the Inoue surfaces, which must be excluded from our main theorems since their universal cover is $`𝔻\times `$, hence are not dominable by $`^2`$, while any nonconstant image of $``$ must be Zariski dense (see proposition 19.1 in \[BPV, V\]). However, it is of interest to note that the Zariski dense holomorphic images of $``$ are constrained by higher order equations on an Inoue surface so that if we relax property C in this sense, we can in fact include Inoue surfaces in the next theorem. Unfortunately, aside from the Hopf surfaces and the Inoue surfaces, the detailed structure of surfaces of this type is not yet clear even though we know the existence of projective affine structures for a special subclass of these surfaces.
We now summarize our investigation in the compact category by giving the following extensions of our main theorems stated in the introduction:
###### THEOREM 4.6
Let $`X`$ be a compact complex surface of Kodaira dimension less than $`2`$. Assume that either $`\kappa (X)\mathrm{}`$ or $`a(X)0`$. In the case that $`X`$ is bimeromorphic to a K3 surface that is not Kummer, assume further that $`X`$ is elliptic. Then $`X`$ is dominable by $`^2`$ if and only if it does not satisfy property C. Equivalently, there is a dominating holomorphic map $`F:^2X`$ if and only if there is a holomorphic image of $``$ in $`X`$ which is Zariski dense.
###### THEOREM 4.7
Let $`X`$ be a compact complex surface not bimeromorphic to a Kodaira surface. Assume that either $`\kappa (X)\mathrm{}`$ or $`a(X)0`$. In the case that $`X`$ is bimeromorphic to a K3 surface that is not Kummer, assume further that $`X`$ is elliptic. Then $`X`$ is dominable by $`^2`$ if and only if it has Kodaira dimension less than two and its fundamental group is a finite extension of an abelian group (of rank $`4`$ or less).
## 5 Non-compact algebraic surfaces
We begin with a key example which motivated the general algebraic setting. This is the example of the complement of a smooth cubic curve in $`^2`$, which we will show to be dominable by $`^2`$.
### 5.1 Complement of a cubic in $`^2`$
Let $`C`$ be a smooth cubic curve in $`^2`$ and let $`X=^2C`$. Then its logarithmic canonical bundle $`K_^2(D)`$ is the trivial line bundle as $`\mathrm{deg}K_^2=3`$. Hence, $`\overline{\kappa }(X)=0`$ and $`X`$ is a logarithmic K3 surface; that is, a non-compact 2-dimensional Calabi-Yau manifold.
###### PROPOSITION 5.1
The surface $`X=^2C`$ is dominable by $`^2`$.
Proof: A tangent line to $`C`$ at a non-inflection point meets $`C`$ at one other point. This gives rise to a holomorphic $`^1`$ bundle with two holomorphic sections. To see that this is actually a bundle (i.e. locally trivial), identify it with the projectivization of the tautological vector bundle of rank two over the dual curve of $`C`$ with the obvious isomorphism. We may pull back this $`^1`$ bundle and the sections to the universal cover $``$ of $`C`$, with two sections $`s_{\mathrm{}}`$ and $`s`$. Hence one may regard the complement of $`s_{\mathrm{}}()`$ of this bundle as a trivial line bundle on $``$ with a meromorphic section $`s`$ (with poles coming from points of inflection of the cubic).
Hence, it suffices to construct a holomorphic map from $`^2`$ onto the complement of the graph of a meromorphic function $`s`$ to give a dominating map to $`X`$. Note that each vertical slice of the complement of the graph is $`^{}`$ except at a pole of $`s`$, where the vertical slice is $``$.
To construct such a map, first define
$`\psi (t,w)`$ $`=`$ $`{\displaystyle \frac{\mathrm{exp}(tw)1}{t}}`$ (5.1)
$`=`$ $`w+{\displaystyle \frac{tw^2}{2!}}+{\displaystyle \frac{t^2w^3}{3!}}+\mathrm{}`$ (5.2)
which is entire on $`^2`$. Note that $`(t,\psi (t,w))`$ is a fiberwise selfmap of $`^2`$ which misses precisely the graph of $`1/t`$, a function with a simple pole at the origin.
Since $``$ is Stein, there exists an entire function $`g`$ such that $`\frac{1}{g}`$ has the same principle parts as $`s`$. This is because we may write $`s=f/f_1`$ where $`f`$ and $`f_1`$ are entire with no common zeros. So $`\mathrm{log}f`$ is well defined in a neighborhood of each zero of $`f_1`$. By Mittag-Leffler and Weierstrass, we can find an entire function $`g_1`$ with the same Taylor expansion as $`\mathrm{log}f`$ to the order of vanishing of $`f_1`$ at each zero of $`f_1`$. Then $`g=f_1/\mathrm{exp}g_1`$ is our desired function. In particular, $`g`$ vanishes precisely when $`s`$ has a pole. Then $`h=s\frac{1}{g}`$ is entire, so
$`\varphi (z,w)`$ $`=`$ $`h(z)\psi (g(z),w)`$
$`=`$ $`s(z){\displaystyle \frac{\mathrm{exp}(wg(z))}{g(z)}}`$
is entire on $`^2`$. For fixed $`z`$ with $`g(z)0`$, we see from the second equality that $`\varphi (z,w)`$ can attain any value in $`\{s(z)\}`$ by varying $`w`$. If $`g(z)=0`$, then $`\varphi (z,w)=h(z)w`$, which can attain any value in $``$ by varying $`w`$.
Hence, the map $`\mathrm{\Phi }:^2^2\text{graph}(s)`$ given by
$$\mathrm{\Phi }(z,w)=(z,\varphi (z,w))$$
is holomorphic and onto. Composing this map with the map into the $`^1`$ bundle over $`C`$, we obtain a dominating map into the complement of the cubic $`C`$.
Note that an important step here is the construction of an entire function $`h`$ whose graph does not intersect the graph of $`s`$. This is certainly analogous to the situation of elliptic fibrations.
Remark: The complement of a smooth cubic does not admit any algebraic map to $`^1`$ whose generic fiber contains $`^{}`$. This is the only example among complements of normal crossing divisors in $`^2`$ with this property. In fact, this is the only meaningful affine example with this property that is dominable by $`^2`$ (see \[M, p. 189\]). Since this is a logarithmic K3 surface, this phenomenon is suggestive of the situation for a generic compact K3 surface.
We isolate the following useful theorem from the above proof.
###### THEOREM 5.2
Let $`s`$ be a meromorphic function on $``$. Then the complement of the graph of $`s`$ admits a dominating fiber-preserving holomorphic map from $`^2`$.
#### 5.1.1 Complements of normal crossing divisors in $`^2`$
Let $`X`$ be the complement of a normal crossing divisor $`D`$ in $`^2`$. If $`\mathrm{deg}D>3`$, then $`\overline{\kappa }(X)=2`$ and hence $`X`$ is not dominable by $`^2`$. If $`\mathrm{deg}D=3`$, then $`D`$ consists of at most three components and it is easy to check that $`X`$ is dominable by $`^2`$ as follows. If $`D`$ has only one component, then it is either a smooth cubic or a cubic with one node. In the first case, the result follows from proposition 5.1. In the second case, blowing up that node gives us a $`^1`$ bundle over $`^1`$ with two sections, one corresponding to the exceptional curve of the blow-up. These two sections intersect precisely at the two fibers of the bundle corresponding to the two tangent directions of the cubic at the node. Hence, removing these two fibers gives us a surface biholomorphic to $`^{}\times ^{}`$, which is dominable by $`^2`$. If $`D`$ has two components, then it consists of a line and a conic (that is, a smooth curve of degree two) intersecting at two points. Blowing up one of the points of their intersection (corresponding to projecting from this point of intersection) gives us a $`^1`$ bundle over $``$ with two sections complemented, one of which is the exceptional curve of the blow-up. If we think of one section as $`\mathrm{}`$, then the other section can be regarded as a meromorphic function on $``$ and so theorem 5.2 applies to give a dominating map from $`^2`$ to $`X`$. An easier way is to delete the fiber containing the only point of intersection of these two sections. The resulting $`X`$ is biholomorphic to $`^{}\times ^{}`$ and hence dominable by $`^2`$. If $`D`$ has three components, then each must be a line and $`X`$ is $`^{}\times ^{}`$, which is dominable by $`^2`$.
From the above argument, we see also that if $`\mathrm{deg}D<3`$, then $`X`$ is dominable by $`^2`$. In summary, we have:
###### THEOREM 5.3
Let $`D`$ be a normal crossing divisor in $`^2`$. Then $`^2D`$ is dominable by $`^2`$ if and only if $`\mathrm{deg}D3`$.
We remark that this theorem is no longer true if $`D`$ is not normal crossing. The unique counterexample in one direction is when $`D`$ consists of three lines intersecting at only one point, which is not dominable by $`^2`$. Another counterexample, but in the opposite direction, is given by the complement of the union of a conic and two lines intersecting at a point of the conic (which we discussed in the two component case of $`\mathrm{deg}D=3`$ above).
### 5.2 The general quasi-projective case
Let $`X`$ be an algebraic surface over $``$. Then $`X=\overline{X}D`$ where $`\overline{X}`$ is projective and $`D`$ is a normal crossing divisor in $`\overline{X}`$. This is the notation set forth in section 2 and we will assume this setup throughout this section. Kawamata (\[K1\],\[K2\],\[K3\]) has considered the structure of $`X`$ and obtained a classification theory analogous to that in the projective case. Much of this is explained in some detail in Miyanishi (\[M\]). We will use their results directly to tackle our problem in this section.
If there is a surjective morphism $`f:XC`$ whose generic fiber is connected, then we say that $`X`$ is fibered over $`C`$. (We remind the reader that morphisms are algebraic holomorphic maps.) More generally, if $`f`$ is required to be only holomorphic rather than a morphism, then we say that $`X`$ is holomorphically fibered over $`C`$. For example, the complement of the graph of a meromorphic function is holomorphically fibered over $`C`$ with generic fiber $`^{}`$. As before, we let $`X_s=f^1(s)`$ be the fiber over $`s`$. We first quote the subadditivity property of (log-)Kodaira dimension due to Kawamata (\[K1\]):
###### PROPOSITION 5.4
If $`X`$ is fibered over a curve $`C`$, then
$$\overline{\kappa }(X)\overline{\kappa }(C)+\overline{\kappa }(X_s)$$
for $`s`$ outside a finite set of points in $`C`$; that is, for the generic fiber $`X_s`$.
From the definitions, a curve of positive genus with punctures has positive Kodaira dimension. An elliptic curve has Kodaira dimension zero. A punctured $`^1`$ has $`\kappa =\mathrm{},0`$ or $`1`$ according to the number of punctures being $`1,2`$ or greater than $`2`$, respectively.
Given a dominating morphism $`f`$ between algebraic varieties, it is clear that $`f^{}`$ is injective on the level of logarithmic forms (see \[Ii\]). Since tensor powers of top dimensional logarithmic forms define the Kodaira dimension, we see that if $`f`$ is equidimensional, then it must decrease Kodaira dimension.
If $`\overline{q}(X)>0`$, then there is a morphism from $`X`$ to a semi-abelian variety (a commutative algebraic Lie group that is an extension of a compact torus by $`(^{})^k`$ for some $`k`$) of dimension $`\overline{q}(X)`$, called the quasi-Albanese map and constructed by Iitaka in \[Ii1\]. One has the simple formula relating $`\overline{q}(X)`$ to the first Betti numbers of $`X`$ and $`\overline{X}`$:
$$\overline{q}(X)q(\overline{X})=b_1(X)b_1(\overline{X}).$$
Note that $``$ does not support any logarithmic form by this formula.
#### 5.2.1 Surfaces fibered by open subsets of $`^1`$
Let $`X`$ be fibered over a curve C by a map $`f`$ whose generic fiber is $`^1`$ (possibly) with punctures. Then, by a finite number of contractions of $`(1)`$-curves that remain on the fiber, the compactification $`\overline{X}`$ of $`X`$ admits a birational morphism $`g`$ to a ruled surface $`\overline{Y}`$ over a projective curve $`\overline{C}`$, the compactification of $`C`$, and $`g`$ is a composition of blowing ups. Hence $`Y=\overline{Y}|_C`$ is a $`^1`$ bundle over $`C`$, whose bundle map will again be denoted by $`g`$. We may write $`f=hg`$, where
$$\stackrel{r}{}X\stackrel{g}{}Y\stackrel{h}{}C.$$
(5.3)
If every holomorphic image of $``$ in $`X`$ is constant in $`C`$ (when composed with $`f`$), then $`X`$ satisfies property C. Otherwise, there exists a holomorphic map $`r:X`$ such that $`fr`$ is not constant. By taking the fiber product with $`fr`$, we can pull back the factorization picture 5.3 to one over $``$
$$\stackrel{\stackrel{~}{r}}{}\stackrel{~}{X}\stackrel{\stackrel{~}{g}}{}\stackrel{~}{Y}\stackrel{\stackrel{~}{h}}{},$$
where $`\stackrel{~}{f}=\stackrel{~}{h}\stackrel{~}{g}`$ is surjective with a holomorphic section $`\stackrel{~}{r}`$. Here, $`\stackrel{~}{X}`$ may be singular, but we will regard it only as an auxiliary space.
We will first deal with the case where the general fiber has at most one puncture; that is, $`X_s=^1`$ or $``$ for $`s`$ in an open subset of $`C`$. We can then regard $`\stackrel{~}{Y}`$ as a trivial $`^1`$ bundle with a section $`D_{\mathrm{}}`$ to which the puncture (if one exists) on the “generic” fiber of $`\stackrel{~}{f}`$ is mapped. Note that $`\stackrel{~}{Y}D_{\mathrm{}}=^2`$ with coordinates $`(z,w)`$, and so we may regard a section of $`\stackrel{~}{h}`$ as a meromorphic function on $``$. In particular, $`\stackrel{~}{g}\stackrel{~}{r}`$ is a meromorphic section of $`\stackrel{~}{h}`$.
Since $`\overline{X}`$ is obtained from $`\overline{Y}`$ by a finite number of blow ups, we can identify points on $`X`$ as infinitely near points on $`Y`$ of order 0 or more as in \[Ha, p. 392\]. Note that the set of fibers in $`Y`$ which contain infinitely near points of order 1 or more is finite (since the set of such fibers in $`\overline{Y}`$ is finite). This finite set of fibers in $`Y`$ pulls back to a discrete set of fibers in $`\stackrel{~}{Y}`$. In $`Y`$, such a higher order infinitely near point corresponds to a point in $`X`$ obtained by finitely many blow-ups, hence to the specification of a finite jet at the point in $`Y`$. Under pull-back, this corresponds to a finite jet in $`\stackrel{~}{Y}`$. Additionally, there is a finite set of fibers in $`Y`$ which may have more than one puncture, and these fibers all pull back to a discrete set of fibers in $`\stackrel{~}{Y}`$. Together, these two types of fibers will be called exceptional fibers.
In order to produce a dominating map into $`X`$, it suffices to produce a fiberwise dominating map $`F(z,w)=(z,H(z,w))`$ into $`\stackrel{~}{Y}`$ which respects these exceptional fibers in the following sense. If $`\stackrel{~}{Y}_s`$ is an exceptional fiber, then $`F(s,w)`$ is a single point independent of $`w`$. Moreover, if $`\stackrel{~}{Y}_s`$ is a fiber having more than one puncture, then the image of the map $`F`$ should avoid all such punctures. If $`\stackrel{~}{Y}_s`$ is a fiber having an infinitely near point, then $`F(s,w)`$ should equal $`\stackrel{~}{g}\stackrel{~}{r}(s)`$. Additionally, if $`\stackrel{~}{g}\stackrel{~}{r}`$ passes through this infinitely near point, and $`\varphi `$ is holomorphic in a neighborhood of $`s`$, then the local curve $`z(z,H(z,\varphi (z)))`$ should agree with the jet given by the infinitely near point on $`\stackrel{~}{Y}_s`$.
Fortunately, the section $`\stackrel{~}{g}\stackrel{~}{r}`$ has the correct jet whenever it intersects one of these exceptional fibers, so we can use this section to obtain such a map. Let $`q(z)=\stackrel{~}{g}\stackrel{~}{r}(z)`$, which is meromorphic. We will define $`H(z,w)=p(z)w+q(z)`$ for some entire $`p(z)`$. For each exceptional fiber $`\stackrel{~}{Y}_s`$, there is an integer $`n_s1`$ such that if $`p`$ vanishes to order $`n_s`$ at $`s`$, then $`F`$ defined with this $`H`$ respects the exceptional fiber as indicated above. By Weierstrass’ theorem, there exists $`p`$ entire vanishing exactly to order $`n_s`$ at each $`s`$. Then $`F(z,w)=(z,H(z,w))`$ gives a dominating map from $`^2`$ into $`\stackrel{~}{Y}`$ respecting the exceptional fibers, and this map pushes forward to $`Y`$, then lifts to give a dominating map into $`X`$, as desired.
We now deal with the case where the generic fiber of $`f`$ is $`^{}`$. In this case, $`\stackrel{~}{Y}`$ can be identified with a $`^1`$ bundle with a double section $`D_Y`$, to which the punctures on the “generic” fibers of $`\stackrel{~}{f}`$ maps to. Now, either $`D_Y`$ consists of two components, both of which are smooth sections of $`\stackrel{~}{h}`$, or $`D_Y`$ consists of one component. In either case, outside of a discrete set of fibers, $`D_Y`$ can be written locally as the union of two mermorphic sections. Moreover, we define the set of exceptional fibers exactly as in the previous case.
First, using a fiber-preserving biholomorphic map of $`\times ^1`$ to itself, we may move $`\stackrel{~}{g}\stackrel{~}{r}`$ to become the $`\mathrm{}`$-section. Then the requirement of agreeing with the jet of $`\stackrel{~}{g}\stackrel{~}{r}`$ at a point $`s`$ is equivalent to having a pole of some given order at $`s`$ in the new coordinate system. Next, let $`E_1`$ be the points in $``$ at which $`D_Y`$ intersects this new infinity section. Near a point $`sE_1`$, $`D_Y`$ can be written as $`w=h(z)\pm \sqrt{g(z)}=u^\pm (z)`$ for some meromorphic $`g`$ and $`h`$. Hence there exists $`n_s>0`$ such that $`u^\pm (z)(zs)^{n_s}`$ converges to 0 as $`z`$ tends to $`s`$. We may assume also that if $`sE_1`$ and $`s`$ is the base point of an exceptional fiber, then the $`n_s`$ obtained here is larger than the $`n_s`$ obtained above for this exceptional fiber.
Let $`E`$ be the union of $`E_1`$ and the set of base points corresponding to exceptional fibers. Let $`p`$ be entire with a zero of order $`n_s`$ at each $`sE`$ and no other zeros, and let $`\mathrm{\Phi }(z,w)=(z,p(z)w)`$. Then $`\mathrm{\Phi }(D_Y)`$ is a double section in $`\times ^1`$, and a dominating map from $`^2`$ to $`^2\mathrm{\Phi }(D_Y)`$ followed by $`\mathrm{\Phi }^1`$ gives a dominating map to the complement of $`D_Y`$ which respects the exceptional fibers.
Hence it suffices to construct a dominating map into the complement of $`\mathrm{\Phi }(D_Y)`$. Note that $`\mathrm{\Phi }(D_Y)`$ can be written as $`w=v^\pm (z)=p(z)u^\pm (z)`$, where $`v^\pm `$ are holomorphic except possibly for square root singularities at branch points.
For complex numbers $`u`$ and $`v`$, define a Mobius transformation $`N_{u,v}(w)=(uwv)/(w1)`$, which takes 0 to $`v`$ and $`\mathrm{}`$ to $`u`$, and define $`G_{u,v}(w)=\mathrm{exp}(w(uv))`$. Note that $`N_{u,v}(w)=N_{v,u}(1/w)`$ and that $`G_{u,v}(w)=1/G_{v,u}(w)`$. Hence $`H_0(u,v,w)=N_{u,v}(G_{u,v}(w))`$ satisfies $`H_0(u,v,w)=H_0(v,u,w)`$. Since symmetric functions of $`v^+`$ and $`v^{}`$ are holomorphic, we see that $`H(z,w)=H_0(v^+(z),v^{}(z),w)`$ is well-defined and holomorphic from $`^2`$ to $`\times ^1`$. Moreover, for fixed $`s`$ such that $`v^\pm (s)`$ are distinct, $`H(s,)`$ is nonconstant from $``$ to $`^1\{v^\pm (s)\}`$. If $`v^\pm (s)`$ are equal, then assuming without loss that $`s=0`$, we have $`v^\pm (z)=h(z)\pm \sqrt{g(z)}`$ for some holomorphic $`g(z)=z^mg_1(z)`$ with $`g_1(0)0`$, $`m1`$. Then $`v^+v^{}=2\sqrt{g}`$, so multiplying the numerator and denominator of $`H`$ by $`\mathrm{exp}(w(v^+v^{})/2)`$ and using the Taylor expansion of $`\mathrm{exp}`$ gives
$`H`$ $`={\displaystyle \frac{(h+\sqrt{g})(1+w\sqrt{g})(h\sqrt{g})(1w\sqrt{g})+O(|z|^m)}{(1+w\sqrt{g})(1\sqrt{g})+O(|z|^m)}}`$
$`={\displaystyle \frac{2hw\sqrt{g}+2\sqrt{g}+O(|z|^m)}{2w\sqrt{g}+O(|z|^m)}}.`$
As $`z0`$, this last expression tends to $`h(0)+1/w`$, and hence $`H(0,)`$ maps $``$ onto $`^1\{v\pm (z)\}`$.
Thus $`H`$ is a dominating map from $`^2`$ to the complement of $`\mathrm{\Phi }(D_Y)`$, hence as noted before, $`\mathrm{\Phi }^1H`$ is a dominating map from $`^2`$ to the complement of $`D_Y`$ which respects the exceptional fibers. As before, this map pushes forward to $`Y`$ and lifts to give a dominating map into $`X`$, as desired.
We can now summarize with the following theorem.
###### THEOREM 5.5
Assume that $`X`$ is fibered over a curve $`C`$ and that the generic fiber is $`^1`$ with at most two punctures. Then $`X`$ is dominable by $`^2`$ if and only if there is a Zariski dense image of $``$ in $`X`$.
The arguments given in this paper are not sufficient to resolve the question of dominability for open fibered surfaces. As an example, we have the following question.
###### QUESTION 5.6
Let $`X`$ be the complement of a double section in a conic bundle over $``$, $`^{}`$, or an elliptic curve. Is $`X`$ dominable by $`^2\mathrm{?}`$
We will consider this and related questions in a forthcoming paper.
#### 5.2.2 The $`\overline{\kappa }=\mathrm{}`$ case
Let $`\overline{\kappa }(X)=\mathrm{}`$. Then $`\kappa (\overline{X})=\mathrm{}`$ as well. Hence $`\overline{X}`$ is either rational or birationally ruled over a curve of non-negative genus. In the latter case, proposition 5.4 says that $`X`$ is fibered over a curve $`C`$ with $`\kappa (C)0`$ where the generic fiber is $`^1`$ with at most one puncture. Hence theorem 5.5 applies in this case to give us the equivalence of dominability by $`^2`$ and the failure of property C. Note that property C holds in the case $`\kappa (C)>0`$ (which include the case $`\overline{q}(X)2`$), corresponding to $`C`$ being hyperbolic.
In the remaining case when $`\overline{X}`$ is rational, we can again divide into two cases according to whether $`\overline{q}(X)`$ is zero or not. In the latter case, we again have a fibering of $`X`$ over a curve $`C`$ via the quasi-Albanese map with the generic fiber having at most one puncture by proposition 5.4, as before. This is because there are no logarithmic $`2`$-forms on $`X`$ since $`\overline{\kappa }(X)=\mathrm{}`$. By the same token, every logarithmic $`1`$-form on $`X`$ is the pull back of a logarithmic form on $`C`$ (One can also see this from the fact that $`^1`$ with at most one puncture has no logarithmic forms so that any logarithmic form on $`X`$ becomes trivial when restricted to the generic fiber. Hence, $`\overline{q}(X)=\overline{q}(C)`$.) So, $`C`$ must be $`^1`$ with at least two punctures. If it has more than two punctures, corresponding to $`\overline{q}(X)2`$, then $`C`$ is hyperbolic. So we have degeneracy of holomorphic maps from $``$ in this case. Otherwise, theorem 5.5 applies.
We are left with the case where $`\overline{q}(X)=0`$ where proposition 5.4 no longer applies, but where much of the analysis has been done in \[M\]. We now quote theorem $`(1^{})`$ of \[M\], (which follows from theorem I.3.11 of \[M\])
###### THEOREM 5.7
With $`X`$ and $`D`$ as before, assume that $`D`$ is connected. Then $`\overline{\kappa }(X)=\mathrm{}`$ if and only if $`X`$ fibers over a curve with generic fiber being $`^1`$ or $``$.
Except in the case where $`X=^2`$, there is, of course, some fibering of $`X`$ to a curve (as is clear from, for example, (1) of the classification list given in section 2) and every such fibering must be to a curve $`C`$ that is either $`^1`$ or $``$. In these fibered cases, we would like to show that the generic fiber is $`^1`$ with at most two punctures so that theorem 5.5 can be applied to show that $`X`$ is dominable by $`^2`$. However, it remains an open question whether or not the generic fiber has this form, and although this question should be resolved by some case checking, this lack prevents us from giving a complete classification in the case $`\overline{\kappa }(X)=\mathrm{}`$ and $`\overline{q}(X)=0`$.
We can now summarize this section as follows.
###### THEOREM 5.8
Let $`X`$ be the complement of a normal crossing divisor $`D`$ in a projective surface. Assume $`\overline{\kappa }(X)=\mathrm{}`$. If $`\overline{q}(X)2`$, then $`X`$ satisfies property C and hence is not dominable by $`^2`$. If $`\overline{q}(X)=1`$ or if $`\overline{q}(X)=0`$ and $`D`$ is connected, then $`X`$ is dominable by $`^2`$ if and only if there exists a holomorphic map of $``$ to $`X`$ whose image is Zariski dense.
#### 5.2.3 The $`\overline{\kappa }=1`$ case
Here, we can directly apply the basic Iitaka fibration theorem, theorem 11.8 in \[Ii\] (see also \[Ue\]):
###### THEOREM 5.9
Assume $`\overline{\kappa }(X)0`$. Then $`X`$ is properly birational to a variety $`X^{}`$ which is fibered over a variety of dimension $`\overline{\kappa }(X)`$ and whose generic fiber has Kodaira dimension zero.
This theorem holds for $`X`$ of any dimension. But for our situation at hand, it says that $`X`$ is properly birational to a surface $`X^{}`$ which is fibered over a curve with generic fiber that is either an elliptic curve, or $`^1`$ with two punctures. Now, we have already shown that for such a fibered variety, dominability is unchanged for any variety properly birational to it. The latter case is already resolved by theorem 5.5. The former case can also be resolved to give the same conclusion by the same analysis as that of theorem 5.5 with the help of the Jacobian fibration as in section 3. Thus, combining with theorem 5.5, we have:
###### THEOREM 5.10
Assume $`X`$ is fibered over a curve with generic fiber that is either an elliptic curve or $`^1`$ with at most two punctures. This is the case, for example, when $`\overline{\kappa }(X)=1`$. Then $`X`$ is dominable by $`^2`$ if and only if there exists a holomorphic map of $``$ to $`X`$ whose image is Zariski dense.
#### 5.2.4 The $`\overline{\kappa }=0`$ case
It remains to look at the case where $`\overline{\kappa }(X)=0`$. If $`\overline{q}(X)2`$, then a well known theorem of Kawamata (\[K4\]) says that $`X`$ has a birational morphism to a semi-abelian surface. Hence, $`X`$ is dominable by $`^2`$. If $`\overline{q}(X)=1`$, then $`X`$ is fibered over a curve and the generic fiber is an elliptic curve or is $`^1`$ with at most two punctures by proposition 5.4. Hence theorem 5.10 applies in this case. When $`\overline{q}(X)=0`$, our problem remains with some K3 surfaces as explained in section 4.2.
Finally, if $`X`$ is affine rational and $`D`$ has a component that is not a rational curve, then Lemma II.5.5 of \[M\] says that either $`X`$ is fibered over a curve with generic fiber $`^1`$ with at most two punctures or $`X`$ is the complement of a smooth cubic in $`^2`$. The former is handled by theorem 5.5 while the latter is dominated by $`^2`$ as shown in section 5.1. This resolves the case of the complement of a reduced divisor $`C`$ in $`^2`$ unless $`C`$ is a rational curve, which one can resolve as well when $`C`$ has either only one singular point or is of low degree (and it is easy to check all the cases for degree less than $`4`$). This is a good exercise for the case when $`C`$ is a rational curve of high degree, which we will not attempt here. Note that, if $`C`$ is normal crossing with dominable complement, then $`C`$ is again a smooth cubic in $`^2`$, being the unique non-rational component.
###### THEOREM 5.11
Assume $`\overline{\kappa }(X)=0`$. If $`\overline{q}(X)`$ is positive or if $`X`$ is affine and $`D`$ has a component that is not a rational curve, then $`X`$ is dominable by $`^2`$ if and only if there exists a holomorphic map of $``$ to $`X`$ whose image is Zariski dense.
## 6 Compact complex surface minus small balls
For the compact complex surfaces which we showed to be dominable by $`^2`$, a surprisingly stronger result can be achieved, thanks to the theory of Fatou-Bieberbach domains. We can show that these surfaces remain dominable after removing any finite number of sufficiently small open balls. In this section we show how this can be done in the most difficult case, the case of a two dimensional compact complex torus. We show that given any finite set of points in a torus $`T`$, it is possible to find some open set, $`U`$, containing this finite set, and a holomorphic map $`F:^2TU`$ with non-vanishing Jacobian determinant. In fact, $`F`$ lifts to an injective holomorphic map from $`^2`$ to $`^2`$. For the statement of the following theorem, we focus only on this lifted map. For notation, $`\mathrm{\Delta }^2(p;r)`$ is the bidisk with center $`p`$ and radii $`r`$ in both coordinate directions and $`\pi ^j`$ represents projection to the $`j`$th coordinate axis.
###### THEOREM 6.1
Let $`\mathrm{\Lambda }^2`$ be a discrete lattice, let $`p_1,\mathrm{},p_m^2`$, let $`\mathrm{\Lambda }_0=_{j=1}^m\mathrm{\Lambda }+p_j`$, and for $`r>0`$, let $`\mathrm{\Lambda }_{0,r}=_{p\mathrm{\Lambda }_0}\mathrm{\Delta }^2(p;r)`$. For some $`r>0`$, there exists an injective holomorphic map $`F:^2^2\mathrm{\Lambda }_{0,r}`$.
In fact, the proof will show that any discrete set contained in $`\mathrm{\Lambda }_{0,r}`$ is a tame set in the sense of section 4.1. As an immediate corollary, we obtain the following result, as mentioned in the introduction. An $`n`$-dimensional version of this result is found in \[Bu\].
###### COROLLARY 6.2
Let $`T`$ be a complex 2-torus and let $`ET`$ be finite. Then there exists an open set $`U`$ containing $`E`$ and a dominating map from $`^2`$ into the complement of $`U`$.
For the remainder of this section, $`\mathrm{\Lambda }`$, $`\mathrm{\Lambda }_0`$ and $`\mathrm{\Lambda }_{0,r}`$ will be as in the statement of this theorem.
### 6.1 Preparatory lemmas
In this subsection we state some necessary lemmas. The proofs are straightforward and perhaps even standard, but they are provided for completeness.
##### Notation:
For $`ϵ>0`$, let $`S_ϵ=\{x+iy:x,ϵ<y<ϵ\}`$.
###### LEMMA 6.3
Let $`C>0`$, let $`f:[0,C]`$ be measurable, and let $`ϵ(0,1)`$. Then there exists a function $`g`$ holomorphic on $`S_ϵ`$ such that if $`\delta >0`$ and $`z_0=x_0+iy_0S_ϵ`$ with $`f(x)=c_0`$ for $`x_0\delta <x<x_0+\delta `$, then
$$|g(z_0)f(x_0)|\frac{2Cϵ}{\pi \delta }.$$
Moreover, $`\mathrm{Re}g(z)0`$ for all $`zS_ϵ`$.
Proof: For $`n`$, let
$`g_n(z)`$ $`={\displaystyle \frac{1}{2\pi i}}{\displaystyle _n^n}\left({\displaystyle \frac{f(x)}{xiϵz}}{\displaystyle \frac{f(x)}{x+iϵz}}\right)𝑑x`$
$`={\displaystyle \frac{1}{\pi }}{\displaystyle _n^n}f(x)\left({\displaystyle \frac{ϵ}{(xz)^2+ϵ^2}}\right)𝑑x.`$
I.e., $`g_n`$ is obtained via the Cauchy integral using the function $`f`$ on the two boundary components of $`S_ϵ`$ and truncating at $`x=\pm n`$. By \[R, Thm 10.7\], each $`g_n`$ is holomorphic in $`S_ϵ`$. Moreover, for $`z_0=x_0+iy_0S_ϵ`$, we have $`|y_0|<ϵ`$, so
$$|(xz)^2+ϵ^2|\mathrm{Re}(x(x_0+iy_0))^2+ϵ^2(xx_0)^2.$$
(6.1)
Using this last inequality and the boundedness of $`f`$, it follows immediately that $`g_n`$ converges uniformly on compact subsets of $`S_ϵ`$ to the holomorphic function
$$g(z)=\frac{1}{\pi }_{\mathrm{}}^{\mathrm{}}f(x)\left(\frac{ϵ}{(xz)^2+ϵ^2}\right)𝑑x.$$
(6.2)
A simple contour integration shows that if $`f`$ is replaced by the constant $`c_0`$, then the integral in (6.2) is $`c_0`$ for all $`zS_ϵ`$. Hence, if $`z_0=x_0+iy_0S_ϵ`$ with $`f(x)=c_0`$ for $`x_0\delta xx_0+\delta `$, then using (6.1),
$`|g(z_0)f(x_0)|`$ $`=\left|{\displaystyle \frac{1}{\pi }}{\displaystyle _{\mathrm{}}^{\mathrm{}}}(f(x)c_0)\left({\displaystyle \frac{ϵ}{(xz)^2+ϵ^2}}\right)𝑑x\right|`$
$`{\displaystyle \frac{C}{\pi }}\left({\displaystyle _{\mathrm{}}^{x_0\delta }}+{\displaystyle _{x_0+\delta }^{\mathrm{}}}{\displaystyle \frac{ϵ}{(xx_0)^2}}𝑑x\right)`$
$`{\displaystyle \frac{2ϵC}{\pi \delta }}.`$
To show that $`\mathrm{Re}g(z)0`$, note that the second part of (6.1) implies that $`\mathrm{Re}(ϵ/((xz)^2+ϵ^2))0`$ for all $`zS_ϵ`$, and since $`f`$ is real, (6.2) implies $`\mathrm{Re}g(z)0`$.
###### LEMMA 6.4
Let $`V=\{(z,w):|w|<1+|z|^2\}`$. Then there exists an injective holomorphic map $`\mathrm{\Phi }:^2V`$.
Proof: Let $`H(z,w)=(w,w^2z/2)`$. Then $`H`$ is a polynomial automorphism of $`^2`$, and $`(0,0)`$ is an attracting fixed point for $`H`$. By \[RR1, appendix\], there is an injective holomorphic map $`\mathrm{\Psi }`$ from $`^2`$ onto the basin of attraction of $`(0,0)`$, which is defined as $`B=\{p^2:lim_n\mathrm{}H^n(p)=(0,0)\}`$. By \[BS\], there exists $`R>0`$ such that $`B`$ is contained in
$$V_R=\{|z|R,|w|<R\}\{|z|R,|w|<|z|\}.$$
Hence taking $`\mathrm{\Phi }=\mathrm{\Psi }/R`$ gives an injective holomorphic map from $`^2`$ into $`V_1V`$.
### 6.2 Proof of theorem 6.1
We will construct an automorphism of $`^2`$ mapping $`\mathrm{\Lambda }_{0,r}`$ into the complement of the set $`V`$ of lemma 6.4. This will be sufficient to prove the theorem, and by \[RR1\] this implies that any discrete set contained in $`\mathrm{\Lambda }_{0,r}`$ is tame.
Choose an invertible, complex linear $`A`$ as in lemma 4.2. Without loss of generality, we may replace $`\mathrm{\Lambda }`$ by $`A(\mathrm{\Lambda })`$, $`p_j`$ by $`A(p_j)`$, and $`\mathrm{\Lambda }_0`$ by $`A(\mathrm{\Lambda }_0)`$. Then $`\pi ^1\mathrm{\Lambda }_0`$ is contained in $`_{k=1}^{\mathrm{}}L_k`$, where each $`L_k`$ is a line of the form $`+i\gamma _k`$, $`\gamma _k`$ real. Moreover, $`\mathrm{dist}(L_j,L_k)1`$ if $`jk`$, and $`|pq|1`$ if $`p,q\mathrm{\Lambda }_0`$ with $`pq`$.
Let $`\{q_j\}_{j=1}^{\mathrm{}}`$ be an enumeration of the set
$$\{q\mathrm{\Lambda }_0:|\pi ^2q|1/8\}=\{q\mathrm{\Lambda }_0:\overline{\mathrm{\Delta }^2(q;1/8)}(\times \{0\})\mathrm{}\}.$$
Let $`C=\mathrm{log}32`$, and define $`f_k:[0,C]`$ for each $`k`$ by
$$f_k(x)=\{\begin{array}{cc}0\hfill & \mathrm{if}(x+i\gamma _k,0)\overline{\mathrm{\Delta }^2(q_j;1/8)}\mathrm{for}\mathrm{some}q_j\hfill \\ C\hfill & \mathrm{otherwise}.\hfill \end{array}$$
Let $`\delta =1/16`$, and choose $`ϵ\delta /2`$ small enough that $`2Cϵ/\pi \delta \mathrm{log}(3/2)`$. Let $`r=ϵ/2`$, and recall that $`\mathrm{\Lambda }_{0,r}=_{p\mathrm{\Lambda }_0}\mathrm{\Delta }^2(p;r)`$.
Let $`S_ϵ^k=\{x+i(y+\gamma _k):ϵ<y<ϵ\}`$ and $`U_ϵ=_{k=1}^{\mathrm{}}S_ϵ^k`$. Define $`g`$ holomorphic on $`U_ϵ`$ by applying lemma 6.3 with $`f=f_k`$ to define $`g`$ on $`S_ϵ^k`$. By Arakelian’s Theorem (e.g. \[RR2\]), there exists $`h`$ entire such that if $`zU_{ϵ/2}`$, then $`|h(z)g(z)|\mathrm{log}(4/3)`$. Define
$$F_1(z,w)=(z,w\mathrm{exp}(h(z))).$$
Then $`F_1:^2^2`$ is biholomorphic.
We show next that there is a complex line in the complement of $`F_1(\mathrm{\Lambda }_{0,r})`$. To do this, let $`p\mathrm{\Lambda }_{0,r}`$, and suppose first that $`p\mathrm{\Delta }^2(q_j;r)`$ for some $`q_j`$. Choose $`k`$ so that $`\gamma _k=\mathrm{Im}\pi ^1q_j`$, and write $`\pi ^1p=x_0+iy_0`$.
Note that $`|y_0\gamma _k|<r=ϵ/2`$. Also, since $`|\pi ^2q_j|1/8`$, we see that if $`|xx_0|<(1/8)r`$, then $`(x+i\gamma _k,0)\overline{\mathrm{\Delta }^2(q_j;1/8)}`$. Since $`\delta <(1/8)r`$, we have $`f_k(x)=0`$ for $`x_0\delta xx_0+\delta `$, and hence by lemma 6.3 and the choice of $`ϵ`$ and $`h`$,
$`|h(\pi ^1p)|`$ $`|g(\pi ^1p)|+\mathrm{log}(4/3)`$
$`{\displaystyle \frac{2Cϵ}{\pi \delta }}+\mathrm{log}(4/3)`$
$`\mathrm{log}2.`$
Hence
$$|\pi ^2F_1(p)|2|\pi ^2p|2(|\pi ^2q_j|+r)\frac{1}{3}.$$
(6.3)
In the remaining case, $`p\mathrm{\Lambda }_{0,r}`$ but $`p\mathrm{\Delta }^2(q_j;r)`$ for any $`j`$, in which case $`|\pi ^2p|(1/8)r`$. Let $`q\mathrm{\Lambda }_0`$ such that $`p\mathrm{\Delta }^2(q;r)`$, and choose $`k`$ so that $`\gamma _k=\mathrm{Im}\pi ^1q`$.
Suppose first that $`x_0=\mathrm{Re}\pi ^1p`$ satisfies $`f_k(x)=C`$ for $`|xx_0|\delta `$. Since $`|y_0\gamma _k|<r=ϵ/2`$, we have by lemma 6.3 and choice of $`ϵ`$ and $`h`$ that
$`\mathrm{Re}h(\pi ^1p)`$ $`\mathrm{Re}g(\pi ^1p)\mathrm{log}(4/3)`$
$`C{\displaystyle \frac{2Cϵ}{\pi \delta }}\mathrm{log}(4/3)`$
$`\mathrm{log}16.`$
Hence
$$|\pi ^2F_1(p)|16|\pi ^2p|16((1/8)r)>1.$$
(6.4)
Otherwise, $`f_k(x)=0`$ for some $`x`$ with $`|xx_0|\delta `$, so there exists $`j`$ such that $`|\pi ^1p\pi ^1q_j|(1/8)+\delta +r`$, hence
$$|\pi ^1q\pi ^1q_j|(1/8)+\delta +2r1/4.$$
Since $`q`$ and $`q_j`$ are distinct points of $`\mathrm{\Lambda }_0`$, we have $`|qq_j|1`$ by assumption, so $`|\pi ^2q\pi ^2q_j|^21(1/4)^2`$, and hence
$$|\pi ^2q||\pi ^2q\pi ^2q_j||\pi ^2q_j|\frac{\sqrt{15}}{4}\frac{1}{8}$$
and
$$|\pi ^2p||\pi ^2q|r\frac{3}{4}.$$
Since $`\mathrm{Re}g(\pi ^1p)0`$ by lemma 6.3, we have $`\mathrm{Re}h(\pi ^1p)\mathrm{log}(4/3)`$, and hence
$$|\pi ^2F_1(p)|\frac{3}{4}|\pi ^2p|\frac{9}{16}.$$
(6.5)
From (6.3), (6.4) and (6.5), we conclude that if $`p\mathrm{\Lambda }_{0,r}`$, then either $`|\pi ^2F_1(p)|1/3`$ or $`|\pi ^2F_1(p)|9/16`$. In particular,
$$\mathrm{dist}(F_1(\mathrm{\Lambda }_{0,r}),\times \{\frac{1}{2}\})\frac{1}{16}.$$
Note also that $`\pi ^1F_1(p)=\pi ^1p`$ for all $`p^2`$.
To finish the proof, we will construct $`F_2`$ similar to $`F_1`$ so that $`F_2(F_1(\mathrm{\Lambda }_{0,r}))`$ is contained in $`^2V`$, where $`V`$ is as in lemma 6.4.
First note that for $`z=x+iyS_ϵ^k`$, we have $`|yi\gamma _k|<ϵ`$, so
$$\mathrm{Re}[(zi\gamma _k)^2+(|\gamma _k|+r)^2+1+ϵ^2]x^2+(|\gamma _k|+ϵ)^2+1>0.$$
(6.6)
Hence we can choose a branch of $`\mathrm{log}`$ so that
$$g_2(z)=\mathrm{log}((zi\gamma _k)^2+(|\gamma _k|+r)^2+1+ϵ^2)+1+\mathrm{log}16$$
(6.7)
is holomorphic on $`_kS_ϵ^k`$. Again by Arakelian’s Theorem, there exists $`h_2`$ entire such that if $`zS_{ϵ/2}^k`$, then $`|g_2(z)h_2(z)|1`$, so $`\mathrm{Re}h_2(z)\mathrm{Re}g_2(z)1`$. Let
$$F_2(z,w)=(z,\left(w\frac{1}{2}\right)\mathrm{exp}(h_2(z))).$$
Again, $`F_2:^2^2`$ is biholomorphic. Moreover, if $`pF_1(\mathrm{\Lambda }_{0,r})`$, then $`|\pi ^2p\frac{1}{2}|1/16`$, and $`\pi ^1p=z=x+iy`$ with $`|y\gamma _k|<r`$ for some $`k`$, so by (6.6) and (6.7), we have
$`|\pi ^2F_2(p)|`$ $`\left|\pi ^2p{\displaystyle \frac{1}{2}}\right|\mathrm{exp}(\mathrm{Re}h_2(z))`$
$`{\displaystyle \frac{1}{16}}\mathrm{exp}(\mathrm{Re}g_2(z)1)`$
$`x^2+(|\gamma _k|+r)^2+1`$
$`1+|\pi ^1p|^2`$
$`1+|\pi ^1F_2(p)|^2.`$
Hence $`F_2F_1(\mathrm{\Lambda }_{0,r})V=\mathrm{}`$, where $`V\mathrm{\Phi }(^2)`$ is as in lemma 6.4, so taking $`F=F_1^1F_2^1\mathrm{\Phi }`$ gives an injective holomorphic map $`F:^2F_1^1F_2^1(V)^2\mathrm{\Lambda }_{0,r}`$ as desired.
### 6.3 The general case of complements of small open balls
It is now easy to deduce the following corollary from theorem 6.1.
###### COROLLARY 6.5
Let $`X`$ be bimeromorphic to a compact complex torus or to a Kummer K3 surface. Then, given any finite set of points in $`X`$, the complement of a neighborhood of this set is dominable by $`^2`$. In particular, such a complement is not measure hyperbolic.
The case of elliptic fibrations over $`^1`$ or over an elliptic curve can be handled in the same way as that of theorem 6.1. This is because removing a finite number of small open balls (plus a smooth fiber away from them if the base is $`^1`$) is tantamount to removing via the Jacobian fibration a discrete set of contractible open sets in $`^2`$ bounded away from the axis by fixed constants and whose projection to the first factor $``$ is also a discrete set of contractible open subsets of $``$. See also theorem 2.3 in \[Bu\].
Gregery T. Buzzard
Department of Mathematics
Cornell University
Ithaca, NY 14853
USA
Steven Shin-Yi Lu
Department of Mathematics
University of Waterloo
Waterloo, Ontario, N2L3G1
Canada
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# 1 Introduction
## 1 Introduction
One of the most interest features of the charm photoproduction at HERA, which had been discovered by ZEUS Collaboration, is the correlation between jet energies in the events with the $`D^{}`$-meson .
It was shown experimentally that in the kinematical region $`130<W<280`$ GeV, $`Q^2<1\mathrm{GeV}^2`$ (where $`W`$ is the energy of the $`\gamma p`$-interaction, and $`Q^2`$ is the photon virtuality) in $`30`$% of all events with $`D^{}`$-mesons two jets with the highest transverse energy take away only a half of the photon energy. The distribution over $`x_\gamma ^{\mathrm{OBS}}`$ has been used for the quantitative analysis of this phenomenon. The value of $`x_\gamma ^{\mathrm{OBS}}`$ means the part of the photon energy which is taken away by two jets with the highest transverse energy. It is defined as below:
$$x_\gamma ^{\mathrm{OBS}}=\frac{_{\mathrm{jet}=1,2}E_T^{\mathrm{jet}}e^{\eta _{\mathrm{jet}}}}{2E_ey},$$
(1)
where the $`\eta ^{\mathrm{jet}}`$ – pseudorapidity of the jet, $`E_T^{\mathrm{jet}}`$– transverse energy of the jet. One takes the sum over two jets with the highest $`E_T^{\mathrm{jet}}`$.
The following equation allows to clarify the physical meaning of $`x_\gamma ^{\mathrm{OBS}}`$:
$$2E_ey=\underset{\mathrm{jet}}{\overset{\mathrm{all}}{}}(E^{\mathrm{jet}}p_z^{\mathrm{jet}})=\underset{\mathrm{jet}}{\overset{\mathrm{all}}{}}E_T^{\mathrm{jet}}e^{\eta _{\mathrm{jet}}}.$$
(2)
One can see that for events with two jets $`x_\gamma ^{\mathrm{OBS}}1`$. For events with many jets always $`x_\gamma ^{\mathrm{OBS}}<1`$.
In the framework of the pQCD there are two types of contribution into charm photoproduction:
1) contribution from the scattering of the parton from the hadron on the photon (”direct photon”);
2) contribution from the scattering of the parton from hadron on the parton from photon (”resolved photon”).
For the ”direct photon” events $`x_\gamma ^{\mathrm{OBS}}1`$. This effect can be easily understood if one keeps in mind that in the LO such events have only two jets.
In the ”resolved photon” events or, by other words in the interaction of the hadronic parton with the charm quark from photon, the charm quark carries only a part of the photon energy, and that is why $`x_\gamma ^{\mathrm{OBS}}`$ can be essentially less than 1.
In the NLO approach one can not distinguish between the contributions from ”direct photon” and ”resolved photon” on the level of the cross-sections, and NLO matrix element includes both of these contributions .
Nevertheless, these contributions have a different kinematical behavior and can be distinguished one from another by using the distribution over $`x_\gamma ^{\mathrm{OBS}}`$. So, ZEUS collaboration divides the $`x_\gamma ^{\mathrm{OBS}}`$ range into two parts:
1) $`0<x_\gamma ^{\mathrm{OBS}}<0.75`$ (”small” $`x_\gamma ^{\mathrm{OBS}}`$) which corresponds to the ”resolved photon” events;
2) $`0.75<x_\gamma ^{\mathrm{OBS}}<1`$ (”large” $`x_\gamma ^{\mathrm{OBS}}`$) which corresponds to the ”direct photon” events.
As it follows from the ZEUS data, for the kinematical region $`p_T^D^{}>3`$ GeV, $`|\eta ^D^{}|<1.5`$, $`130<W<280`$ GeV, $`Q^2<1\mathrm{GeV}^2`$, $`|\eta ^{\mathrm{jet}}|<2.4`$, $`E_T^{\mathrm{jet}_1}>7`$ GeV, $`E_T^{\mathrm{jet}_2}>6`$ GeV a noticeable discrepancy between the NLO predictions for the distribution over $`x_\gamma ^{\mathrm{OBS}}`$ and the experimental one exits . The experimental value of the cross-section in the region $`0<x_\gamma ^{\mathrm{OBS}}<1`$ (region of ”resolved photon”) is larger than that predicted by the NLO approach.
In this paper we will try to describe the distribution over $`x_\gamma ^{\mathrm{OBS}}`$ by taking into account the hadronic component of the photon in the framework of the Vector Dominance Model (VDM) .
## 2 Description of the model and results
First of all, it is worth to mention that we have to take into account only photon charm quark fluctuation, because we are interested in the events with the charm jets only. In our calculation we suggest that the main contribution into the $`D^{}`$-meson photoproduction at HERA in the region of ”small” $`x_\gamma ^{\mathrm{OBS}}`$ is due to the scattering of gluon on the $`c`$-quark from the $`J/\psi `$, $`\psi ^{}`$ and other vector $`c\overline{c}`$-mesons from photon in the wavefunction of the initial photon.
In the case under consideration the photon wave function can be approximated by a sum of the wave functions of the vector mesons as bellow :
$$|\gamma =\underset{V=\psi ,\psi ^{},\psi (3770)\mathrm{}}{}\frac{4\pi \alpha }{\gamma _V^2}|V,$$
(3)
where a $`\gamma _V`$ constant is determined by the vector meson width of decay into the $`e^+e^{}`$-pair:
$$\mathrm{\Gamma }(Ve^+e^{})=\frac{\alpha ^2}{3}\frac{4\pi }{\gamma _V^2}M_V$$
(4)
From the equations above the charm photoproduction cross section can be expressed through cross-sections of the vector meson interaction with hadron:
$$\sigma (\gamma Nc\overline{c})=\underset{V=\psi ,\psi ^{},\psi (3770)\mathrm{}}{}\frac{4\pi \alpha }{\gamma _V^2}\sigma _{\mathrm{inel}.}(VN)$$
(5)
As it was mentioned above, one should take into account only the $`c\overline{c}`$-mesons contribution, as the contribution from other vector mesons is small .
So, in the framework of VDM, the following equation can be written for the charm photoproduction with two large transverse energy jets:
$$\sigma (\gamma p)=\underset{V=\psi ,\psi ^{},\psi (3770)\mathrm{}}{}\frac{4\pi \alpha }{\gamma _V^2}K(E_\gamma ,m_V,m_p)\sigma (Vp),$$
(6)
where
$$\sigma (Vp)=f_c^V(x_\gamma ^{\mathrm{OBS}})f_g^N(x_g)\sigma (cg2\mathrm{jet}_{\mathrm{higt}\mathrm{E}_\mathrm{T}})𝑑x_\gamma ^{\mathrm{OBS}}𝑑x_g,$$
(7)
and
$$K(E_\gamma ,m_V,m_p)=\sqrt{\frac{E_V^2m_p^2}{E_\gamma ^2}}.$$
(8)
To evaluate the $`\sigma (cg2\mathrm{jet}_{\mathrm{higt}\mathrm{E}_\mathrm{T}})`$ we use the Born approximation for the $`cgcg`$ elastic scattering. So, the cross-section for the events with two high $`E_T`$ jets is determined by the hard scattering of the gluon on the $`c`$-quark from vector meson, as it has been shown in Fig. 1.
For the $`p_T`$-distribution of $`D^{}`$ we convolute this cross section with the fragmentation function of $`c`$-quark into $`D^{}`$-meson from .
For the $`c`$-quark distribution in $`J/\psi `$ and any other vector $`c\overline{c}`$-meson we use the Regge parametrization :
$$f_c^\mathrm{\Psi }(x)=Nx^{\alpha _\mathrm{\Psi }}(1x)^{\gamma \alpha _\mathrm{\Psi }},$$
(9)
where $`x`$ is the $`J/\mathrm{\Psi }`$ momentum fraction of $`c`$-quark, $`\alpha _\psi `$ is the intercept of $`J/\psi `$, $`N`$ is the normalization coefficient and $`\gamma =1/4`$.
The best description of the charm production at small $`p_T`$ can be achieved for $`\alpha _\psi =2.2`$ . With this value of $`\alpha _\psi `$ the parametrization has the following form
$$f_c^\mathrm{\Psi }(x)=49.8x^{2.2}(1x)^{2.45}.$$
(10)
It is worth to mention that this quark distribution is a complete analog of the valence quark distribution in $`\rho `$-meson:
$$f_q^\rho (x)=Nx^{\alpha _\rho }(1x)^{\gamma \alpha _\rho }.$$
(11)
The only difference is the different values for Regge intercept and parameter $`\gamma `$.
Average fraction of the total momentum carred by $`c`$-quarks is close to 1:
$$x_c+x_{\overline{c}}0.96,$$
(12)
and each quark get approximately a half of the photon momentum.
So, we get the result we needed: in average, only a half of the photon momentum participates in the hard $`cg`$-interaction followed by the production of two jets with large transverse energies. Namely such events give the contribution into region of ”small” $`x_\gamma ^{\mathrm{OBS}}`$.
The characteristic time of the $`c\overline{c}`$-fluctuation is about $`1/m_\psi `$ in the rest frame of fluctuation. This value multiplied by the Lorence factor $`E_\gamma /m_\psi `$ is essentially larger than the characteristic time of the hard jet production:
$$\frac{E_\gamma }{m_\mathrm{\Psi }^2}\frac{1}{E_T^{\mathrm{jet}}}.$$
(13)
This circumstance allows us to calculate the charm production cross-section incoherentely, summarizing the contributions from hard scattering on each $`c`$-quark. It is clear that for small $`E_T`$ such approach is not valid and it is worth to consider the $`c\overline{c}`$-meson as a color dipole and take into account transverse quark motion and interference between different contributions. In this kinematical region the dipole model gives a good description of the data.
In this paper we limited ourselves to the consideration of the incoherent contribution and describe the distributions in $`p_T`$ and $`\eta `$ for $`p_T>4`$ GeV.
The cuts on the jets transverse energy for the distribution in $`x_\gamma ^{\mathrm{OBS}}`$ ($`E_T^{\mathrm{jet}_1}>7`$ GeV and $`E_T^{\mathrm{jet}_2}>6`$ GeV) ensure large transverse momentum of $`c`$-quark before its fragmentation into $`D^{}`$-meson (or in other words after hard scattering $`cgcg`$). That is why we use our model for all investigated values of $`p_T`$ to predict the $`x_\gamma ^{\mathrm{OBS}}`$ distribution.
Our model does not contain any new parameters. Indeed:
1) the photon coupling constants with $`J/\psi `$ and other vector $`c\overline{c}`$-mesons are known from the width of meson decay into $`e^+e^{}`$;
2) the $`c`$-quark distribution in the vector $`c\overline{c}`$-meson (9) contains the Regge intercept $`\alpha _\psi (0)`$ which determines fragmentation function of $`c`$-quark into $`D^{}`$, and value of $`\alpha _\psi (0)=2.2`$ allows to describe the experimental data on $`c`$-quark fragmentation into $`D^{}`$-meson and charm photoproduction at low energies ;
3) the scale $`\mu _R`$ in the determination of the strong coupling constant $`\alpha _s(\mu _R)`$ in the hard scattering matrix element for the process $`cgcg`$ and the scale $`\mu _F`$ for the gluonic structure function are common parameters for such calculations. We use the following values for the scales: $`\mu _R=\mu _F=2m_D^{}`$.
It is worth to mention that the main contribution into the charm photoproduction calculated in the framework of VDM is due to $`J/\psi `$ (about 60%). Other vector $`c\overline{c}`$-mesons give 40% of the total cross-section.
In the frame of the model under consideration the ZEUS experimental cuts $`E_T^{jet_1}>6`$ GeV, $`E_T^{jet_2}>7`$ GeV are equivalent to the following ones: $`E_T^{jet_1},E_T^{jet_2}>7`$ GeV. The model does not account for the initial transverse momentum of the $`c`$-quark into the $`c\overline{c}`$-fluctuation, which is about $`1`$ GeV. That is why the both high energy jets have the same transverse energy.
To evaluate the cross-section uncertainty which is due to the transverse momentum of the initial $`c`$-quark we have calculated the distribution in $`x_\gamma ^{OBS}`$ for $`E_T^{jet_1},E_T^{jet_2}>7`$ GeV and for $`E_T^{jet_1},E_T^{jet_2}>6`$ GeV. The difference between two this calculations estimates the theory uncertainty.
One can see from Fig. 2a,b that this uncertainty is rather large.
As one can see from Fig. 2a the description of the experimental data on $`x_\gamma ^{\mathrm{OBS}}`$-distribution in the region of the ”resolved photon” has been essentially improved by adding the VDM predictions with the BKL ones .
The slightly better description of the data can be achieved by adding the VDM predictions to NLO ones from (see Fig. 2b). However, one must keep in mind that the adding of the VDM contribution can dramatically enlarge the normalization of the distribution in transverse momentum of the $`D^{}`$-meson for the model .
It is important to mention that hard interaction with the hadronic component of the photon ($`J/\psi `$, $`\psi ^{}`$ and so on) gives the noticeable contribution into the region of large $`p_T`$ (see Fig. 3). In other words, the attempt to describe $`x_\gamma ^{\mathrm{OBS}}`$ distribution in the framework of VDM leads to additional contribution in the region of large $`p_T`$.
So, the following conclusion can be drawn: the nonperturbative $`c\overline{c}`$-fluctuations of photon are important for the charm photoproduction. Such fluctuations are described in the framework of VDM and can not be described by simple formula of the pQCD
$$f_c(x)(x^2+(1x)^2)$$
(14)
or improved formula :
$$f_c(x)\left(x^2+(1x)^2+\frac{2m_c^2}{E_T^2}x(1x)\right)$$
(15)
or Bethe-Heitler formula :
$$\begin{array}{ccc}\hfill f_c(x)& =& \frac{4\alpha }{3\pi }[\beta (8x(1x)1\frac{4m_c^2}{Q^2}x(1x))\hfill \\ & & +(x^2+(1x)^2+\frac{4m_c^2}{Q^2}x(13x)\frac{8m_c^4}{Q^4}x^2)\mathrm{ln}\left(\frac{1+\beta }{1\beta }\right)],\hfill \end{array}$$
(16)
if $`\beta ^2=1\frac{4m_c^2x}{(1x)Q^2}>0`$ and $`f_c(x)=0`$ if $`\beta ^2<0`$.
This part of the fluctuations has a typical hadronic structure (9) and such fluctuations provide about 30% in the cross-section distribution in $`x_\gamma ^{\mathrm{OBS}}`$ for the $`D^{}`$-meson photoproduction.
The topology of the events in the region of the ”small” $`x_\gamma ^{\mathrm{OBS}}`$ essentially differs from one of the events in the region of the ”large” $`x_\gamma ^{\mathrm{OBS}}`$. In the latter case both jets carry a charm quark, in contradiction with the former case, where only one jet has charm quark. This circumstance leads to the large azimuthal correlation between charmed particles for the ”direct photon” and to negligible one for ”resolved photon” (see Fig. 4).
## 3 Conclusion
It has been shown that scattering of virtual vector $`c\overline{c}`$-mesons from photon on the proton gives noticeable contribution into the $`D^{}`$-meson photoproduction for all values of the transverse momentum investigated at HERA. Furthermore this contribution allows essentially improve the description of the experimental distribution in $`x_\gamma ^{\mathrm{OBS}}`$ in the region of the ”resolved photon” ($`0<x_\gamma ^{\mathrm{OBS}}<0.75`$).
We are grateful L. Gladilin, V Kiselev, D. Kharzeev, I. Korzhavina for the useful discussion of the materials presented in this article.
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# DISPERSION RELATIONS IN ULTRADEGENERATE RELATIVISTIC PLASMAS
## I Introduction
The study of hot relativistic plasmas is nowadays a very active field of research . This is due to the existence of experimental programs to test the existence of the quark-gluon plasma phase of QCD. The physics of some astrophysical settings, such as those of neutron stars and supernovas, also requires knowledge of a different regime of relativistic plasmas, less hot but still very dense. This cold and ultradegenerate regime of QED and QCD has been much less explored. However, the fact that matter at very high baryonic densities behaves as a color superconductor has given us a strong motivation to study ultradegenerate relativistic plasmas, as several new phenomena occurs in this phase of QCD (see and , and references therein).
One of the central concepts in a plasma is that of a quasiparticle. A particle immersed in a medium modifies its propagation properties by interacting with the surrounding medium. In field theoretical language, we would say that the particle is “dressed” by a self-energy cloud. In the ultradegenerate plasma the relevant degrees of freedom are those of quasiparticles or quasiholes (absences of particles in the Fermi sea) living close to the Fermi surface. Because of the exclusion principle, quasiparticles/quasiholes can only live if they are outside/inside the Fermi sea, respectively. These excitations tend to lower their energy, by undergoing collisions with the particles in the Fermi sea. They decay, and thus have a finite lifetime. The concept of quasiparticle, however, only makes sense if its lifetime is long enough, or in other words, if its damping rate is much smaller than its energy.
Here we will mainly be concerned with electromagnetic plasmas. There is a vast literature on the quasiparticle properties in non-relativistic cold plasmas . The same does not hold true for the relativistic ones, though. There are two main differences in these two energy regimes of a plasma. In the non-relativistic domain, the electric interactions are dominant, while the magnetic ones are suppressed by a factor $`(v/c)^2`$, where $`v`$ is the velocity of the particle, and $`c`$ is the velocity of light. Thus, magnetic interactions start to be relevant only when quasiparticles are fast enough or, in other words, when the Fermi velocity $`v_F`$ approaches the velocity of light. This is an important difference, as electric interactions are not long-ranged in the medium, because of Debye screening, while magnetic interactions are. The relevance of this last point has already been stressed in the condensed matter literature , just noticing that magnetic interactions spoil the normal Fermi liquid behavior of the plasma. The second main difference is due to the fact that in a relativistic plasma there are also antiparticle excitations. Their contribution to any physical process is in general suppressed, since it takes more energy to excite an antiparticle than a particle of the Fermi-Dirac sea. Nevertheless, in the context of the color flavor locking phase of QCD , some of the properties of the antiparticles determine the mass spectrum of the Goldstone modes which arise from the spontaneous breaking of chiral symmetry, so those cannot be neglected.
In this paper we study the quasiparticle and antiquasiparticle dispersion relations in a full relativistic framework, generalizing the results of a previous publication to the case where the Fermi velocity $`v_Fc`$. We can thus explore all the energy domains of the system, and in particular, we can take the non-relativistic limit $`v_Fc`$, and match the results obtained in the condensed matter literature . The dispersion relations are obtained by computing the on-shell one-loop self-energy. While the one-loop self-energy is in general gauge-dependent, it is not when evaluated on the particles mass-shell. For quasiparticles with momenta close to the Fermi momentum, the one-loop self-energy is dominated by a diagram in which the photon is soft. When the photon is soft, it also needs to be dressed to take properly into account the effects of the medium. This can be done by using the resummation techniques proposed by Braaten and Pisarski , and considering hard thermal loop photon propagators, or hard dense loop (HDL) ones for the ultradegenerate case . We first compute the on-shell imaginary part of the one-loop self-energy for electrons and positrons, which can be interpreted in terms of their scattering with particles of the Fermi sea, via an exchange of a soft photon. The on-shell real part of the self-energy can be reconstructed from the on-shell imaginary part, just by using a Kramers-Kroning relation.
This paper is structured as follows. Section II introduces the notation of the paper. We work in natural units, $`c=\mathrm{}=k_B=1`$, unless otherwise stated. In Sect. IIIA we compute the on-shell one-loop self-energy of the fermion. We take the non-relativistic limit of our results in Sect. IIIB, and conclude in Section IV. In Appendix A the spectral functions of the HDL photon propagators are given for $`v_Fc`$, and in Appendix B Luttinger’s theorem is recalled.
## II Dispersion Relations for the quasiparticles
We consider a plasma with a finite density of electrons, characterized by a chemical potential $`\mu `$. In order to guarantee its stability, we assume that the electrons are immersed in a uniform background of positive charges, of density equal to the average electron density. These background charges can be due to positively charged ions, which are very heavy.
In a plasma with chemical potential $`\mu `$, the propagation properties of the quasiparticles are modified by medium effects. The dressed fermion propagator $`S(P)`$, where $`P=(p_0,𝐩)`$ is the four momentum, obeys the Schwinger-Dyson equation
$$S^1(P)=S_0^1(P)+\mathrm{\Sigma }(P),$$
(1)
where $`S_0^1`$ is the inverse free propagator
$$S_0^1(P)=P\text{/}+\mu \gamma _0m,$$
(2)
with $`P\text{/}=P^\mu \gamma _\mu `$ and $`\mathrm{\Sigma }(P)`$ is the one-loop self-energy.
Because of the clear asymmetry between electrons and positrons in the electromagnetic plasma, it is convenient to treat them separately, as their propagation properties will be modified in different ways. Introducing the positive and negative energy projectors
$$\mathrm{\Lambda }_𝐩^\pm =\frac{E_p\pm \left(\gamma _0\gamma 𝐩+m\gamma _0\right)}{2E_p},$$
(3)
where $`E_p=\sqrt{p^2+m^2}`$, we can rewrite
$`S_0^1(P)`$ $`=`$ $`\gamma _0\mathrm{\Lambda }_𝐩^+\left(p_0+\mu E_p\right)+\gamma _0\mathrm{\Lambda }_𝐩^{}\left(p_0+\mu +E_p\right),`$ (4)
$`\mathrm{\Sigma }(P)`$ $`=`$ $`\gamma _0\mathrm{\Lambda }_𝐩^+\mathrm{\Sigma }_+(P)\gamma _0\mathrm{\Lambda }_𝐩^{}\mathrm{\Sigma }_{}(P).`$ (5)
After inverting (1) one gets
$$S(P)=S_+(P)\mathrm{\Lambda }_𝐩^+\gamma _0+S_{}(P)\mathrm{\Lambda }_𝐩^{}\gamma _0,$$
(6)
where
$$S_\pm (P)=\frac{1}{p_0+\mu \left(E_p\mathrm{\Sigma }_\pm (P)\right)},$$
(7)
and the upper/lower subscripts refer to electrons/positrons, respectively.
Every energy eigenstate can be projected onto states of given helicity, with the projectors
$$𝒫^\pm (𝐩)=\frac{1\pm \gamma _5\gamma _0\gamma \widehat{𝐩}}{2}.$$
(8)
In principle, the most general structure of the one-loop self-energy contains four unknown functions, according to the energies and helicities of the quasiparticles. However, the effects which will be discussed in this article do not depend on the helicity of the quasiparticles, and thus we would not explicitly take into account the helicity projectors.
The value of the one-loop self-energy $`\mathrm{\Sigma }`$ is gauge-dependent. However, when it is evaluated on the particles mass-shell, it should be gauge independent. This is so because the poles of (7) give the physical dispersion relations of electrons and positrons which define their propagation properties in the plasma.
The dispersion relations obtained from (7) are
$`\omega _\pm `$ $`=`$ $`\mu \pm \left(E_p\mathrm{Re}\mathrm{\Sigma }_\pm (\omega _\pm +i\gamma _\pm ,𝐩)\right),`$ (9)
$`\gamma _\pm `$ $`=`$ $`\mathrm{Im}\mathrm{\Sigma }_\pm (\omega _\pm +i\gamma _\pm ,𝐩),`$ (10)
where $`\omega _\pm `$ and $`\gamma _\pm `$ define the energy and damping rates of the electrons/positrons, respectively. For the concept of quasiparticle to make sense, it is necessary that $`\gamma _\pm \omega _\pm `$, so that the quasiparticles are long-lived enough.
In the remaining part of the paper the dispersion relations for quasiparticles and antiquasiparticles with momentum close to the Fermi momentum will be studied. In this case, the dominant contribution to their one-loop self-energy arises when the photon in the loop is soft, that is, of order $`e\mu `$, where $`e`$ is the electromagnetic coupling constant. When the photon is soft it has also to be dressed, in order to take into account properly the medium effects of Debye screening and Landau damping.
## III The on-shell fermion Self-Energy
### A Relativistic Domain
For a plasma at temperature $`T`$ and chemical potential $`\mu `$, we compute the one-loop self-energy $`\mathrm{\Sigma }`$ using the imaginary time formalism. It is convenient to use the spectral function representation of the fermion and photon propagators in the computation. The free fermion propagator is given by
$$S_0(i\omega _n,𝐤)=_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{d}k_0}{2\pi }\frac{(K\text{/}+m)\rho _f(K)}{k_0i\omega _n\mu },$$
(11)
with
$$\rho _f(K)=\frac{\pi }{E_k}\left(\delta (k_0E_k)\delta (k_0+E_k)\right).$$
(12)
In (11), $`\omega _n=\pi (2n+1)T`$ is a fermionic Matsubara frequency. The (resummed) photon propagator $`\mathrm{\Delta }_{\mu \nu }(Q)`$, where $`Q=(i\omega _s,𝐪)`$, and $`\omega _s=2\pi sT`$ is a bosonic Matsubara frequency, is written in the Coulomb gauge
$$\mathrm{\Delta }_{\mu \nu }(Q)=\delta _{\mu 0}\delta _{\nu 0}\mathrm{\Delta }_L(Q)+𝒫_{\mu \nu }^T\mathrm{\Delta }_T(Q)+\xi _C\frac{Q_\mu Q_\nu }{q^4},$$
(13)
where $`𝒫_{ij}^T=(\delta _{ij}\widehat{q}_i\widehat{q}_j)`$, $`\widehat{q}^i=𝐪^i/|𝐪|`$, $`𝒫_{i0}^T=𝒫_{0i}^T=𝒫_{00}^T=0`$, and $`\xi _C`$ is the gauge parameter. The longitudinal and transverse propagators are written in terms of their spectral functions
$`\mathrm{\Delta }_L(i\omega _s,q)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}q_0}{2\pi }}{\displaystyle \frac{\rho _L(q_0,q)}{q_0i\omega _s}}{\displaystyle \frac{1}{q^2}},`$ (15)
$`\mathrm{\Delta }_T(i\omega _s,q)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}q_0}{2\pi }}{\displaystyle \frac{\rho _T(q_0,q)}{q_0i\omega _s}}.`$ (16)
Analytical expressions for $`\rho _{L,T}`$ can be found in for the case of an ultrarelativistic ($`m=0`$) plasma. At $`T=0`$, it is also possible to derive the spectral functions for $`m0`$ . We present analytical expressions for the spectral functions in this case in Appendix A.
The one-loop self-energy
$$\mathrm{\Sigma }(P)=e^2T\underset{s}{}\frac{\mathrm{d}^3q}{(2\pi )^3}\gamma _\mu S_0(PQ)\gamma _\nu \mathrm{\Delta }_{\mu \nu }(Q),$$
(17)
when expressed in terms of the spectral functions, reads
$`\mathrm{\Sigma }(i\omega ,p)`$ $`=`$ $`e^2T{\displaystyle \underset{n}{}}{\displaystyle \frac{\mathrm{d}^3q}{(2\pi )^3}_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{d}k_0}{2\pi }\rho _f(K)\gamma _\mu (K\text{/}+m)\gamma _\nu }`$ (18)
$`\times `$ $`\{\left({\displaystyle _{\mathrm{}}^{\mathrm{}}}{\displaystyle \frac{\mathrm{d}q_0}{2\pi }}{\displaystyle \frac{\delta _{\mu 0}\delta _{\nu 0}\rho _L(q_0,q)+𝒫_{\mu \nu }^T\rho _T(q_0,q)}{\left(q_0i\omega _n\right)\left(k_0i\omega +i\omega _n\mu \right)}}\right){\displaystyle \frac{1}{q^2}}{\displaystyle \frac{\delta _{\mu 0}\delta _{\nu 0}}{\left(k_0i\omega +i\omega _n\mu \right)}}`$ (19)
$`+`$ $`\xi _C{\displaystyle \frac{Q_\mu Q_\nu }{q^4}}\}.`$ (20)
The sum over Matsubara frequencies is now easily performed. After analytical continuation $`i\omega _n+\mu p_0+i\eta `$, with $`\eta 0^+`$ to Minkowski space, one can evaluate the on-shell imaginary part. It is very easy to realize that none of the last two pieces of Eq. (18) contribute to this on-shell imaginary part. Therefore, the result of the computation is gauge independent. One finds
$`\mathrm{Im}\mathrm{\Sigma }(p_0+i\eta ,p)`$ $`=`$ $`e^2\mathrm{Im}{\displaystyle \frac{\mathrm{d}^3q}{(2\pi )^3}_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{d}k_0}{2\pi }\rho _f(K)_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{d}q_0}{2\pi }\frac{1+f(q_0)\stackrel{~}{f}(k_0\mu )}{p_0k_0q_0+i\eta }}`$ (21)
$`\times `$ $`\gamma _\mu (K\text{/}+m)\gamma _\nu \left[\delta _{\mu 0}\delta _{\nu 0}\rho _L(q_0,q)+𝒫_{\mu \nu }^T\rho _T(q_0,q)\right].`$ (22)
In (21), $`f`$ and $`\stackrel{~}{f}`$ are Bose-Einstein and Fermi-Dirac distribution functions ($`\beta =1/T`$)
$$f(q_0)=\frac{1}{\mathrm{e}^{\beta q_0}1},\stackrel{~}{f}(k_0\mu )=\frac{1}{\mathrm{e}^{\beta (k_0\mu )}+1}.$$
(23)
The damping rates for the quasiparticles and antiquasiparticles are thus obtained after multiplying Eq. (21) by the corresponding projectors and taking a Dirac trace, evaluating the final expression on the particles mass-shellAt this point, one can check that the damping rate of the quasiparticles does not depend on their helicities, by using the helicity projectors of Eq.(8).
$$\gamma _\pm =\mathrm{Im}\mathrm{\Sigma }_\pm (p_0+i\eta ,𝐩)|_{p_0\mathrm{on}\mathrm{shell}}.$$
(24)
To obtain the damping rate of a quasiparticle one has to evaluate the imaginary part of its self-energy on the pole of the dressed propagator. However, up to corrections of order $`e^2`$, it would be enough to consider the above expressions at $`p_0=\pm E`$, as the corrections introduced by $`\mathrm{\Sigma }_\pm `$ only displace these poles by an amount proportional to $`e^2`$. In this case, after evaluating the Dirac traces we find, with $`𝐤=𝐩𝐪`$,
$`\gamma _\pm `$ $`=`$ $`\pm {\displaystyle \frac{\pi e^2}{E}}{\displaystyle \frac{\mathrm{d}^3q}{(2\pi )^3}_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{d}k_0}{2\pi }\rho _f(k_0)_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{d}q_0}{2\pi }}`$ (25)
$`\times `$ $`(1+f(q_0)\stackrel{~}{f}(k_0\mu ))\delta (p_0k_0q_0)\{[p_0k_0+𝐩𝐤+m^2]`$ (26)
$`\times `$ $`\rho _L(q_0,q)+2[p_0k_0(𝐩\widehat{𝐪})(𝐤\widehat{𝐪})m^2]\rho _T(q_0,q)\left\}\right|_{p_0=\pm E}.`$ (27)
Equation (25) gives the general expression for the damping rates for any value of $`T`$, $`\mu `$ and $`m`$. At very high temperature, these damping rates are infrared (IR) logarithmic divergent, even after including the screening corrections. This is due to the soft photon contribution, as $`f(q_0)T/q_0`$ for $`q_0T`$. Perturbation theory fails to provide the damping rates at high $`T`$ . A non-perturbative treatment to resum the leading order divergences was proposed in to find a non-exponential decay law for the quasiparticles.
In this paper we are concerned with the ultradegenerate limit, when $`T=0`$. In this case several simplifications occur. The damping rates are IR finite after the inclusion of the screening effects , as opposed to what happens at high $`T`$. For $`T=0`$, $`(1+f(q_0))=\mathrm{\Theta }(q_0)`$, where $`\mathrm{\Theta }`$ is the step function. For $`T=0`$ the fermion distribution function is $`\stackrel{~}{f}(E_k\mu )=\mathrm{\Theta }(\mu E_k)`$, while $`\stackrel{~}{f}(E_k\mu )=1\stackrel{~}{f}(E_k+\mu )=1`$.
From this point on, it is convenient to treat separately the electron and positron damping rates, as different phase-space restrictions arise in the two cases. If we concentrate in the soft photon region, we can approximate
$$E_k=\sqrt{|𝐩𝐪|^2+m^2}E𝐯𝐪,$$
(28)
where $`𝐯=𝐩/E`$ is the velocity of the fermion. We thus find
$`\gamma _+`$ $``$ $`{\displaystyle \frac{e^2}{8\pi ^2v}}{\displaystyle _{q\mathrm{soft}}}qdqdq_0\left[\left(\mathrm{\Theta }(q_0)\mathrm{\Theta }(\mu E+q_0)\right)\left\{\rho _L(q_0,q)+v^2(1\mathrm{cos}^2\theta )\rho _T(q_0,q)\right\}\right],`$ (30)
$`\gamma _{}`$ $``$ $`{\displaystyle \frac{e^2}{8\pi ^2v}}{\displaystyle _{q\mathrm{soft}}}qdqdq_0\left[\mathrm{\Theta }(q_0)\left\{\rho _L(q_0,q)+v^2(1\mathrm{cos}^2\theta )\rho _T(q_0,q)\right\}\right],`$ (31)
where $`q_0=qv\mathrm{cos}\theta `$.
The damping rates for the electron and the positron thus only differ in the phase-space restrictions of these two types of particles. One can interpret the above equations as follows. A particle/antiparticle, with energy $`\pm E`$ is scattered to a state of energy $`\pm E_k`$, respectively, creating a particle-hole pair. For the electron, $`E_k`$ is forced to be above the Fermi energy, because of Pauli blocking. This last restriction is absent in the case of the positron.
For a quasiparticle with velocity close to the Fermi velocity, we can further approximate $`vv_F`$ in Eq. (30). From the fact that the spectral functions $`\rho _{L,T}`$ in Eqs. (28) are evaluated for values of $`q_0^2q^2v_F^2`$, we see that it is only the part of the spectral function corresponding to Landau damping (the functions $`\beta _{L,T}`$ in (A)), that contributes to the integrals. Using the explicit values of spectral densities as given in Appendix A, one can evaluate the above integrals numerically. Analytical expressions can be obtained for the interesting case $`|E\mu |M`$ (so this includes the case of quasiholes). In this regime we find at leading order
$$\gamma _+\frac{e^2}{24\pi }|E\mu |+\frac{e^2}{64v_F^2M}(E\mu )^2+𝒪(|E\mu |^3),$$
(32)
which generalizes the expressions obtained in Refs. for the case $`v_F1`$. In the above equations $`M=\sqrt{e^2\mu ^2v_F/\pi ^2}`$ is the Debye mass. The first terms in the r.h.s. of Eqs. (32) are due to scattering processes with exchange of soft magnetic photons, while the second is due to the exchange of soft electric ones. As can be seen the magnetic contribution is suppressed with respect to the electric one by a factor $`v_F^2`$. Therefore, the electric contribution is dominant for $`v_F1`$. In the ultrarelativistic limit, $`v_F=1`$, the damping rate of the electron is dominated by the magnetic contribution.
The damping rate of a quasiparticle, or a quasihole, which lives close to the Fermi surface can then be expressed as a power series in $`|E\mu |`$. If this parameter is large, then $`\gamma _+`$ is large, and the lifetime of these excitations is so short, that it does not make sense to talk about quasiparticles or quasiholes. On the contrary, when the fermion energy approaches the Fermi energy, its lifetime tends to infinity. In particular, from Eq. (32) one deduces that the Fermi sea is stable. We should point out also the very different contribution to the energy dependence of $`\gamma _+`$ from the electric and magnetic interactions. The quadratic dependence on $`(E\mu )`$ of $`\gamma _+`$ can be entirely understood as arising from the short-ranged character of the electric interactions in the plasma, and also the phase-space restrictions of electron-electron scattering (see Appendix B). Magnetic interactions are not short-ranged, but only suffer a weak dynamical screening due to Landau damping. The linear dependence on $`(E\mu )`$ is also a product of Landau damping and phase-space restrictions.
We now consider the damping rate of the antiquasiparticle. Let us first stress that for a positron pair annihilation also contributes to its damping rate. However, this is a process that occurs at order $`e^4`$, and it can be computed by taking the imaginary part of a two-loop correction to the fermion self-energy. In a weak coupling expansion, the damping rate of the positron is dominated by the scattering of the positron with the electrons of the Fermi sea. For a positron with velocity $`vv_F`$, we can evaluate numerically Eq. (31). The only difference with respect to the computation of $`\gamma _+`$ comes from the different phase-space restrictions for antifermions, or in other words, the different domain of integration of the integrals. We find at leading order
$$\gamma _{}e^2M\left(\frac{v_F^2}{24\pi }+\frac{1}{64}\right),$$
(33)
which agrees for $`v_F=1`$ with the result of Ref..
The on-shell real part of the self-energy can be obtained from the general expression (18), just by using a principal value prescription to evaluate the integral after the analytical continuation to Minkowski space is done. However, it is much simpler to reconstruct it from the value of the on-shell imaginary part, using a Kramers-Kroning dispersion relation, which gives the value of the real part, up to a constant.
If $`f_\pm (\omega )`$ is an analytic function in the upper/lower complex plane, respectively, then from the Cauchy theorem, its real part is given as a function of its imaginary part as
$$\mathrm{Re}f_\pm (\omega )=\pm \frac{PP}{\pi }_{\mathrm{}}^{\mathrm{}}\frac{\mathrm{d}\omega ^{}}{\omega ^{}\omega }\mathrm{Im}f_\pm (\omega ^{})+C_{\mathrm{}},$$
(34)
where $`PP`$ denotes the principal value of the integral along the real axis, and $`C_{\mathrm{}}`$ is a subtraction constant needed in case that $`f_\pm `$ does not vanish for $`|\omega |\mathrm{}`$.
Since we have computed the damping rate $`\gamma _+`$ for values $`|E\mu |M`$, the only energy domain where the concept of quasiparticle makes sense, we will use the dispersion relation using a cutoff which implements this constraint, that is with cutoffs $`\mathrm{\Lambda }_\pm =\mu \pm M`$. We thus find
$$\mathrm{Re}\mathrm{\Sigma }_+(E,p)\mathrm{Re}\mathrm{\Sigma }_+(\mu ,p)+\frac{e^2}{12\pi ^2}(E\mu )\mathrm{ln}\frac{M}{|E\mu |}+\frac{e^2}{32\pi v_F^2}|E\mu |+𝒪((E\mu )^2).$$
(35)
The value of the energy-independent constant $`\mathrm{Re}\mathrm{\Sigma }_+(E=\mu ,p)`$, which renormalizes the chemical potential, can only be determined from the explicit evaluation of Eq. (18).
We now use the Kramers-Kroning dispersion relation for the antifermions, also imposing a cutoff in the dispersion relation that guarantees that the momentum of the particle is not far away from the Fermi momentum. We thus find
$$\mathrm{Re}\mathrm{\Sigma }_{}(E,p)\mathrm{Re}\mathrm{\Sigma }_{}(\mu ,p)+\left(\frac{e^2v_F^2}{12\pi ^2}+\frac{e^2}{32\pi }\right)\left(\mu E\right)+𝒪((\mu E)^2).$$
(36)
The leading logarithmic behavior of the real part of the one-loop self-energy of a quark in the high baryonic limit of QCD has been obtained in the ultrarelativistic limit in Ref. . There the same leading logarithmic dependence in the energy as in Eq. (35) has been found. We should stress here that this can only be valid for the quark excitations, but not for the antiquarks ones.
### B Non-Relativistic Limit
In this subsection we take the non-relativistic (nr) limit of the expressions computed previously, restoring the fundamental constant $`c`$ in the equations. The nr limit corresponds to $`v_Fc`$. The antiparticles then decouple. In such a case, the contribution from the magnetic sector to the fermion self-energy is suppressed by a factor $`(v_F/c)^2`$ with respect to the electric sector one. The electric effects are thus dominant.
The lifetime of an electron $`\tau `$ is defined as $`1/2\gamma _+`$. We neglect the magnetic contribution to the damping rate, and express the electric contribution in terms of the plasma frequency $`\omega _p^2=\frac{1}{3}M^2v_F^2`$
$$\frac{1}{\tau }=\frac{\sqrt{3}\pi ^2\omega _p}{32}\left(\frac{E\mu }{\mu }\right)^2\frac{c^4}{v_F^4}.$$
(37)
The relativistic and non-relativistic chemical potentials differ by the rest mass of the particle, $`\mu ^2=\mu _{nr}^2+m^2c^4`$. For $`p^2m^2c^2`$, $`E=mc^2+ϵ_{nr}+𝒪(\frac{p^4}{m^4c^2})`$, where $`ϵ_{nr}=\frac{p^2}{2m}`$. Therefore
$$\frac{1}{\tau }\frac{\sqrt{3}\pi ^2\omega _p}{32}\frac{\left(ϵ_{nr}\mu _{nr}\right)^2}{m^2c^4}\frac{c^4}{v_F^4}=\frac{\sqrt{3}\pi ^2\omega _p}{128}\left(\frac{ϵ_{nr}\mu _{nr}}{\mu _{nr}}\right)^2,$$
(38)
where in the last equality we have used $`\mu _{nr}=ϵ_F=\frac{1}{2}mv_F^2`$. The above expression agrees with the computation of the lifetime of an electron in a non relativistic quantum liquid, using the random phase approximation (see Eq. (5.134c) of ).
Since in the nr limit the difference $`(E\mu )\left(ϵ_{nr}\mu _{nr}\right)+𝒪(\frac{p^4}{m^4c^2})`$, we also reproduce the dispersion relations due to magnetic interactions of non relativistic electrons computed in .
## IV Conclusions
We have derived the one-loop dispersion relations of quasiparticles and antiquasiparticles with momentum close to the Fermi momentum in a relativistic electromagnetic plasma, recovering in the non-relativistic limit the results of Refs. . As already emphasized in those papers, the long-ranged character of the magnetic interactions spoils the normal Fermi liquid behavior of the plasma. This effect is fully dominant when the Fermi velocity $`v_F`$ is close to the velocity of light. We have also found that the medium modifies in a different way the propagation properties of particles and antiparticles. This can be simply understood from their different phase-space restrictions when they scatter with the electrons of the Fermi sea. We should also emphasize that our results are gauge independent. This is because we have computed the one-loop self-energy on mass-shell. Off-shell, the one-loop self-energy (18) is a gauge-dependent function.
While we have concentrated our study to QED plasmas, our results can be easily transported to QCD, only by replacing the electromagnetic coupling constant by the QCD one and taking into account some additional color factors. In the superconducting phase of QCD, the dispersion relations of quarks and antiquarks would be modified in a different way according to whether or not these form Cooper pairs. The dispersion relation for antiquarks in the presence of a color gap has not yet been determined, while it is still necessary to understand how antiquarks propagate in a color superconducting medium.
Acknowledgements: I want to thank M. Le Bellac, with whom this work was initiated, M. Thoma and M. Tytgat for useful discussions.
## A Spectral functions for HDL photon propagators
The spectral functions $`\rho _{L,T}`$ for the resummed propagators $`\mathrm{\Delta }_{L,T}`$ can be found in for the case of an ultrarelativistic ($`m=0`$) plasma. In the case of a ultradegenerate plasma, they can also be determined when $`m0`$ . For completeness, we will present them below. In this case, the Debye mass is $`M^2=e^2\mu ^2v_F/\pi ^2`$, where $`v_F`$ is the Fermi velocity, defined as the ratio between the Fermi momentum and the Fermi energy, $`v_F=p_F/\mu `$. The spectral functions of the resummed HDL propagators are computed from their imaginary part
$$\rho _{L,T}(q_0,𝐪)=2\mathrm{Im}\mathrm{\Delta }_{L,T}(q_0+iϵ,𝐪).$$
(A1)
These functions can be written in terms of a contribution of the poles of the propagators, plus another one arising from Landau damping:
$$\frac{\rho _{L,T}(q_0,q)}{2\pi }=Z_{L,T}\left[\delta (q_0\omega _{L,T}(q))\delta (q_0+\omega _{L,T}(q))\right]+\beta _{L,T}(q_0,q).$$
(A2)
The poles $`\omega _{L,T}`$ are solutions of the dispersion relations
$`\omega _L^2(q)`$ $`=`$ $`\omega _p^2{\displaystyle \frac{3\omega _T^2(q)}{v_F^2q^2}}\left[{\displaystyle \frac{\omega _T(q)}{2v_Fq}}\mathrm{ln}{\displaystyle \frac{\omega _T(q)+v_Fq}{\omega _T(q)v_Fq}}1\right],0q<q_{max},`$ (A4)
$`\omega _T^2(q)`$ $`=`$ $`q^2+\omega _p^2{\displaystyle \frac{3\omega _T^2(q)}{2v_F^2q^2}}\left[1+{\displaystyle \frac{1}{2}}\left({\displaystyle \frac{v_Fq}{\omega _T(q)}}{\displaystyle \frac{\omega _T(q)}{v_Fq}}\right)\mathrm{ln}{\displaystyle \frac{\omega _T(q)+v_Fq}{\omega _T(q)v_Fq}}\right],0q<\mathrm{},`$ (A5)
where $`\omega _p^2=\frac{1}{3}M^2v_F^2`$ is the plasma frequency, and
$$q_{max}=\left(\frac{1}{2v_F}\mathrm{ln}\frac{1+v_F}{1v_F}1\right)^{\frac{1}{2}}M,$$
(A6)
in the maximum momentum at which the plasmon can propagate. The above dispersion relations have to be solved numerically. It is possible to obtain analytically the small $`q`$ behavior of their solutions. For $`q\omega _p`$
$$\omega _T^2(q)\omega _p^2+q^2\left(1+\frac{v_F^2}{5}\right),\omega _L^2(q)\omega _p^2+\frac{3}{5}v_F^2q^2.$$
(A7)
The functions $`Z_{L,T}`$ are the residues of $`\mathrm{\Delta }_{L,T}`$ evaluated at their poles, and are given by
$`Z_L(q)`$ $`=`$ $`{\displaystyle \frac{\omega _L\left(\omega _L^2v_F^2q^2\right)}{q^2\left(3\omega _p^2\left(\omega _L^2v_F^2q^2\right)\right)}},`$ (A9)
$`Z_T(q)`$ $`=`$ $`{\displaystyle \frac{\omega _T\left(\omega _T^2v_F^2q^2\right)}{3\omega _p^2\omega _T^2+\left(\omega _T^2+q^2\right)\left(\omega _T^2v_F^2q^2\right)2\omega _T^2\left(\omega _T^2q^2\right)}}.`$ (A10)
The pole contribution to the spectral functions is only non-vanishing above the light-cone. The Landau damping pieces of the spectral functions are non-vanishing only for $`q_0^2q^2v_F^2`$ and are given by
$`\beta _L(q_0,q)`$ $`=`$ $`{\displaystyle \frac{M^2x\mathrm{\Theta }(1x^2)}{2\left[q^2+M^2\left(1\frac{x}{2}\mathrm{ln}\left|\frac{x+1}{x1}\right|\right)\right]^2+\frac{M^4\pi ^2x^2}{4}}},`$ (A12)
$`\beta _T(q_0,q)`$ $`=`$ $`{\displaystyle \frac{M^2v_F^2x(1x^2)\mathrm{\Theta }(1x^2)}{\left[2q^2(x^2v_F^21)M^2x^2v_F^2\left(1+\frac{(1x^2)}{2x}\mathrm{ln}\left|\frac{x+1}{x1}\right|\right)\right]^2+\frac{M^4v_F^4\pi ^2x^2(1x^2)^2}{4}}},`$ (A13)
where $`x=q_0/qv_F`$.
## B Luttinger’s Theorem
The dependence on $`(E\mu )^2`$ of the damping rate of a fermion with energy above $`\mu `$ can be understood completely as arising from phase-space restrictions of fermion-fermion scattering, in the case where the interactions are short-ranged and repulsive. The argument, due to Luttinger , is simple. We present it below. Let us consider the decay rate of a fermion with energy $`E`$ which interacts with a fermion with energy $`E_k`$ inside the Fermi sea. As a result, two new particles appear, with energies $`E_k^{}`$ and $`E_p^{}`$, respectively, which are outside the Fermi sea. The total decay rate is then given by
$`\mathrm{\Gamma }(E)`$ $`=`$ $`{\displaystyle \frac{1}{E}}{\displaystyle \frac{\mathrm{d}^3p^{}}{(2\pi )^3}\frac{\left(1\mathrm{\Theta }(\mu E_p^{})\right)}{2E_p^{}}\frac{\mathrm{d}^3k}{(2\pi )^3}\frac{\mathrm{\Theta }(\mu E_k)}{2E_k}\frac{\mathrm{d}^3k^{}}{(2\pi )^3}\frac{\left(1\mathrm{\Theta }(\mu E_k^{})\right)}{2E_k^{}}}`$ (B1)
$`\times `$ $`(2\pi )^4\delta ^{(4)}(P+KP^{}K^{})||^2,`$ (B2)
where $`||^2`$ is the scattering matrix element squared. After performing the $`p^{}`$ integral
$`\mathrm{\Gamma }(E)`$ $`=`$ $`{\displaystyle \frac{2\pi }{E}}{\displaystyle \frac{\mathrm{d}^3k}{(2\pi )^3}\frac{\mathrm{\Theta }(\mu E_k)}{2E_k}\frac{\mathrm{d}^3k^{}}{(2\pi )^3}\frac{\left(1\mathrm{\Theta }(\mu E_k^{})\right)}{2E_k^{}}\frac{\left(1\mathrm{\Theta }(\mu E_{𝐩+𝐤𝐤^{}})\right)}{2E_{𝐩+𝐤𝐤^{}}}}`$ (B3)
$`\times `$ $`\delta (E+E_kE_{𝐩+𝐤𝐤^{}}E_k^{})||^2.`$ (B4)
We now make the change of variables
$$E_k=\mu t_k,E_k^{}=\mu +t_k^{},E_{𝐩+𝐤𝐤^{}}=\mu +t_{𝐩+𝐤𝐤^{}},$$
(B5)
where the $`t_i`$ variables are positive quantities. The delta function of energy conservation imposes
$$E\mu =t_k+t_k^{}+t_{𝐩+𝐤𝐤^{}},$$
(B6)
which is only valid for $`E\mu 0`$. The maximum value that each one of the variables $`t_i`$ can achieve is $`E\mu `$, while the minimum is zero. Using the energies of the particles as integration variables, we see that the integration is always performed over an energy shell of thickness $`E\mu `$. If $`E\mu \mu `$, then the values of the energy variables inside the integral can be substituted by the Fermi energy. One then finally reaches
$$\mathrm{\Gamma }(E)_{\mu (E\mu )}^\mu 𝑑E_k_\mu ^{\mu +(E\mu )}𝑑E_k^{},$$
(B7)
and thus $`\mathrm{\Gamma }(E)(E\mu )^2`$. A similar argument can be applied for a quasihole to get the energy dependence of its damping rate.
In the case we studied in this article, the electric interactions can be considered as short-ranged, because of Debye screening; they thus give a contribution to the damping rate of electrons as expected from Luttinger’s theorem. The above arguments fail in the case of magnetic interactions, as those are not short-ranged, but rather suffer a weak dynamical screening, where the energies themselves play the role of infrared cutoffs in the above integrals.
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# 1 Introduction
## 1 Introduction
Recently there has been a growing interest concerning supersymmetric models with non–universal soft–breaking terms . The main theoretical reason is that some string-inspired models naturally favour SUSY models with non–universality in the soft-breaking sector . Within this class of string-inspired models particularly interesting are that ones with non–universal $`A`$–terms . This is mainly due to the relevant role they can play in solving the SUSY CP problem, especially in the light of the recent experimental results on the direct CP violation parameter $`\epsilon ^{}/\epsilon `$ .
The new measurements of $`\epsilon ^{}/\epsilon `$ at KTeV and NA48 lead to a world average of Re $`\epsilon ^{}/\epsilon `$ =$`(21.4\pm 4.0)\times 10^4`$ . This result is higher than the Standard Model (SM) predictions , opening the way to the interpretation that it may be a signal of new physics beyond the SM. Clearly it is still premature to claim that this is a genuine new physics effect, mainly due to the large theoretical uncertainties in the non-perturbative hadronic sector which affect the SM predictions for $`\epsilon ^{}/\epsilon `$ . However if one accepts the point of view that new CP sources are needed in order to obtain large values for $`\epsilon ^{}/\epsilon `$ , one may wonder if the minimal supersymmetric extension of the SM (MSSM) can help in solving this problem. The answer is no, mainly due to the assumption of universal boundary conditions of the soft-breaking terms . Moreover we stress that without new flavor structure beyond the usual Yukawa couplings, general SUSY models with phases of the soft terms of order $`𝒪(1)`$ (but with a vanishing CKM phase $`\delta _{\mathrm{CKM}}=0`$) can not give a sizeable contribution to the CP violating processes .
The impact of the new flavor structure in the non–universality of the $`A`$–terms have been studied in Refs. . In these works it was emphasized that the non-degenerate $`A`$–terms can generate the experimentally observed CP violation $`\epsilon `$ and $`\epsilon ^{}/\epsilon `$ even with a vanishing $`\delta _{CKM}`$, i.e., fully supersymmetric CP violation in the kaon system is possible in a class of models with non–universal $`A`$–terms. This effect can be simply understood by making use of the mass-insertion approximation : the non-degenerate $`A`$–terms enhance the gluino contributions to $`\epsilon `$ and $`\epsilon ^{}/\epsilon `$ through enhancing the imaginary parts of the L-R mass insertions $`\mathrm{Im}(\delta _{LR}^d)_{12}`$ and $`\mathrm{Im}(\delta _{RL}^d)_{12}`$ .
It is well known that the experimental $`bs\gamma `$ constraints cause a dramatic reduction of the allowed parameter space in case of universal soft terms . Hence one may worry if these constraints would be even more severe in the case of the non-degenerate $`A`$–terms. However a complete analysis of the $`bs\gamma `$ constraints has not been considered in Refs. . In particular in these works it was roughly checked, by using the results of Ref. , that the $`bs\gamma `$ constraints are satisfied, namely $`(\delta _{LR}^d)_{23}1.6\times 10^2`$ and $`(\delta _{LL}^d)_{23}8.2`$. Note that these constraints are obtained by assuming that the gluino amplitude is the dominant contribution to $`bs\gamma `$ (dominant also with respect to the SM).
Although the gluino contribution to $`bs\gamma `$ decay is usually very small in the universal case, being proportional to the mass insertions $`(\delta _{LR}^d)_{23}`$, we could expect that this contribution may be enhanced by the non–universality of the $`A`$–terms. However the non–universal $`A`$–terms could also give large contributions to the chargino amplitude through the $`(\delta _{LR}^u)_{23}`$. Even though the non-degenerate $`A`$–terms enhance the gluino contribution, one can not a priori expect that it will be the dominant effect. Therefore a careful analysis of the $`bs\gamma `$ predictions, including the full SUSY contributions, is necessary in this scenario.
The purpose of this work is to perform a complete analysis of the $`bs\gamma `$ constraints for the SUSY models with non–universal $`A`$–terms studied in Refs. . Indeed we will show that our results are different from the naive expectations based only on the gluino-dominance approximation.
The paper is organized as follows. In section 2 we present two models with non–universal $`A`$–terms that have recently been considered for solving the SUSY CP problem. These two models are based on weakly coupled heterotic string and type I string theories respectively. In section 3 we present formulas for the total branching ratio taking into account the QCD corrections and we discuss the different SUSY contributions to the $`bs\gamma `$ decay in these models. Our numerical results for the predictions of the $`bs\gamma `$ branching ratio are presented in section 4. The last section is devoted to conclusions. Finally, various formulas are summarized in the appendix.
## 2 String inspired models with non-degenerate <br>$`A`$–terms
In this work we consider the class of string inspired model which has been recently studied in Refs. . In this class of models, the trilinear $`A`$–terms of the soft SUSY breaking are non–universal. It was shown that this non–universality among the $`A`$–terms plays an important role on CP violating processes. In particular, it has been shown that non-degenerate $`A`$-parameters can generate the experimentally observed CP violation $`\epsilon `$ and $`\epsilon ^{}/\epsilon `$ even with a vanishing $`\delta _{\mathrm{CKM}}`$.
Here we consider two models for non-degenerate $`A`$–terms. The first model (model A) is based on weakly coupled heterotic strings, where the dilaton and the moduli fields contribute to SUSY breaking . The second model (model B) is based on type I string theory where the gauge group $`SU(3)\times U(1)_Y`$ is originated from the $`9`$ brane and the gauge group $`SU(2)`$ is originated from one of the $`5`$ branes .
In order to fix the conventions, the following Lagrangian $`_{SB}`$ for the soft-breaking terms is assumed
$$_{SB}=\frac{1}{6}h_{ijk}\varphi _i\varphi _j\varphi _k+\frac{1}{2}(\mu B)^{ij}\varphi _i\varphi _j+\frac{1}{2}(m^2)_i^j\varphi ^i\varphi _j+\frac{1}{2}M_a\lambda _a\lambda _a+h.c.$$
(1)
where the $`\varphi _i`$ are the scalar parts of the chiral superfields $`\mathrm{\Phi }_i`$ and $`\lambda _a`$ are the gaugino fields. In the notation for the trilinear couplings, the $`A`$–terms are defined as $`h_{ijk}=Y_{ijk}A_{ijk}`$ (indices not summed) where $`Y_{ijk}`$ are the corresponding Yukawa couplings.
### 2.1 Model A
We start with the weakly coupled string-inspired supergravity theory. In this class of models, it is assumed that the superpotential of the dilaton ($`S`$) and moduli ($`T`$) fields is generated by some non-perturbative mechanism and the $`F`$-terms of $`S`$ and $`T`$ contribute to the SUSY breaking. Then one can parametrize the $`F`$-terms as
$$F^S=\sqrt{3}m_{3/2}(S+S^{})\mathrm{sin}\theta ,F^T=m_{3/2}(T+T^{})\mathrm{cos}\theta .$$
(2)
Here $`m_{3/2}`$ is the gravitino mass, $`n_i`$ is the modular weight and $`\mathrm{tan}\theta `$ corresponds to the ratio between the $`F`$-terms of $`S`$ and $`T`$. In this framework, the soft scalar masses $`m_i`$ and the gaugino masses $`M_a`$ are given by
$`m_i^2`$ $`=`$ $`m_{3/2}^2(1+n_i\mathrm{cos}^2\theta ),`$ (3)
$`M_a`$ $`=`$ $`\sqrt{3}m_{3/2}\mathrm{sin}\theta .`$ (4)
The $`A^{u,d}`$-terms are written as
$`(A^{u,d})_{ij}`$ $`=`$ $`\sqrt{3}m_{3/2}\mathrm{sin}\theta m_{3/2}\mathrm{cos}\theta (3+n_i+n_j+n_{H_{u,d}}),`$ (5)
where $`n_{i,j,k}`$ are the modular weights of the fields that are coupled by this $`A`$–term. As shown in Eqs.(3-5), the values of the soft SUSY breaking parameters depend on the modular weight of the matter fields. These modular weights $`n_i`$ are negative integers, and their ‘natural’ values (in case of $`Z_N`$ orbifolds) are $`1,2`$, and $`3`$ . If we assign $`n_i=1`$ for the third family and $`n_i=2`$ for the first and second families (we also assume that $`n_{H_1}=1`$ and $`n_{H_2}=2`$). Note that with this choice of modular weights we have $`m_{H_2}^2<m_{H_1}^2`$ which is favored for the electroweak breaking (EW) and all the squark mass matrices are equal. Also we find the following texture for the $`A`$-parameter matrix at the string scale
$$A^{u,d}=\left(\begin{array}{ccc}x_{u,d}& x_{u,d}& y_{u,d}\\ x_{u,d}& x_{u,d}& y_{u,d}\\ y_{u,d}& y_{u,d}& z_{u,d}\end{array}\right),$$
(6)
where
$`x_u`$ $`=`$ $`m_{3/2}(\sqrt{3}\mathrm{sin}\theta +3\mathrm{cos}\theta ),`$ (7)
$`x_d`$ $`=`$ $`y_u=m_{3/2}(\sqrt{3}\mathrm{sin}\theta +2\mathrm{cos}\theta ),`$ (8)
$`y_d`$ $`=`$ $`z_u=m_{3/2}(\sqrt{3}\mathrm{sin}\theta +\mathrm{cos}\theta ),`$ (9)
$`z_d`$ $`=`$ $`\sqrt{3}m_{3/2}\mathrm{sin}\theta .`$ (10)
By fixing the value of $`\mathrm{tan}\beta `$ we can determine the values of $`\mu `$ and $`B`$ from the radiatively EW breaking conditions. Then all the SUSY particle spectrum is completely determined in terms of $`m_{3/2}`$ and $`\theta `$. The non–universality of this model is parameterized by the angle $`\theta `$ and the value $`\theta =\pi /2`$ corresponds to the universal limit for the soft terms. In order to avoid negative mass squared in the scalar masses we restrict ourselves to the case with $`\mathrm{cos}^2\theta <1/2`$. Such restriction on $`\theta `$ makes the non–universality in the whole soft SUSY breaking terms very limited. However, as shown in , this small range of variation for the non–universality is enough to generate sizeable SUSY CP violations in K system. We emphasize that choosing modular weights different from those assigned above, the allowed range of the soft SUSY breaking terms do not essentially change. For instance, if we assign $`n_i=3`$ for the first family instead of $`2`$, it may appear that the non–universality among the entries of the $`A`$–terms is enhanced. However in this case the angle $`\theta `$ is more constrained than before ($`\mathrm{cos}^2\theta <1/3`$). We have checked that different choices for $`n_i`$ do not significantly affect our results for the $`bs\gamma `$ branching ratio.
### 2.2 Model B
As mentioned in the introduction, this model is based on type I string theory. Like model A, this is a good candidate for generating sizeable SUSY CP violations. Recently, there has been considerable interest in studying the phenomenological implications of this class of models . In type I string theory, non universality in the scalar masses, $`A`$–terms and gaugino masses can be naturally obtained . Type I models contain either 9 branes and three types of $`5_i(i=1,2,3)`$ branes or $`7_i`$ branes and 3 branes. From the phenomenological point of view there is no difference between these two scenarios. Here we consider the same model used in Ref. , where the gauge group $`SU(3)_C\times U(1)_Y`$ is associated with 9 brane while $`SU(2)_L`$ is associated with $`5_1`$ brane.
If SUSY breaking is analysed, as in model A, in terms of the vevs of the dilaton and moduli fields
$$F^S=\sqrt{3}m_{3/2}(S+S^{})\mathrm{sin}\theta ,F^{T_i}=m_{3/2}(T_i+T_i^{})\mathrm{\Theta }_i\mathrm{cos}\theta ,$$
(11)
where the angle $`\theta `$ and the parameter $`\mathrm{\Theta }_i`$ with $`_i\left|\mathrm{\Theta }_i\right|^2=1`$, just parametrize the direction of the goldstino in the $`S`$ and $`T_i`$ fields space . Within this framework, the gaugino masses are
$`M_1`$ $`=`$ $`M_3=\sqrt{3}m_{3/2}\mathrm{sin}\theta ,`$ (12)
$`M_2`$ $`=`$ $`\sqrt{3}m_{3/2}\mathrm{\Theta }_1\mathrm{cos}\theta .`$ (13)
In this case the quark doublets and the Higgs fields are assigned to the open string which spans between the $`5_1`$ and $`9`$ branes. While the quark singlets correspond to the open string which starts and ends on the $`9`$ brane, such open string includes three sectors which correspond to the three complex compact dimensions. If we assign the quark singlets to different sectors we obtain non–universal $`A`$–terms. It turns out that in this model the trilinear couplings $`A^u`$ and $`A^d`$ are given by
$$A^u=A^d=\left(\begin{array}{ccc}x& y& z\\ x& y& z\\ x& y& z\end{array}\right),$$
(14)
where
$`x`$ $`=`$ $`\sqrt{3}m_{3/2}\left(\mathrm{sin}\theta +(\mathrm{\Theta }_1\mathrm{\Theta }_3)\mathrm{cos}\theta \right),`$ (15)
$`y`$ $`=`$ $`\sqrt{3}m_{3/2}\left(\mathrm{sin}\theta +(\mathrm{\Theta }_1\mathrm{\Theta }_2)\mathrm{cos}\theta \right),`$ (16)
$`z`$ $`=`$ $`\sqrt{3}m_{3/2}\mathrm{sin}\theta .`$ (17)
The soft scalar masses for quark-doublets and Higgs fields $`(m_L^2)`$, and the quark-singlets $`(m_{R_i}^2)`$ are given by
$`m_L^2`$ $`=`$ $`m_{3/2}^2\left(1{\displaystyle \frac{3}{2}}(1\mathrm{\Theta }_1^2)\mathrm{cos}^2\theta \right),`$ (18)
$`m_{R_i}^2`$ $`=`$ $`m_{3/2}^2\left(13\mathrm{\Theta }_i^2\mathrm{cos}^2\theta \right),`$ (19)
where $`i`$ refers to the three families. For $`\mathrm{\Theta }_i=1/\sqrt{3}`$ the $`A`$–terms and the scalar masses are universal while the gaugino masses could be non–universal. The universal gaugino masses are obtained at $`\theta =\pi /6`$.
It is worth mentioning that in these models (A and B) the gaugino masses, the $`A`$–terms, and the $`\mu `$–term are in general complex. However, by using $`R`$-rotation we can make the gaugino masses real and we end up, in addition to the phase of $`\mu `$, with the phases of the $`A`$–terms. The phase of $`\mu `$ is severely constrained by the electric dipole moment (EDM) of the electron and the neutron , while the phases of the $`A`$–terms are essentially unconstrained. Thus one can set the phase of $`\mu `$ to zero, as done in Ref. . Moreover it has been shown that the phases of the $`A`$–terms can lead to sizeable supersymmetric contribution to CP observables, in particular on the direct CP violation $`\epsilon ^{}/\epsilon `$. However, since the total branching ratio $`bs\gamma `$ decay is a CP conserving observable, this should not be very effective in constraining the phases of the SUSY soft-breaking terms. For this reason we have made in our analysis the simplifying assumption to set to zero all the phases.
### 2.3 Yukawa textures
As emphasized in Refs. , in models with non-degenerate $`A`$–terms we have to fix the Yukawa matrices to completely specify the model. In fact, with universal $`A`$–terms the textures of the Yukawa matrices at GUT scale affect the physics at EW scale only through the quark masses and usual CKM matrix, since the extra parameters contained in the Yukawa matrices can be eliminated by unitary fields transformations. This is no longer true with non-degenerate $`A`$–terms since in the scalar potential the $`A`$–terms enter through the tensorial product $`(Y_q^A)_{ij}=(Y_q)_{ij}(A_q)_{ij}`$. Thus the diagonalization of $`Y_q`$ can not be done simultaneously with $`Y_q^A`$ (unlike the universal case). Thus in the models with non–universal $`A`$–terms, some extra degrees of freedom (in addition to the quark masses and CKM matrix) contained in the Yukawa matrices become observable. Hence, the analysis of the non-degenerate $`A`$–terms could shed some light on the favoured Yukawa textures. For instance in Ref. , using a symmetric Yukawa texture with a symmetric $`A`$–terms, an accidental cancellation between the different SUSY contributions to $`\epsilon ^{}/\epsilon `$ was found, cancellation which leads to a very small value for $`\epsilon ^{}/\epsilon `$. On the contrary, by using asymmetric Yukawa matrices with symmetric $`A`$–terms, this cancellation does not occur and it is found that $`\epsilon ^{}/\epsilon `$ can be easily of the order of the KTeV result.
Here we show two realistic examples of Yukawa matrix textures that have already been used in Refs. . In the first one we have the following symmetric Yukawa matrices
$`\begin{array}{cc}Y^d=y^b\left(\begin{array}{ccc}0& V_{12}\frac{m_s}{m_b}& V_{13}\\ V_{12}\frac{m_s}{m_b}& \frac{m_s}{m_b}& V_{23}\\ V_{13}& V_{23}& 1\end{array}\right),\hfill & \hfill Y^u=y^t\left(\begin{array}{ccc}0& 0& V_{13}\\ 0& \frac{m_c}{m_t}& 0\\ V_{13}& 0& 1\end{array}\right),\end{array}`$ (27)
where $`y^{b,t}`$ are the Yukawa couplings of the bottom and top respectively, and $`V`$ is the CKM matrix. The second example is based on the assumption that the CKM mixing matrix originates from the down Yukawa couplings and that the Yukawa matrices are hermitian.
$$Y^u=\frac{1}{v\mathrm{cos}\beta }\mathrm{diag}(m_u,m_c,m_t),Y^d=\frac{1}{v\mathrm{sin}\beta }V^{}\mathrm{diag}(m_d,m_s,m_b)V$$
(28)
Although the analysis of the CP violation is quite sensitive to the specific Yukawa matrix, we found that the branching ratio of $`bs\gamma `$ does not essentially depend on it. We checked this property by using different Yukawa textures (the two examples presented here and others). We will only present through all the paper the results concerning the second example just as a representative case.
Now we present the general expressions for the squark mass matrices in the non-universal case in the SCKM basis, in this basis the unitary matrices $`S^{U_{R,L}}`$ and $`S^{D_{R,L}}`$ (obtained by a superfield rotation) are chosen to diagonalize the up– and down– Yukawa couplings $`Y^{u,d}`$
$$m_U=\frac{v\mathrm{sin}\beta }{\sqrt{2}}S_{U_R}(Y^{u,d})^TS_{U_L}^{},m_D=\frac{v\mathrm{cos}\beta }{\sqrt{2}}S_{D_R}(Y^{u,d})^TS_{D_L}^{}$$
(29)
where $`T`$ stands for the transpose and $`m_{U,D}`$ are the diagonal up– and down–quark mass matrices respectively. In this basis the up and down squark mass matrices at low energy (respectively $`M_{\stackrel{~}{u}}^2`$ and $`M_{\stackrel{~}{d}}^2`$) are given by
$`\begin{array}{cc}M_{\stackrel{~}{u},\stackrel{~}{d}}^2=\left(\begin{array}{cc}\left(M_{\stackrel{~}{u},\stackrel{~}{d}}^2\right)_{LL}& \left(M_{\stackrel{~}{u},\stackrel{~}{d}}^2\right)_{LR}\\ \left(M_{\stackrel{~}{u},\stackrel{~}{d}}^2\right)_{RL}& \left(M_{\stackrel{~}{u},\stackrel{~}{d}}^2\right)_{RR}\end{array}\right),\hfill & \end{array}`$ (33)
where for the up-sector
$`\left(M_{\stackrel{~}{u}}^2\right)_{LL}`$ $`=`$ $`S_{U_L}M_{\stackrel{~}{Q}}^2S_{U_L}^{}+m_U^2+{\displaystyle \frac{m_Z^2}{6}}(34\mathrm{sin}^2\theta _W)\mathrm{cos}2\beta ,`$
$`\left(M_{\stackrel{~}{u}}^2\right)_{RR}`$ $`=`$ $`S_{U_R}(M_{\stackrel{~}{u}^c}^2)^TS_{U_R}^{}+m_U^2+{\displaystyle \frac{2m_Z^2}{3}}\mathrm{sin}^2\theta _W\mathrm{cos}2\beta ,`$
$`\left(M_{\stackrel{~}{u}}^2\right)_{LR}`$ $`=`$ $`\left(M_{\stackrel{~}{u}}^2\right)_{RL}^{}=\mu m_U\mathrm{cot}\beta +{\displaystyle \frac{v\mathrm{sin}\beta }{\sqrt{2}}}S_{U_L}Y_u^AS_{U_R}^{},`$ (34)
and for the down-sector
$`\left(M_{\stackrel{~}{d}}^2\right)_{LL}`$ $`=`$ $`S_{D_L}M_{\stackrel{~}{Q}}^2S_{D_L}^{}+m_D^2{\displaystyle \frac{m_Z^2}{6}}(32\mathrm{sin}^2\theta _W)\mathrm{cos}2\beta ,`$
$`\left(M_{\stackrel{~}{d}}^2\right)_{RR}`$ $`=`$ $`S_{D_R}(M_{\stackrel{~}{d}^c}^2)^TS_{D_R}^{}+m_D^2+{\displaystyle \frac{2m_Z^2}{3}}\mathrm{sin}^2\theta _W\mathrm{cos}2\beta ,`$
$`\left(M_{\stackrel{~}{d}}^2\right)_{LR}`$ $`=`$ $`\left(M_{\stackrel{~}{d}}^2\right)_{RL}^{}=\mu m_D\mathrm{tan}\beta +{\displaystyle \frac{v\mathrm{cos}\beta }{\sqrt{2}}}S_{D_L}Y_d^AS_{D_R}^{}`$ (35)
where $`M_{\stackrel{~}{Q}}^2`$ and $`M_{\stackrel{~}{u}^c,\stackrel{~}{d}^c}^2`$ are the soft-breaking ($`3\times 3`$) mass matrices for the squark doublet and singlets respectively. Our convention for the sign of the $`\mu `$–term is opposite to the same one of Ref. . Note that the matrices $`S_{U,D}`$, unlike in the universal case, can not be re–absorbed in the definition of diagonal Yukawa couplings.
## 3 The $`bs\gamma `$ decay in SUSY models
In this section we analyse the $`bs\gamma `$ decay in SUSY models with non–universal $`A`$–terms. As pointed out previously, these models are particularly interesting because they can give sizeable contributions to the CP violating processes through their large contributions to the mass insertions $`(\delta _{LR}^{u,d})_{ij}`$ or $`(\delta _{RL}^{u,d})_{ij}`$. For this reason one may expect that large effects can be also induced in the processes mediated by the dipole-magnetic operators, such as the rare decay $`bs\gamma `$ . Indeed, if the $`\delta _{LR}^{u,d}`$ are large enough, then the SUSY contributions to these operators are enhanced since the typical chiral suppression is removed by the insertion of the internal gaugino mass, thus allowing for a competition with the chiral-suppressed SM amplitude.
Let us start with the experimental results. The most recent result reported by CLEO collaboration for the total (inclusive) B meson branching ratio $`BX_s\gamma `$ is
$$\mathrm{BR}(BX_s\gamma )=(3.15\pm 0.35\pm 0.32\pm 0.26)\times 10^4$$
(36)
where the first error is statistical, the second systematic, and the third one accounts for model dependence. From this result the following bounds (each of them at 95% C.L.) are obtained
$$2.0\times 10^4<\mathrm{BR}(BX_s\gamma )<4.5\times 10^4.$$
(37)
In addition the ALEPH collaboration at LEP reported a compatible measurement of the corresponding branching ratio for b hadrons at the Z resonance .
The starting point for the theoretical study of $`bs\gamma `$ decay is given by the effective Hamiltonian
$$H_{eff}=\frac{4G_F}{\sqrt{2}}V_{32}^{}V_{33}\underset{i=1}{\overset{8}{}}C_i(\mu _b)Q_i(\mu _b)$$
(38)
where the complete basis of operators in the SM can be found in Ref. . Recently the main theoretical uncertainties present in the previous leading order (LO) SM calculations have been reduced by including the NLO corrections to the $`bs\gamma `$ decay, through the calculation of the three-loop anomalous dimension matrix of the effective theory . The relevant SUSY contributions to the effective Hamiltonian in Eq.(38) affect only the $`Q_7`$ and $`Q_8`$ operators, the expression for these operators are given (in the usual notation) by
$`Q_7`$ $`=`$ $`{\displaystyle \frac{e}{16\pi ^2}}m_b\left(\overline{s}_L\sigma ^{\mu \nu }b_R\right)F_{\mu \nu },`$
$`Q_8`$ $`=`$ $`{\displaystyle \frac{g_s}{16\pi ^2}}m_b\left(\overline{s}_L\sigma ^{\mu \nu }T^ab_R\right)G_{\mu \nu }^a.`$ (39)
The Wilson coefficients $`C_i(\mu )`$ are evaluated at the renormalization scale $`\mu _bO(m_b)`$ by including the NLO corrections . They can be formally decomposed as follows
$$C_i(\mu )=C_i^{(0)}(\mu )+\frac{\alpha _s(\mu )}{4\pi }C_i^{(1)}(\mu )+𝒪(\alpha _s^2).$$
(40)
where $`C_i^{(0)}`$ and $`C_i^{(1)}`$ stand for the LO and NLO order respectively. Finally the branching ratio $`\mathrm{BR}(BX_s\gamma )`$, conventionally normalized to the semileptonic branching ratio $`\mathrm{BR}^{exp}(BX_ce\nu )=(10.4\pm 0.4)\%`$ , is given by
$`\mathrm{BR}^{\mathrm{NLO}}(BX_s\gamma )`$ $`=`$ $`\mathrm{BR}^{exp}(BX_ce\nu ){\displaystyle \frac{|V_{32}^{}V_{33}|^2}{|V_{23}|^2}}{\displaystyle \frac{6\alpha _e}{\pi g(z)k(z)}}\left(1{\displaystyle \frac{8}{3}}{\displaystyle \frac{\alpha _s(m_b)}{\pi }}\right)`$ (41)
$`\times `$ $`\left(|D|^2+A\right)(1+\delta _{np}),`$
with
$`D`$ $`=`$ $`C_7^{(0)}(\mu _b)+{\displaystyle \frac{\alpha _s(\mu _b)}{4\pi }}\left(C_7^{(1)}(\mu _b)+{\displaystyle \underset{i=1}{\overset{8}{}}}C_i^{(0)}(\mu _b)\left[r_i(z)+\gamma _{i7}^{(0)}\mathrm{log}{\displaystyle \frac{m_b}{\mu _b}}\right]\right),`$
$`A`$ $`=`$ $`\left(e^{\alpha _s(\mu _b)\mathrm{log}\delta (7+2\mathrm{log}\delta )/3\pi }1\right)|C_7^{(0)}(\mu _b)|^2+{\displaystyle \frac{\alpha _s(\mu _b)}{\pi }}{\displaystyle \underset{ij=1}{\overset{8}{}}}C_i^{(0)}(\mu _b)C_j^{(0)}(\mu _b)f_{ij}(\delta ),`$
where $`z=m_c^2/m_b^2`$. The expressions for $`C_i^{(0)}`$, $`C_i^{(1)}`$, and the anomalous dimension matrix $`\gamma `$, together with the functions $`g(z)`$, $`k(z)`$, $`r_i(z)`$ and $`f_{ij}(\delta )`$, can be found in Ref. . The term $`\delta _{np}`$ (of order a few percent) includes the non-perturbative $`1/m_b`$ and $`1/m_c`$ corrections. From the formula above we obtain the theoretical result for BR($`BX_s\gamma `$) in the SM which is given by
$$\mathrm{BR}^{\mathrm{NLO}}(BX_s\gamma )=(3.29\pm 0.33)\times 10^4$$
(42)
where the main theoretical uncertainty comes from uncertainties in the SM input parameters, namely $`m_t,\alpha _s(M_Z),\alpha _{em},m_c/m_b,m_b,V_{ij}`$, and the small residual scale dependence.<sup>§</sup><sup>§</sup>§ Recently in Ref. the current method of extracting the inclusive rate for $`bs\gamma `$ from the currently published CLEO data has been criticized arguing that the theoretical uncertainties have so far been underestimated and only a precise measurement of the photon spectrum would be help in reducing these uncertainties. The central value in Eq.(42) corresponds to the following central values for the SM parameters $`m_t^{\mathrm{pole}}m_t^{\overline{\mathrm{MS}}}(m_Z)174\mathrm{GeV}`$, $`m_b^{\mathrm{pole}}=4.8\mathrm{GeV}`$, $`m_c^{\mathrm{pole}}=1.3\mathrm{GeV}`$, $`\mu _b=m_b`$, $`\alpha _s(m_Z)=0.118`$, $`\alpha _e^1(m_Z)=128`$, $`\mathrm{sin}^2\theta _W=0.23`$ and a photon energy resolution corresponding to $`\delta =0.9`$ is assumed. Note that in Eq.(41) the (small) $`1/m_c`$ corrections have not been included.
The SUSY contributions to the Wilson coefficients $`C_{7,8}^{(0,1)}`$ are obtained by calculating the $`bs\gamma `$ and $`bsg`$ amplitudes at EW scale respectively. The LO contributions to these amplitudes are given by the 1-loop magnetic-dipole and chromomagnetic dipole penguin diagrams respectively, mediated by charged Higgs boson, chargino, gluino, and neutralino exchanges. The corresponding results for these amplitudes can be found in Ref.. It is known that the charged Higgs contribution always interferes with the SM contribution . The chargino contribution could give rise to a substantial destructive interference with SM and charged Higgs amplitudes, depending on the sign of $`\mu `$, the value of $`\mathrm{tan}\beta `$, and the mass difference between the stop masses .
We point out that the SUSY models with non–universal $`A`$–terms may induce non-negligible contributions to the dipole operators $`\stackrel{~}{Q}_{7,8}`$ which have opposite chirality with respect to $`Q_{7,8}`$. It is worth mentioning that these operators are also induced in the SM and in the MSSM with supergravity scenario, but their contributions are negligible being suppressed by terms of order $`𝒪(m_s/m_b)`$. In particular in MSSM, due to the universality of the $`A`$–terms, the gluino and chargino contributions to $`\stackrel{~}{Q}_{7,8}`$ turn out to be of order $`𝒪(m_s/m_b)`$. This argument does not hold in the models with non–universal $`A`$–terms and in particular in our case. It can be simply understood by using the mass insertion method . For instance, the gluino contributions to $`Q_7`$ and $`\stackrel{~}{Q}_7`$ operators are proportional to $`(\delta _{LR}^d)_{23}(S_{D_L}Y_d^AS_{D_R}^{})_{23}/m_{\stackrel{~}{q}}^2`$ and $`(\delta _{RL}^d)_{23}(S_{D_R}Y_d^AS_{D_L}^{})_{23}/m_{\stackrel{~}{q}}^2`$ respectively. Since the $`A^D`$ matrix is symmetric in model A and $`A_{ij}^DA_{ji}^D`$ in model B, then $`(\delta _{LR}^d)_{23}(\delta _{RL}^d)_{23}`$. Then in our case we should consistently take into account the SUSY contributions to $`\stackrel{~}{Q}_7`$ in $`bs\gamma `$ . Analogous considerations hold for the operator $`\stackrel{~}{Q}_8`$.
By taking into account the above considerations regarding the operators $`\stackrel{~}{Q}_{7,8}`$, the new physics effects in $`bs\gamma `$ can be parametrized in a model independent way by introducing the so called $`R_{7,8}`$ and $`\stackrel{~}{R}_{7,8}`$ parameters defined at EW scale as
$$R_{7,8}=\frac{\left(C_{7,8}^{(0)}C_{7,8}^{(0)SM}\right)}{C_{7,8}^{(0)SM}},\stackrel{~}{R}_{7,8}=\frac{\stackrel{~}{C}_{7,8}^{(0)}}{C_{7,8}^{(0)SM}},$$
(43)
where $`C_{7,8}`$ include the total contribution while $`C_{7,8}^{SM}`$ contains only the SM ones. Note that in $`\stackrel{~}{C}_{7,8}`$, which are the corresponding Wilson coefficients for $`\stackrel{~}{Q}_{7,8}`$ respectively, we have set to zero the SM contribution. In Ref. only the expressions for the $`R_{7,8}`$ are given, for completeness we report the corresponding expressions for $`\stackrel{~}{R}_{7,8}`$ in the appendix.
Inserting these definitions into the BR($`BX_s\gamma `$) formula in Eq.(41) yields a general parametrization of the branching ratio in terms of the new physics contributions Note that the SM central value for BR($`BX_s\gamma `$) in Ref. slightly differs from the result in Eq.(42), since in Ref. the non-perturbative $`\mathrm{\Lambda }/m_c`$ corrections have been included.
$`\mathrm{BR}(BX_s\gamma )`$ $`=`$ $`(3.29\pm 0.33)\times 10^4(1+0.622R_7+0.090(R_7^2+\stackrel{~}{R}_7^2)`$ (44)
$`+`$ $`0.066R_8+0.019(R_7R_8+\stackrel{~}{R}_7\stackrel{~}{R}_8)+0.002(R_8^2+\stackrel{~}{R}_8^2)),`$
where the overall SM uncertainty has been factorized outside. We have checked explicitly that the result in Eq.(44) is in agreement with the corresponding one used in Ref. . Recently the leading EW corrections in the SM have been included in the $`C_{7,8}`$ coefficients . However, we have not included them since they should affect the results in a few percent which is within the theoretical uncertainties present in the SUSY sector. In particular, in order to be consistent with the NLO calculations, one should also include the corresponding two-loop QCD corrections to the SUSY amplitudes, namely $`C_{7,8}^{(1)}`$. In this respect recent works have been done in this direction. They calculated the two-loop $`\alpha _s`$ corrections to the chargino and gluino amplitudes. However these corrections can be consistently applied only in a restricted SUSY scenario with low $`\mathrm{tan}\beta `$and large gluino masses, since on the contrary other two-loop diagrams with no QCD interactions at all could become relevant. In our work we are going to explore SUSY scenarios in the full range of $`\mathrm{tan}\beta `$($`2\mathrm{tan}\beta m_t/m_b`$) and gluino mass, so we do not include these corrections since they would not change our main conclusions.
## 4 Numerical results and discussions
In this section we present our results for the total branching ratio BR($`BX_s\gamma `$) as a function of the fundamental parameters of the soft-breaking sector of models A and B. We start our discussion by analysing the constraints on these models set by the condition of vacuum stability and the experimental bounds on the SUSY particle spectrum.
We calculate the low energy SUSY spectrum by running the soft-breaking terms and the other SUSY parameters (by means of the renormalization group equations (RGE) of MSSM generalized to the non–universal soft-breaking terms ) from the GUT scale $`(2\times 10^{16}`$ GeV) to the EW scale ($`M_Z`$), by using the boundary conditions given in section 2. For fixed $`\mathrm{tan}\beta `$and sign of $`\mu `$, we restrict the parameter space by imposing the present experimental bounds on the SUSY spectra . Since the spectrum of these models is given in terms of few parameters ($`m_{3/2},\theta `$ and/or $`\mathrm{\Theta }_i`$) we find that by requiring the lightest chargino mass to be $`m_{\chi ^\pm }`$ $`90\mathrm{GeV}`$ implies that all the other SUSY particle and the lightest Higgs masses are above their experimental bounds.
In order to avoid vacuum instabilities and color-charge breaking, we require that all the square scalar masses in Eqs.(3), (19) should be positive. All these constraints set strong restrictions on the parameter space of both models. As pointed out in section 2, in model A this leads to $`\mathrm{sin}\theta >1/\sqrt{2}0.7`$. The bounds on $`m_{3/2}`$ from the lightest chargino mass are as follows: for low $`\mathrm{tan}\beta `$($`\mathrm{tan}\beta `$$`=2`$) $`m_{3/2}80(60)\mathrm{GeV}`$ for $`\mathrm{sin}\theta 1/\sqrt{2}(1)`$. For large $`\mathrm{tan}\beta `$these bounds are increased by $`20`$ GeV in both cases.
In model B the non–universality is parameterized by the angle $`\theta `$ and the $`\mathrm{\Theta }_i`$’s. The values $`\mathrm{\Theta }_i=1/\sqrt{3}`$ give universal $`A`$ terms and scalar masses, the gaugino masses are universal only at $`\theta =\pi /6`$. To avoid negative scalar masses in this model one needs to impose the constraints
$`\mathrm{cos}^2\theta \mathrm{\Theta }_i`$ $`<`$ $`1/3,(i=1,2,3),`$
$`(1\mathrm{\Theta }_1^2)\mathrm{cos}^2\theta `$ $`<`$ $`2/3.`$ (45)
In the limit of universal $`A`$ terms these two constraints are satisfied for any value of $`\theta `$.
Further constraints on the parameter space of model B are obtained from the gaugino sector. As shown in Eq.(13), for a non very large $`m_{3/2}`$, the limit of $`\theta \pi /2`$ can not be reached since $`M_2`$ would approach zero and the mass of the lightest chargino would be too small. Moreover, the lower bound on the gluino mass and the condition of having the EW breaking at the correct scale, lead also to a lower bound on the angle $`\theta `$. In the case of universal $`A`$–terms we find that $`\theta >0.1`$, while in the non–universal case Eqs.(45) set more severe constraints on $`\theta `$. For example, for $`\mathrm{\Theta }_11`$ we obtain that $`0.9\theta <\pi /2`$ and for $`\theta `$ close to $`\pi /2`$ the gravitino mass $`m_{3/2}`$ should be quite heavy (of order TeV) to make the lightest chargino higher than the experimental bound. We find that these constraints together strongly reduce the allowed parameter space of model B.
We have checked that, in both models, the $`B\overline{B}`$ mixing measurements do not set further constraints on the allowed ranges of the parameter space. This mixing, being a $`\mathrm{\Delta }B=2`$ process, is proportional to $`(\delta _{AB}^d)_{13}\times (\delta _{CD}^d)_{13}`$ where $`A,B,C,D=(L,R)`$. We have found that, in our case, the values of these mass insertions satisfy the constraints given in Ref..
Our results for the partial SUSY amplitude contributions are presented in Figs.\[1-3\]. The total branching ratio BR($`BX_s\gamma `$) is shown in Figs.\[4-5\] and for model A and B respectively. In Figs.\[1-3\] we show, for model A, the individual SUSY contributions to the $`R_7`$ variable, (see Eq.(43) for its definition) versus $`\mathrm{sin}\theta `$. We see that the chargino and charged Higgs contributions give the dominant effect in all the range of $`\theta `$, while the gluino is sub–dominant, such as in the universal case. For model B, as expected from its tightly constrained parameter space, we have found that the results of the separate amplitude contributions do not differ from the corresponding ones in the universal scenario, and will not be presented here.
We can understand this behavior by using the mass insertion method. The gluino amplitude gets two leading contributions: one is proportional to the single–mass–insertion $`(\delta _{LR}^d)_{23}`$ and the another one to the double–mass–insertion, namely $`(\delta _{LR}^d)_{22}(\delta _{LL}^d)_{23}`$. In the low and intermediate $`\mathrm{tan}\beta `$regions these two mass insertions are comparable, so a possible destructive or constructive interference between them may appear depending on the sign of $`\mu `$. In the large $`\mathrm{tan}\beta `$region the double mass insertion becomes dominant, since $`(\delta _{LR}^d)_{22}`$ is proportional to $`\mu \mathrm{tan}\beta `$ and the amplitude (normalized to the SM one) is
$$\frac{A_{\stackrel{~}{g}}}{A_{SM}}\frac{\alpha _s}{\alpha _WV_{32}}\mathrm{tan}\beta (m_W^2M_{\stackrel{~}{g}}\mu )\frac{M_{\stackrel{~}{b}_L\stackrel{~}{s}_L}^2}{m_{\stackrel{~}{q}}^6},$$
(46)
where $`m_{\stackrel{~}{q}}`$ is the average squark mass in the down sector and $`M_{\stackrel{~}{b}_L\stackrel{~}{s}_L}^2`$ is the off diagonal element of the down-squark mass matrix defined in Eqs.(33,35). In this way, the sensitivity of the gluino amplitude to the sign of $`\mu `$ for large $`\mathrm{tan}\beta `$can be understood. The chargino amplitude, like in the gluino case, is enhanced by $`\mathrm{tan}\beta `$and the dominant contributions to it are the (Higgsino) $`A_{\stackrel{~}{h}^{}}`$ and (Wino-Higgsino) $`A_{\stackrel{~}{W}\stackrel{~}{h}^{}}`$ amplitudes, given by
$`{\displaystyle \frac{A_{\stackrel{~}{h}^{}}}{A_{SM}}}`$ $``$ $`\mathrm{tan}\beta (\mu m_t){\displaystyle \frac{M_{\stackrel{~}{t}_L\stackrel{~}{t}_R}^2}{m_{\stackrel{~}{t}_L}^2m_{\stackrel{~}{t}_R}^2}},`$ (47)
$`{\displaystyle \frac{A_{\stackrel{~}{W}\stackrel{~}{h}^{}}}{A_{SM}}}`$ $``$ $`{\displaystyle \frac{1}{V_{32}}}\mathrm{tan}\beta (\mu M_2){\displaystyle \frac{M_{\stackrel{~}{t}_L\stackrel{~}{t}_L}^2}{m_{\stackrel{~}{t}_L}^2m_{\stackrel{~}{t}_R}^2}}.`$ (48)
We see that these contributions are large with respect to the gluino one, such as in the universal case of MSSM, mainly because the mass insertions get a large enhancement of the light-stop mass $`m_{\stackrel{~}{t}_L}^2`$ in the denominator, unlike in the gluino case where the $`1/m_{\stackrel{~}{q}}^6`$ suppression is effective.
Regarding the contribution of the $`\stackrel{~}{Q}_{7,8}`$ operators to the total branching ratio, we checked that their effect is negligible, almost an order of magnitude smaller than the total contribution to $`Q_{7,8}`$. This can be explained by observing that the dominant effect to these operators comes from the gluino amplitude which is much smaller than the chargino or charged Higgs one.
In Figs. we plot the results for the branching ratio BR($`BX_s\gamma `$), in model A, versus $`\mathrm{sin}\theta `$ for different values of $`\mathrm{tan}\beta `$, the sign of $`\mu `$ and for two representative values of gravitino mass $`m_{3/2}`$ , namely $`m_{3/2}=150,300`$ GeV. The main message arising from these results is that the sensitivity of BR($`BX_s\gamma `$) respect to $`\mathrm{sin}\theta `$ increases with $`\mathrm{tan}\beta `$. In particular for the low $`\mathrm{tan}\beta `$region the $`bs\gamma `$ result does not differ significantly from the universal case. In the large $`\mathrm{tan}\beta `$region, $`\mathrm{tan}\beta =1540`$, the CLEO measurement of $`bs\gamma `$ set severe constraints on the angle $`\theta `$ for low gravitino masses. For $`\mu <0`$ almost the whole range of parameter space is excluded as shown in Fig. .
Comparing Top and Bottom plots in Figs., we see that that for positive (negative) sign of $`\mu `$ the branching ratio BR($`BX_s\gamma `$) decreases (increase) when the departure from the universality increases ($`\mathrm{sin}\theta 1/\sqrt{2}`$). Clearly, due to decoupling effects, the deviations from universality tend to be reduced for large gravitino masses, as can be seen by comparing the plots at $`m_{3/2}=150\mathrm{GeV}`$ with the corresponding ones at $`m_{3/2}=300\mathrm{GeV}`$.
Now we discuss results in model B. In Fig. we plot the branching ratio BR($`BX_s\gamma `$) versus $`\mathrm{tan}\beta `$for three different values of $`\mathrm{\Theta }_1,\mathrm{\Theta }_2`$ (see the figure caption) which are representative examples for universal and highly non–universal cases. From these figures it is clear that BR($`BX_s\gamma `$) is not very sensitive to the values of $`\mathrm{\Theta }_i`$’s parameters, even at very large $`\mathrm{tan}\beta `$, unlike model A. The constraints from CLEO measurement are almost the same in the universal and non–universal cases. For $`\mu >0`$ the branching ratio is constrained from the lower bound of CLEO only at very large $`\mathrm{tan}\beta `$, while for $`\mu <0`$ the branching ratio is almost excluded except at low $`\mathrm{tan}\beta `$.
## 5 Conclusions
The recent CP violation measurements of $`\epsilon ^{}/\epsilon `$ indicate a large deviation from the SM predictions which might be interpreted like a signal of new CP violation sources beyond the SM. It is unlikely that the minimal supersymmetric standard model with universal boundary conditions at GUT scale can explain these large enhancements in the direct CP violations. On the contrary, non-minimal SUSY scenarios with non–universal $`A`$–terms, derived from some string inspired models, have been found effective in explaining large values for $`\epsilon ^{}/\epsilon `$ while keeping the electric dipole moments below the experimental bounds.
In this paper we have considered two models based on weakly coupled heterotic string (model A) and type I string theories (model B). In this framework we carefully analysed the constraints set by the $`bs\gamma `$ decay on these two models by taking into account the relevant set of SUSY diagrams. In the calculation of the total branching ratio we take into account the NLO QCD corrections to the SM. We found that in the model based on the weakly coupled heterotic string, the $`bs\gamma `$ branching ratio is more sensitive to the non–universality at large $`\mathrm{tan}\beta `$and that the dominant SUSY contribution comes from the chargino amplitude for any value of $`\mathrm{tan}\beta `$. For type I string-derived model, we found that the sensitivity of the $`bs\gamma `$ branching ratio to the non–universality parameters $`\theta `$ and $`\mathrm{\Theta }_i`$ is quite weak. The main reason for this weakness is because in this model the allowed ranges for these parameters are strongly constrained by the vacuum stability bounds and the experimental limits on the lightest chargino mass.
We conclude that the recent CLEO measurements on the total inclusive B meson branching ratio BR($`BX_s\gamma `$) do not set severe constraints on the non–universality of these models. Moreover the constraints set on $`\mathrm{tan}\beta `$and gravitino mass are almost the same as in the universal case. In this respect we have found that the parameter regions which are important for generating sizeable contributions to $`\epsilon ^{}/\epsilon `$ , in particular the low $`\mathrm{tan}\beta `$regions, are not excluded by $`bs\gamma `$ decay.
## Acknowledgments
We would like to thank C. Muñoz for useful discussions. S. K. acknowledges the financial support of a Spanish Ministerio de Educacion y Cultura research grant. E.G. acknowledges the financial support of the TMR network, project “Physics beyond the standard model”, FMRX-CT96-0090. The work of E.T. was supported by DGICYT grant AEN97-1678.
## Appendix
Here we give the expressions for the dominant SUSY contributions to $`\stackrel{~}{R}_{7,8}`$ defined in Eq.(43), namely the chargino and gluino ones, in the approximation $`𝒪(m_s/m_b)=0`$
$`\stackrel{~}{R}_{7,8}`$ $`=`$ $`\stackrel{~}{R}_{7,8}^\chi +\stackrel{~}{R}_{7,8}^{\stackrel{~}{g}}`$
$`\stackrel{~}{R}_7^\chi `$ $`=`$ $`{\displaystyle \frac{2}{3V_{32}^{}V_{33}x_{tW}F_7(x_{tW})}}`$
$`\times `$ $`{\displaystyle \underset{I=1}{\overset{2}{}}}{\displaystyle \underset{k=1}{\overset{6}{}}}x_{W\stackrel{~}{u}_k}\left(X_I^L\right)_{k3}\left(X_I^R\right)_{k2}^{}{\displaystyle \frac{m_{\chi _I}}{m_b}}\left(F_3(x_{\chi _I\stackrel{~}{u}_k})+{\displaystyle \frac{2}{3}}F_4(x_{\chi _I\stackrel{~}{u}_k})\right)`$
$`\stackrel{~}{R}_8^\chi `$ $`=`$ $`{\displaystyle \frac{2}{3V_{32}^{}V_{33}x_{tW}F_1(x_{tW})}}{\displaystyle \underset{I=1}{\overset{2}{}}}{\displaystyle \underset{k=1}{\overset{6}{}}}x_{W\stackrel{~}{u}_k}\left(X_I^L\right)_{k3}\left(X_I^R\right)_{k2}^{}{\displaystyle \frac{m_{\chi _I}}{m_b}}F_4(x_{\chi _I\stackrel{~}{u}_k})`$
$`\stackrel{~}{R}_7^{\stackrel{~}{g}}`$ $`=`$ $`{\displaystyle \frac{16\alpha _S}{27\alpha _WV_{32}^{}V_{33}x_{tW}F_7(x_{tW})}}{\displaystyle \underset{k=1}{\overset{6}{}}}x_{W\stackrel{~}{d}_k}\left(\mathrm{\Gamma }^{D_L}\right)_{k3}\left(\mathrm{\Gamma }^{D_R}\right)_{k2}^{}{\displaystyle \frac{m_{\stackrel{~}{g}}}{m_b}}F_4(x_{\stackrel{~}{g}\stackrel{~}{d}_k})`$
$`\stackrel{~}{R}_8^{\stackrel{~}{g}}`$ $`=`$ $`{\displaystyle \frac{2\alpha _S}{9\alpha _WV_{32}^{}V_{33}x_{tW}F_1(x_{tW})}}`$ (49)
$`\times `$ $`{\displaystyle \underset{k=1}{\overset{6}{}}}x_{W\stackrel{~}{d}_k}\left(\mathrm{\Gamma }^{D_L}\right)_{k3}\left(\mathrm{\Gamma }^{D_R}\right)_{k2}^{}{\displaystyle \frac{m_{\stackrel{~}{g}}}{m_b}}\left(9F_3(x_{\stackrel{~}{g}\stackrel{~}{d}_k})+F_4(x_{\stackrel{~}{g}\stackrel{~}{d}_k})\right),`$
with
$`\left(X_I^L\right)_{ki}`$ $`=`$ $`𝐕_{I1}\left(\mathrm{\Gamma }^{U_L}V\right)_{ki}+{\displaystyle \frac{1}{\sqrt{2}m_W\mathrm{sin}\beta }}𝐕_{I2}\left(\mathrm{\Gamma }^{U_R}M_UV\right)_{ki}`$
$`\left(X_I^R\right)_{ki}`$ $`=`$ $`{\displaystyle \frac{1}{\sqrt{2}m_W\mathrm{cos}\beta }}𝐔_{I2}\left(\mathrm{\Gamma }^{U_L}VM_D\right)_{ki}`$ (50)
where $`x_{ij}m_i^2/m_j^2`$, $`F_7(x)=\frac{2}{3}F_1(x)+F_2(x)`$, $`𝐔`$ and $`𝐕`$ are the $`2\times 2`$ diagonalization matrices for the chargino mass matrix defined as in Ref. , and the $`6\times 6`$ matrices $`\mathrm{\Gamma }^{(U,D)}[\mathrm{\Gamma }_{6\times 3}^{(D_L,U_L)},\mathrm{\Gamma }_{6\times 3}^{(D_R,U_R)}]`$ respectively diagonalize the Up- and Down-squark mass matrices in Eqs.(33,35). The expressions for the functions $`F_i(x)`$ can be found in Ref..
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# Extracting and using photon polarization information in radiative 𝐵 decays
## I Introduction
Rare decays of $`B`$ mesons have attracted a lot of attention due to their ability to probe the existence of new physics. The most accessible such processes are weak radiative decays mediated by the quark process $`bs(d)\gamma `$. The CLEO Collaboration reported measurements of both exclusive and inclusive branching ratios for the $`bs\gamma `$ process, with results in good agreement with the Standard Model (SM) predictions . This still leaves open the possibility that new physics could be present, but it manifests itself only in the details of the decay process, such as polarization effects or differential distributions. A number of methods have been proposed which could detect deviations from the SM predictions along these lines .
One particular class of methods is based on the SM prediction that the photons emitted in $`bs\gamma `$ decays are predominantly left-handed. (Long-distance effects and the light-quark masses introduce a small right-handed component, which can be neglected to a first approximation.) This property does not hold true in extensions of the SM such as the left-right symmetric model (LRM) , and therefore, can be used to signal the presence of new physics. Unfortunately, in $`B`$ decays all photon polarization information contained in the final hadron is lost. Since the photon polarization in $`\overline{B}K^{}\gamma `$ decays is difficult to detect, most polarization-based methods focused on the related $`bs\mathrm{}^+\mathrm{}^{}`$ decay, where the angular distributions and lepton polarizations can probe the chiral structure of the short-distance matrix element .
Detecting an unambiguous signal for new physics entails a good control over the SM prediction for the respective exclusive modes. While the short-distance contribution can be parametrized in terms of (more or less well-known) hadronic form-factors, the charm- and up-quark loops introduce hard-to-calculate long-distance contributions. In $`bs\mathrm{}^+\mathrm{}^{}`$ decays, these effects become significant when $`q^2`$ (the $`e^+e^{}`$ invariant mass) is in the region of the charmonium excitations, and are minimal at the lower end ($`q^20`$). Furthermore, the differential rate is enhanced at this point due to the small photon propagator denominator. The remaining long-distance contributions, connected with weak annihilation or $`W`$-exchange topologies, can be computed in an expansion in $`1/q_0`$ . This suggests that the $`bs\gamma `$ decay, respectively the $`bs\mathrm{}^+\mathrm{}^{}`$ decay in the small-$`q^2`$ region, is a good place to search for new physics through photon polarization effects.
Measuring mixing-induced CP asymmetries in the inclusive $`bs\gamma `$ decay was proposed as an indirect method for probing photon polarization effects in . Since both $`B`$ and $`\overline{B}`$ must decay to a common final state, the resulting asymmetry measures the interference of right- and left-handed photon amplitudes. As the SM predicts a very small right-handed admixture of photons in $`bq\gamma `$ decays, a large mixing-induced CP asymmetry is a signal of new physics.
We explore in this paper an alternative way of measuring the photon polarization in the exclusive $`BV\gamma `$ decay, which makes use of the conversion electron pairs formed by the primary photon. Electron–positron pairs from photons that were produced in the inner part of the detector can be traced and their decay plane can be reconstructed with high accuracy. For example, at BaBar about 3% of the photons are expected to convert on the beam pipe and the silicon detector . We show in Sec. II how the conversion process can be used to extract information about the photon polarization in $`BV\gamma `$ decays, in analogy to the classical example of determining the $`\pi ^0`$ parity through the chain decay $`\pi ^0\gamma (e^+e^{})\gamma (e^+e^{})`$ . The main ingredient of this analysis is the fact that in $`B`$ decay into two vectors, the subsequent decay (or conversion) of the vectors contains their relative polarization information. In particular, the angular distribution in the relative angle of the $`K^{}K\pi `$ decay plane and that of the conversion pair can be used to determine the helicity amplitudes in the $`BV\gamma `$ decay. In Sec. III we discuss time-dependent CP asymmetries in the angular distribution of the conversion pairs produced in neutral $`B`$ decays, which measure $`\mathrm{sin}2\beta `$ and $`\mathrm{cos}2\beta `$ with little hadronic uncertainty. We comment in Sec. IV on the experimental aspects of the methods proposed here.
## II Conversion lepton pairs
There are two possible mechanisms by which the photon (real or virtual) emitted in $`bs\gamma `$ decay can convert into a lepton pair. In the first one, a virtual photon with momentum $`q^2(2m_e)^2`$ produces an electron-positron pair, which are emitted without any other interaction. In the second mechanism, a real photon produces a lepton pair which subsequently interacts with a nucleus by Coulomb forces (the so-called Bethe-Heitler process ).
The lepton pair produced in these two processes are seen very differently in a detector. The first mechanism produces prompt lepton pairs, which originate practically from the interaction region. The pairs would be produced even in vacuum, in the absence of any matter content of the detector. On the other hand, the Bethe-Heitler process produces lepton pairs within the volume of the detector with a probability proportional to the density of matter.
These arguments can be illustrated with a simple estimate as follows. The lifetime of a virtual photon contributing to the first mechanism is (in its rest frame) of the order $`\tau _01/\sqrt{q^2}O(1/m_e)`$. In the lab frame the lifetime is longer by a factor of $`\gamma =q_0/\sqrt{q^2}O(m_B/m_e)`$. Thus the photon travels a distance $`\mathrm{\Delta }xm_B/m_e^210^6`$ $`mm`$ before decaying. Clearly, any photon that has $`q^2>(2m_e)^2`$ travels a distance that is too short to be measured. For this reason we will refer to the lepton pairs produced in these two mechanisms as to short-distance and long-distance conversion leptons, respectively. In the following we examine them in turn.
### A Short-distance lepton pairs
In the Standard Model, the decays $`\overline{B}X_se^+e^{}`$ are mediated by a combination of the short-distance Hamiltonian
$`_{\mathrm{s}.\mathrm{d}.}={\displaystyle \frac{4G_F}{\sqrt{2}}}V_{tb}V_{ts}^{}{\displaystyle \underset{i=710}{}}C_i(\mu )𝒪_i(\mu )`$ (1)
with
$`𝒪_7`$ $`=`$ $`{\displaystyle \frac{e}{16\pi ^2}}\overline{s}\sigma _{\mu \nu }(m_bP_R+m_sP_L)bF^{\mu \nu },𝒪_8={\displaystyle \frac{g}{16\pi ^2}}\overline{s}\sigma _{\mu \nu }(m_bP_R+m_sP_L)T^abG^{a\mu \nu }`$ (2)
$`𝒪_9`$ $`=`$ $`{\displaystyle \frac{e^2}{16\pi ^2}}(\overline{s}\gamma _\mu P_Lb)(\overline{e}\gamma ^\mu e),𝒪_{10}={\displaystyle \frac{e^2}{16\pi ^2}}(\overline{s}\gamma _\mu P_Lb)(\overline{e}\gamma ^\mu \gamma _5e)`$ (3)
and nonlocal contributions introduced by the usual weak nonleptonic Hamiltonian
$`_W`$ $`=`$ $`{\displaystyle \frac{4G_F}{\sqrt{2}}}\{V_{ub}V_{us}^{}[C_1𝒪_1^{(u)}+C_2𝒪_2^{(u)}]+V_{cb}V_{cs}^{}[C_1𝒪_1^{(c)}+C_2𝒪_2^{(c)}]`$ (5)
$`V_{tb}V_{ts}^{}{\displaystyle \underset{i=3}{\overset{10}{}}}C_i(\mu )𝒪_i^{(s)}+(sd)\}`$
with
$`𝒪_1^{(q)}=(\overline{q}\gamma _\mu P_Lb)(\overline{s}\gamma ^\mu P_Lq),𝒪_2^{(q)}=(\overline{s}\gamma _\mu P_Lb)(\overline{q}\gamma ^\mu P_Lq).`$ (6)
With new physics other operators can also contribute. In particular, the right handed operators, which can be obtained from the SM operators by $`RL`$, are denoted by $`\stackrel{~}{𝒪}_i`$.
The amplitude for $`\overline{B}Ve^+e^{}`$ decay, with $`V`$ a light vector meson, can be written as
$$=\frac{4G_F}{\sqrt{2}}V_{tb}V_{ts}^{}e\left\{(A_\mu +\frac{1}{q^2}H_\mu )[\overline{u}(p_{e_+})\gamma ^\mu v(p_e^{})]+B_\mu [\overline{u}(p_{e_+})\gamma ^\mu \gamma _5v(p_e^{})]\right\},$$
(7)
with $`A,B,H`$ hadronic matrix elements. $`B_\mu `$ receives only contributions from the operator $`𝒪_{10}`$ above, and can be expressed in terms of form-factors alone. $`A`$ and $`H`$ can receive contributions from all the other operators. The part which contains a pole at $`q^2=0`$, arising from the propagator of an intermediate photon, dominates in the small-$`q^2`$ region, in which we will be interested in the following. Therefore, we will neglect the nonpole terms proportional to $`A_\mu `$ and $`B_\mu `$.
The matrix element $`H_\mu `$ is parametrized in the most general form in terms of three invariant form factors
$`H_\mu (q^2)=A_{}(q^2)\left(q_0ϵ_{V\mu }^{}v_\mu (qϵ_V^{})\right)+A_0(q^2)v_\mu (qϵ_V^{})+A_{}(q^2)i\epsilon (q,\mu ,ϵ_V^{},v).`$ (8)
The amplitudes $`A_{}(q^2)`$ (CP-even) and $`A_{}(q^2)`$ (CP-odd) correspond to the virtual photon polarization being parallel, respectively transverse to that of the vector meson. $`A_0(q^2)`$ (CP-even) is related to the longitudinal polarization. The values of these form factors at $`q^2=0`$ describe the coupling of a real photon in the weak radiative decay $`\overline{B}V\gamma `$. Thus, in general, $`A_0(0)=0`$ and in the SM, where the photon is mostly left handed, we have $`A_{}(0)A_{}(0)`$.
An alternative description of the photon coupling is in terms of helicity amplitudes, giving the amplitude for the $`\overline{B}`$ meson to decay into a virtual photon of well-defined helicity. They are related to the invariant form factors in (8) by
$`A_{\mathrm{}}(q^2)`$ $`=`$ $`\sqrt{q^2}{\displaystyle \frac{m_Bq_0}{m_V}}A_{}(q^2)+{\displaystyle \frac{\stackrel{}{q}^2m_B}{m_V\sqrt{q^2}}}A_0(q^2)`$ (9)
$`A_{R,L}(q^2)`$ $`=`$ $`q_0A_{}(q^2)|\stackrel{}{q}|A_{}(q^2).`$ (10)
Here $`A_R`$ is the right (left) handed polarization amplitude and $`A_{\mathrm{}}`$ is the longitudinal polarization amplitude which clearly vanish in the $`q^2=0`$ limit. Expressed in terms of helicity amplitudes, the rate for radiative decay $`\overline{B}V\gamma `$ is given by
$`\mathrm{\Gamma }(\overline{B}V\gamma )={\displaystyle \frac{G_F^2|V_{ts}V_{tb}^{}|^2}{\pi m_B^2}}E_\gamma (|A_R(0)|^2+|A_L(0)|^2),`$ (11)
where $`E_\gamma `$ is the photon energy in the $`\overline{B}`$ rest frame.
In the SM, the $`\overline{B}V\gamma `$ amplitude is dominated by the operator $`𝒪_7`$, which contributes mainly to $`A_L`$, with a small right-handed admixture due to long-distance effects and light quark masses. In certain extensions of the SM, such as the left-right symmetric model , a new penguin operator $`\stackrel{~}{𝒪}_7`$ is introduced involving right-handed photons. This operator can make a significant contribution to $`A_R`$. Thus, the photon amplitudes are given, in the general case, by
$`A_R(0)`$ $`=`$ $`(\stackrel{~}{C}_7m_b+C_7m_q){\displaystyle \frac{e}{16\pi ^2}}2(m_B^2m_V^2)g_+(0)+a_R,`$ (12)
$`A_L(0)`$ $`=`$ $`(C_7m_b+\stackrel{~}{C}_7m_q){\displaystyle \frac{e}{16\pi ^2}}2(m_B^2m_V^2)g_+(0)+a_L.`$ (13)
where the form factor $`g_+(q^2)`$ is defined by
$`V(p^{})|\overline{q}\sigma _{\mu \nu }b|B(p)`$ $`=`$ $`g_+(q^2)\epsilon _{\mu \nu \lambda \sigma }ϵ_\lambda ^{}(p+p^{})_\sigma +g_{}(q^2)\epsilon _{\mu \nu \lambda \sigma }ϵ_\lambda ^{}(pp^{})_\sigma `$ (15)
$`+h(q^2)\epsilon _{\mu \nu \lambda \sigma }(p+p^{})_\lambda (pp^{})_\sigma (ϵ^{}p).`$
$`m_q`$ is the mass of the strange or down quark depending on the specific decay. The long-distance amplitudes $`a_{L,R}`$ are introduced by the $`bc\overline{c}s`$ part of the weak Hamiltonian and are expected to be about 5% of the short-distance contribution .
We emphasize here that in the SM the photon in $`BV\gamma `$ is almost pure left-handed. The small right handed component due to long distance effects and light quark masses is at the 5% level. Thus, any measurement of a significant right handed amplitude will be an unambiguous signal for new physics.
A measurement of the individual helicity amplitudes $`A_{L,R}(0)`$ can therefore give useful information about the short-distance weak radiative Hamiltonian. We will show in the following how a study of the decay $`\overline{B}Ve^+e^{}`$ in the low $`e^+e^{}`$ invariant mass region can be useful in this respect. The argument is a simple adaptation of the classical analysis of Kroll and Wada given in .
Let us take the virtual photon to be moving along the $`+z`$ axis, and the final state meson $`V`$ in the $`z`$ direction, with momenta $`\stackrel{}{q}`$ and $`\stackrel{}{q}`$ respectively. The photon converts into a $`e^+e^{}`$ pair, with the $`e^+`$ moving at an azimuthal angle $`\theta `$ with respect to the $`+z`$ axis and a polar angle $`\varphi `$. The vector meson decays into two pseudoscalars, which will be denoted generically by $`K(p_K)`$ and $`\pi (p_\pi )`$ (corresponding to the interesting case $`V=K^{}`$). The pion momentum $`p_\pi `$ is parametrized by the angles $`(\psi ,0)`$ with respect to the $`z`$ direction. The differential decay rate in these coordinates is given by
$`{\displaystyle \frac{\text{d}\mathrm{\Gamma }}{\text{d}q^2\text{d}\mathrm{cos}\theta \text{d}\mathrm{cos}\psi \text{d}\varphi }}`$ $`=`$ $`𝒞\{4|A_{\mathrm{}}|^2\mathrm{sin}^2\theta \mathrm{cos}^2\psi +(|A_R|^2+|A_L|^2)(1+\mathrm{cos}^2\theta )\mathrm{sin}^2\psi `$ (18)
$`\mathrm{sin}2\theta \mathrm{sin}2\psi \left\{\mathrm{cos}\varphi [\text{Re}(A_{\mathrm{}}A_R^{})+\text{Re}(A_{\mathrm{}}A_L^{})]+\mathrm{sin}\varphi [\text{Im}(A_{\mathrm{}}A_R^{})\text{Im}(A_{\mathrm{}}A_L^{})]\right\}`$
$`2[\text{Re}(A_RA_L^{})\mathrm{cos}2\varphi \text{Im}(A_RA_L^{})\mathrm{sin}2\varphi ]\mathrm{sin}^2\theta \mathrm{sin}^2\psi \}.`$
The constant $`𝒞`$ is given by
$`𝒞={\displaystyle \frac{3\alpha G_F^2|V_{tb}V_{ts}^{}|^2}{8(2\pi )^3m_B^3}}{\displaystyle \frac{\sqrt{\lambda }}{q^2}},\lambda ={\displaystyle \frac{1}{4}}(m_B^2q^2m_V^2)^2q^2m_V^2.`$ (19)
We assumed in deriving (18) that the final leptons are massless, and neglected the parity-violating effects in the decay $`\gamma ^{}e^+e^{}`$. Such effects are introduced by $`Z`$ boson exchange (the operator $`𝒪_{10}`$), and are parametrized by the hadronic matrix element $`B_\mu `$ in (7). They are negligibly small in the small $`q^2`$ region we consider here. The form factor $`A_0(q^2)`$ vanishes at $`q^2=0`$ as $`A_0(q^2)q^2`$, such that the amplitude for producing longitudinally polarized real photons $`A_{\mathrm{}}(q^2)`$ vanishes for $`q^2=0`$, as expected.
From (18) one obtains, after integrating over $`(\theta ,\psi )`$, the following $`\varphi `$ distribution
$`{\displaystyle \frac{d\mathrm{\Gamma }}{d\varphi }}`$ $`=`$ $`{\displaystyle \frac{32}{9}}{\displaystyle _{(2m_e)^2}^{q_{\mathrm{max}}^2}}dq^2𝒞(q^2)\{(|A_{\mathrm{}}(q^2)|^2+|A_R(q^2)|^2+|A_L(q^2)|^2)`$ (21)
$`[\text{Re}(A_R(q^2)A_L^{}(q^2))\mathrm{cos}2\varphi \text{Im}(A_R(q^2)A_L^{}(q^2))\mathrm{sin}2\varphi ]\}.`$
In the region close to threshold, the helicity amplitude for producing longitudinally polarized photons has the asymptotic form $`|A_{\mathrm{}}(q^2)|^2q^2`$. Furthermore, to a first approximation, one can neglect the $`q^2`$-variation of the transverse helicity amplitudes $`|A_{R,L}(q^2)|`$ in this region and set them equal to their values at $`q^2=0`$. Therefore, the $`q^2`$-integral can be approximated as
$`{\displaystyle \frac{d\mathrm{\Gamma }}{d\varphi }}`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}\mathrm{\Gamma }(\overline{B}V\gamma )\left({\displaystyle \frac{\alpha }{3\pi }}\mathrm{log}{\displaystyle \frac{q_{\mathrm{max}}^2}{(2m_e)^2}}\right)`$ (23)
$`\times \left\{1{\displaystyle \frac{\text{Re}(A_R(0)A_L^{}(0))\mathrm{cos}2\varphi \text{Im}(A_R(0)A_L^{}(0))\mathrm{sin}2\varphi }{|A_R(0)|^2+|A_L(0)|^2}}\right\}+\mathrm{}.`$
The ellipsis denote a neglected contribution from the longitudinal amplitude, proportional to the integral $`𝑑q^2/q^2|A_{\mathrm{}}(q^2)|^2`$. At $`q_{\mathrm{max}}^2=1`$ GeV<sup>2</sup>, it amounts to about 2% of the leading terms which are kept. Note the presence of the large logarithm $`\mathrm{log}\frac{q_{\mathrm{max}}^2}{(2m_e)^2}`$, which can partly compensate the additional suppression of $`\alpha `$ compared to the purely radiative rate (11). Numerically, the value of the suppression factor in brackets is $`\frac{\alpha }{3\pi }\mathrm{log}\frac{q_{\mathrm{max}}^2}{(2m_e)^2}0.01`$ at $`q_{\mathrm{max}}^2=1`$ GeV<sup>2</sup>. This implies an effective branching ratio of few times $`10^7`$ for events with $`q^2`$ in the region of interest.
From (23) one can see that from the $`\varphi `$ dependence in the $`q^2`$-integrated rate we can extract the ratio
$`R{\displaystyle \frac{|A_R(0)||A_L(0)|}{|A_R(0)|^2+|A_L(0)|^2}}.`$ (24)
Combining this with the total exclusive rate (11), the individual helicity amplitudes for right- and left-handed photons can be extracted (up to a $`A_R(0)A_L(0)`$ ambiguity). In the SM we expect $`R5\%`$. Therefore, by measuring $`R`$ we are sensitive to new physics amplitudes that are an order of magnitude smaller than the SM amplitude.
Angular distributions of the type (23) in the low dilepton invariant mass region were also discussed in . There, the full expressions for the helicity amplitudes are kept, including parity-violating effects induced by the operator $`𝒪_{10}`$. The resulting form of the angular distribution depends on many form-factors and is therefore not easily connected to the parameters of the short-distance Hamiltonian. In contrast, the phenomenological analysis presented here is model-independent; by restricting to a sufficiently small region above $`q^2=0`$, the radiative helicity amplitudes can be directly extracted, without need for any additional form factors.
### B Long-distance lepton pairs
Angular correlations in lepton pair production by a real photon have been suggested long ago as a means for measuring photon polarization (for a review see ). Our discussion here will focus on aspects relevant to the photons emitted in exclusive radiative $`B`$ decays.
The cross-section for pair production by a polarized photon was computed in . To lowest order it is given by
$`{\displaystyle \frac{d\sigma }{dE_1d\mathrm{\Omega }_1d\mathrm{\Omega }_2}}`$ $`=`$ $`{\displaystyle \frac{Z^2e^6}{16\pi ^3}}{\displaystyle \frac{|\stackrel{}{p}_1||\stackrel{}{p}_2|}{E_\gamma ^3k^4}}\{{\displaystyle \frac{(k^24E_2^2)(\stackrel{}{e}\stackrel{}{p}_1)(\stackrel{}{e}^{}\stackrel{}{p}_1)}{(E_1|\stackrel{}{p}_1|\mathrm{cos}\theta _1)^2}}+{\displaystyle \frac{(k^24E_1^2)(\stackrel{}{e}\stackrel{}{p}_2)(\stackrel{}{e}^{}\stackrel{}{p}_2)}{(E_2|\stackrel{}{p}_2|\mathrm{cos}\theta _2)^2}}`$ (25)
$``$ $`{\displaystyle \frac{k^2+4E_1E_2}{(E_1|\stackrel{}{p}_1|\mathrm{cos}\theta _1)(E_2|\stackrel{}{p}_2|\mathrm{cos}\theta _2)}}[(\stackrel{}{e}\stackrel{}{p}_1)(\stackrel{}{e}^{}\stackrel{}{p}_2)+(\stackrel{}{e}^{}\stackrel{}{p}_1)(\stackrel{}{e}\stackrel{}{p}_2)]`$ (26)
$`+`$ $`E_\gamma ^2{\displaystyle \frac{\stackrel{}{p}_1^2\mathrm{sin}^2\theta _1+\stackrel{}{p}_2^2\mathrm{sin}^2\theta _2+2|\stackrel{}{p}_1||\stackrel{}{p}_2|\mathrm{sin}\theta _1\mathrm{sin}\theta _2\mathrm{cos}(\varphi _1\varphi _2)}{(E_1|\stackrel{}{p}_1|\mathrm{cos}\theta _1)(E_2|\stackrel{}{p}_2|\mathrm{cos}\theta _2)}}\}.`$ (27)
We denote here with $`\stackrel{}{e}`$ the photon polarization vector and $`p_1(p_2)`$ the positron (electron) momenta. $`k=p_1+p_2q`$ is the momentum transferred to the nucleus, which will be taken to be infinitely heavy. In this limit the nucleus does not carry away any energy and the photon energy is transferred entirely to the electron pair $`E_\gamma =E_1+E_2`$. The angles of the positron and electron with respect to the photon momentum direction are denoted with $`\theta _1`$ and $`\theta _2`$ respectively.
We propose to use as polarization analyzer the rate (25) integrated over the electron direction $`\mathrm{\Omega }_2`$. For a linearly polarized photon, the integrated cross-section for pair production has the following dependence on the angle $`\alpha `$ between the photon polarization vector and the projection of the positron momentum $`\stackrel{}{p}_1`$ on the plane transverse to the photon momentum
$`{\displaystyle \frac{d\sigma }{dE_1d\mathrm{\Omega }_1}}=\sigma _I+{\displaystyle \frac{1}{2}}(\sigma _{II}\sigma _{III})\mathrm{cos}2\alpha +\sigma _{IV}\mathrm{sin}2\alpha .`$ (28)
We used here the notations of for the pair production cross-section by a polarized photon. Only the cross-sections $`\sigma _I,\sigma _{II}`$ and $`\sigma _{III}`$ have been computed in this paper (see Eqs. (17) in ) corresponding to unpolarized photons ($`\sigma _I`$), and linearly polarized photons with the $`(\stackrel{}{p}_1,\stackrel{}{q})`$ plane parallel to the polarization plane ($`\sigma _{II}`$) and transverse to it ($`\sigma _{III}`$). One has $`\sigma _I=\frac{1}{2}(\sigma _{II}+\sigma _{III})`$. The cross-section $`\sigma _{IV}`$ measures the acoplanarity of the three vectors $`\stackrel{}{q},\stackrel{}{p}_1,\stackrel{}{p}_2`$. The lepton pair is predominantly produced such that the photon momentum $`\stackrel{}{q}`$ lies in the plane of the pair. Therefore, $`\sigma _{IV}`$ is very small and will be neglected in the numerical estimates below.
Once the vector meson in $`BV\gamma `$ is observed through its decay to a pair of pseudoscalars (whose decay plane defines the $`x`$ axis), the polarization vector of the emitted photon is fixed to be $`\stackrel{}{e}A_{}\stackrel{}{e}_1+iA_{}\stackrel{}{e}_2`$. The angular distribution of the positron momentum direction is then
$`{\displaystyle \frac{d\sigma }{dE_1d\mathrm{\Omega }_1}}`$ $`=`$ $`\mathrm{\Gamma }(BV\gamma )\sigma _I\{1+\mathrm{cos}2\varphi [{\displaystyle \frac{\sigma _{II}\sigma _{III}}{\sigma _I}}{\displaystyle \frac{\text{Re }(A_RA_L^{})}{|A_R|^2+|A_L|^2}}+{\displaystyle \frac{2\sigma _{IV}}{\sigma _I}}{\displaystyle \frac{\text{Im }(A_RA_L^{})}{|A_R|^2+|A_L|^2}}]`$ (30)
$`+\mathrm{sin}2\varphi [{\displaystyle \frac{\sigma _{II}\sigma _{III}}{\sigma _I}}{\displaystyle \frac{\text{Im }(A_RA_L^{})}{|A_R|^2+|A_L|^2}}+{\displaystyle \frac{2\sigma _{IV}}{\sigma _I}}{\displaystyle \frac{\text{Re }(A_RA_L^{})}{|A_R|^2+|A_L|^2}}]\},`$
where the angle $`\varphi `$ defined before Eq. (18). The dependence on $`\varphi `$ has the form
$$\frac{d\sigma }{d\varphi }1+\xi R\mathrm{cos}(2\varphi +\delta ),$$
(31)
with $`R`$ the ratio of amplitudes defined in (24), and $`\xi `$ an efficiency factor which can be expressed in terms of the cross-sections $`\sigma _{IIV}`$ as
$`\xi (E_1,\theta _1)=\sqrt{\left({\displaystyle \frac{\sigma _{II}\sigma _{III}}{\sigma _I}}\right)^2+4\left({\displaystyle \frac{\sigma _{IV}}{\sigma _I}}\right)^2},\xi 2.`$ (32)
Most of the pairs produced in the Bethe-Heitler process are emitted in a small cone of opening angle $`2\theta `$ with $`\theta m/E_\gamma 0.01^{}`$. We plot in Fig. 1 the efficiency parameter $`\xi (E_1,\theta _1)`$ as a function of $`\theta _1`$ in this range, for three values of the positron energy $`E_1`$. The sensitivity to the photon polarization is smaller by about a factor of 2 than for the short-distance lepton pairs (see Eq. (23)), and is maximal when the photon energy is equally distributed among the two leptons $`(E_1=E_\gamma /2)`$.
In Fig. 1 we show also the effective $`\xi `$ parameter obtained when the positron energy is not measured. This is defined as in (32) but in terms of the cross-sections integrated over $`E_1`$. Averaging over $`E_1`$ further decreases the sensitivity to the photon polarization. An alternative method has been discussed in the literature which improves the sensitivity at the cost of statistics. This method uses as polarization analyzer the rate (25) integrated over a restricted region of the electron direction $`\mathrm{\Omega }_2`$, chosen such that the three vectors $`\stackrel{}{q},\stackrel{}{p}_1,\stackrel{}{p}_2`$ are almost coplanar $`(|\varphi _2\varphi _1\pi |\mathrm{\Delta }\varphi )`$. A detailed discussion of the resulting asymmetry as a function of the width $`\mathrm{\Delta }\varphi `$ can be found in . However, the gain in sensitivity of this method may be offset by a loss in statistics involved by integrating over a restricted region in $`\mathrm{\Omega }_2`$.
To conclude this section we stress again the main point. In the SM, as $`R`$ is very small, there is almost no angular dependence in the electron-positron conversion rate. Any significant measurement of such angular dependence will be an indication of new physics. In principle, if indeed such new physics exists, using the formulae presented in this section one could also determine the relative size of this new physics amplitude. While this may be hard to achieve, the modest goal of demonstrating any angular distribution may be experimentally feasible if $`RO(1)`$.
## III Time-dependent angular distributions
A different aspect of polarization effects in weak radiative decays is manifested through time-dependence in neutral $`B`$ decays. Assuming the validity of the SM, we will show that certain time-dependent CP asymmetries involving real or virtual photons can be used to measure $`\mathrm{sin}2\beta `$ and $`\mathrm{cos}2\beta `$ with very little hadronic uncertainty. While the measurement of the polarization is sensitive to the right-handed operator $`\stackrel{~}{𝒪}_7`$, the CP asymmetry is sensitive to a new CP violating amplitude independent of its helicity. In the presence of any new contribution to the decay amplitude with a weak phase that is different from the SM phase, the “would be” $`\mathrm{sin}2\beta `$ measured in the radiative decay would not agree with the one measured in $`B\psi K_S`$. Such a disagreement will be a clean signal for a new CP violating amplitude in $`bs\gamma `$ . In addition, the sign of $`\mathrm{cos}2\beta `$ can be used to resolve a discrete ambiguity in the value of $`2\beta `$ deduced from the measurement of $`\mathrm{sin}2\beta `$ . Therefore, this measurement is sensitive to new CP violating contributions to the mixing amplitude.
Before developing the formalism we explain below why we gain sensitivity to CP violation phases by using polarization information. Atwood et al. studied the time-dependent CP asymmetries in $`B^0(t)V\gamma `$ where no polarization information is obtained. They conclude that in the SM the asymmetry almost vanishes and only in the presence of right-handed amplitude there is going to be an asymmetry. The reason for this is simple. In the SM, $`\overline{B}^0`$ decays into a left-handed photon, while $`B^0`$ decays into a right-handed photon. However, interference is necessary in order to produce an asymmetry. Therefore, a final state that is accessible only from $`B^0`$ or $`\overline{B}^0`$ does not produce any asymmetry. In contrast, when we consider decays into a linear polarized photon, both $`B^0`$ and $`\overline{B}^0`$ can decay into the same final state. This is because the linear polarization state contains an equal mixture of the circular polarization states. Moreover, unless there are at least two amplitudes with different weak and strong phases, the magnitudes of the decay amplitudes into a linear polarized state for both $`B^0`$ or $`\overline{B}^0`$ are the same.
The situation is very similar to the well known $`B^0\psi K`$ decay. The $`B^0`$ decays only to $`\psi K^0`$ while the $`\overline{B}^0`$ decays into $`\psi \overline{K}^0`$ and there is no interference between the two decays. Indeed, if we measure $`B\psi K`$ without determining any property of the final kaon, we do not get any asymmetry. The situation is very different when we look into final state kaons that are admixtures of $`K`$ and $`\overline{K}`$, namely $`K_S`$ and $`K_L`$. In that case, both $`B^0`$ and $`\overline{B}^0`$ decay to the same final state and the asymmetry can be used to measure $`\mathrm{sin}2\beta `$ in the SM. In both $`BV\gamma `$ and $`B\psi K_S`$ cases the situation is the same: in order to measure the asymmetry we have to observe final states that are accessible from both $`B^0`$ and $`\overline{B}^0`$.
At this point we already know what can be measured in the CP asymmetry in $`BK^{}\gamma `$ with linear polarized photon. Since these final states are CP eigenstates we can use the well known formalism of CP asymmetries in $`B`$ decays into CP eigenstates. In the SM (working in the Wolfenstein parametrization) the $`bs\gamma `$ amplitude has a trivial weak phase, and thus the asymmetry is sensitive to the mixing amplitude, namely to $`2\beta `$. Moreover, since we have many amplitudes that interfere we are sensitive to both $`\mathrm{sin}2\beta `$ and $`\mathrm{cos}2\beta `$. Below we show how to extract $`\mathrm{sin}2\beta `$ and $`\mathrm{cos}2\beta `$ from the angular distribution information.
In $`B^0V\gamma `$ decays, final states of well-defined CP correspond to the amplitudes $`A_{}(0)`$ and $`A_{}(0)`$, rather than to states of well-defined helicity. Furthermore, if $`V=K^0,\overline{K^0}`$, one must require that the final state be identified through the decay $`K^0K_S\pi ^0`$, which is CP-odd. We denote the corresponding amplitudes in $`\overline{B}^0V\gamma `$ decays by $`\overline{A}_{}`$ and $`\overline{A}_{}`$<sup>*</sup><sup>*</sup>*Note the change in notation compared to Sec. II. To conform with usual conventions, $`B(\overline{B})`$ decay amplitudes will be denoted with $`A(\overline{A})`$, whereas in Sec. II we dealt only with $`\overline{B}`$ decay amplitudes..
If the decaying meson is tagged as $`B^0(\overline{B}^0)`$ at time $`t=0`$, then the amplitudes $`A_{}`$ and $`A_{}`$ will depend on $`t`$ at a later time $`t`$. This time dependence is given by
$`A_{}(t)`$ $`=`$ $`A_{}(0)\left(f_+(t)+\lambda _{}f_{}(t)\right)`$ (33)
$`\overline{A}_{}(t)`$ $`=`$ $`{\displaystyle \frac{p}{q}}A_{}(0)\left(f_{}(t)+\lambda _{}f_+(t)\right),`$ (34)
and analogous for $`A_{}`$ with a different parameter $`\lambda _{}`$. Our notation for the $`B\overline{B}`$ mixing parameters $`p`$, $`q`$ is the standard one . The time-dependence is contained in the functions
$$f_\pm (t)=\frac{1}{2}\left\{e^{i(m_1i\mathrm{\Gamma }_1/2)t}\pm e^{i(m_2i\mathrm{\Gamma }_2/2)t}\right\},$$
(35)
with $`m_{1,2}`$ and $`\mathrm{\Gamma }_{1,2}`$ the masses and widths of the mass eigenstates of the $`B\overline{B}`$ system. The parameters $`\lambda _{}`$ and $`\lambda _{}`$ are defined as usual by
$`\lambda _{}={\displaystyle \frac{q}{p}}{\displaystyle \frac{\overline{A}_{}(0)}{A_{}(0)}},\lambda _{}={\displaystyle \frac{q}{p}}{\displaystyle \frac{\overline{A}_{}(0)}{A_{}(0)}}.`$ (36)
In the following we concentrate on $`B^0`$ decay via the $`bs\gamma `$ quark level transition. Within the SM, the dominance of the left-handed amplitude implies
$`A_{}A_{},\overline{A}_{}\overline{A}_{},`$ (37)
where we used $`A_L=A_{}+A_{}0`$ and $`\overline{A}_R=\overline{A}_{}\overline{A}_{}0`$. Moreover, in the SM the $`bs\gamma `$ decay amplitude has a trivial weak phase (in the Wolfenstein parametrization) and the ratio $`q/p=e^{2i\beta }`$, which gives
$`\lambda _{}`$ $`=`$ $`{\displaystyle \frac{q}{p}}{\displaystyle \frac{\overline{A}_L+\overline{A}_R}{A_L+A_R}}e^{2i\beta },\lambda _{}={\displaystyle \frac{q}{p}}{\displaystyle \frac{\overline{A}_R\overline{A}_L}{A_RA_L}}e^{2i\beta }.`$ (38)
The above results can be generalized to many extensions of the SM. When there is new CP conserving contribution to $`A_R`$, they are not modified. When there is new CP violating contribution to $`A_L`$ the above results still hold, where one must replace $`2\beta `$ in both $`\lambda _{}`$ and $`\lambda _{}`$ with the angle between the mixing and the decay amplitude. When there is new CP violating contribution to $`A_R`$, but no strong phase between the left and right handed amplitude, $`\lambda _{}`$ and $`\lambda _{}`$ are still pure phase. However, the phase is not the same. Finally, when there is also a strong phase between the left and right handed amplitudes $`|\lambda |1`$. While we again assume the SM in the following discussion, our results hold also for the first two cases discuss above (with the general interpretation of $`2\beta `$). It is clear how to generalize the results below also to more general cases.
Neglecting, as usual, the lifetime difference of the $`B_d`$ mass eigenstates and assuming that $`|\lambda |1`$ one finds that each of the time-dependent CP asymmetries for the final states of linear polarization measure $`\mathrm{sin}2\beta `$
$`a_{}(t)`$ $`=`$ $`{\displaystyle \frac{|A_{}(t)|^2|\overline{A}_{}(t)|^2}{|A_{}(t)|^2+|\overline{A}_{}(t)|^2}}=\text{Im}\lambda _{}\mathrm{sin}(\mathrm{\Delta }mt)\mathrm{sin}2\beta \mathrm{sin}(\mathrm{\Delta }mt)\text{ (SM)}`$ (39)
$`a_{}(t)`$ $`=`$ $`{\displaystyle \frac{|A_{}(t)|^2|\overline{A}_{}(t)|^2}{|A_{}(t)|^2+|\overline{A}_{}(t)|^2}}=\text{Im}\lambda _{}\mathrm{sin}(\mathrm{\Delta }mt)\mathrm{sin}2\beta \mathrm{sin}(\mathrm{\Delta }mt)\text{ (SM)},`$ (40)
whereas the corresponding CP asymmetry in the unpolarized rate is much suppressed
$`A_{CP}(t)`$ $``$ $`{\displaystyle \frac{|A_{}(t)|^2+|A_{}(t)|^2|\overline{A}_{}(t)|^2|\overline{A}_{}(t)|^2}{|A_{}(t)|^2+|A_{}(t)|^2+|\overline{A}_{}(t)|^2+|\overline{A}_{}(t)|^2}}`$ (41)
$`=`$ $`{\displaystyle \frac{|A_{}(0)|^2|A_{}(0)|^2}{|A_{}(0)|^2+|A_{}(0)|^2}}\mathrm{sin}(2\beta )\mathrm{sin}(\mathrm{\Delta }mt).`$ (42)
It is interesting that much larger asymmetries are obtained for the coefficients of $`\mathrm{sin}2\varphi `$ and $`\mathrm{cos}2\varphi `$ in the angular dependence (23). Inserting the relations (33), (34) into (23) one finds a particularly simple time-dependence for the CP asymmetry of the $`\mathrm{cos}2\varphi `$ coefficient
$`{\displaystyle \frac{\text{Re}(A_R(t)A_L^{}(t))\text{Re}(\overline{A}_R(t)\overline{A}_L^{}(t))}{|A_R(t)|^2+|A_L(t)|^2+|\overline{A}_R(t)|^2+|\overline{A}_L(t)|^2}}={\displaystyle \frac{1}{2}}\mathrm{sin}2\beta \mathrm{sin}(\mathrm{\Delta }mt).`$ (43)
Expressed in terms of the observed time-dependent angular distributions, this asymmetry can be written as
$`4{\displaystyle \frac{\mathrm{cos}2\varphi \frac{d\mathrm{\Gamma }(t)}{d\varphi }\mathrm{cos}2\varphi \frac{d\overline{\mathrm{\Gamma }}(t)}{d\varphi }}{\frac{d\mathrm{\Gamma }(t)}{d\varphi }+\frac{d\overline{\mathrm{\Gamma }}(t)}{d\varphi }}}=\mathrm{sin}2\beta \mathrm{sin}(\mathrm{\Delta }mt).`$ (44)
We denoted here $`f(\varphi )=_0^{2\pi }𝑑\varphi f(\varphi )`$. Note that the result (44) does not depend on the smallness of the right-handed photon amplitude. On the other hand, the coefficient of $`\mathrm{sin}2\varphi `$ is more sensitive to the presence of a right-handed photon amplitude
$`{\displaystyle \frac{\text{Im}(A_R(t)A_L^{}(t))\text{Im}(\overline{A}_R(t)\overline{A}_L^{}(t))}{|A_R(t)|^2+|A_L(t)|^2+|\overline{A}_R(t)|^2+|\overline{A}_L(t)|^2}}=`$ (45)
$`{\displaystyle \frac{\text{Re}(A_{}A_{}^{})}{|A_{}|^2+|A_{}|^2}}\mathrm{cos}2\beta \mathrm{sin}(\mathrm{\Delta }mt){\displaystyle \frac{\text{Im}(A_{}A_{}^{})}{|A_{}|^2+|A_{}|^2}}\mathrm{cos}(\mathrm{\Delta }mt).`$ (46)
Assuming dominance by the left-handed amplitude in the SM (see Eq. (37)), one can use this asymmetry to extract $`\mathrm{cos}2\beta `$
$`4{\displaystyle \frac{\mathrm{sin}2\varphi \frac{d\mathrm{\Gamma }(t)}{d\varphi }\mathrm{sin}2\varphi \frac{d\overline{\mathrm{\Gamma }}(t)}{d\varphi }}{\frac{d\mathrm{\Gamma }(t)}{d\varphi }+\frac{d\overline{\mathrm{\Gamma }}(t)}{d\varphi }}}=\mathrm{cos}2\beta \mathrm{sin}(\mathrm{\Delta }mt).`$ (47)
While less clean theoretically than the determination of $`\mathrm{sin}2\beta `$ from (44), this result is important because it can help us in resolving discrete ambiguities . In the SM, once $`\mathrm{sin}2\beta `$ is measured, we know $`\beta `$ with no ambiguity from the bounds on the sides of the unitarity triangle. However, in the presence of physics beyond the SM the values of the “would be” $`\beta `$ extracted from asymmetry measurements may not fall within its SM allowed range. Such new physics cannot be detected if the values of the asymmetry (i.e., $`\mathrm{sin}2\beta `$) lie within the SM range. By measuring the sign of $`\mathrm{cos}2\beta `$ we are sensitive to yet another kind of new physics: new CP violating contributions to the mixing amplitude.
## IV Conclusions
In this paper we argued that photon polarization information in exclusive weak radiative $`B`$ decay can be used to probe new physics effects. The SM predicts that the photons emitted in $`\overline{B}(B)`$ decays are almost purely left (right) handed. By measuring the photon polarization we may find a signal for right-handed component that could only be generated by new physics. Moreover, since the linear polarization states are also CP eigenstates, the time-dependent CP asymmetries in $`B^0(t)`$ decays are clean. In the SM they measure $`\mathrm{sin}2\beta `$; by comparing to the CP asymmetries in $`B\psi K_S`$ decay, a possible new CP violating amplitude in the $`bs\gamma `$ decay (independent on its helicity) can be found.
We discuss two methods for determining the photon polarization by using the chain $`BV\gamma Ve^+e^{}`$. In the first method, discussed also in , the photon is off-shell and we need to use the corresponding direct decay $`BVe^+e^{}`$ in the region where the dilepton invariant mass is close to the threshold since there photon exchange dominant the decay. The second method makes use of the Bethe-Heitler process where photon collide with matter and produce a lepton pair.
Two other methods were proposed in the past that are also sensitive to a right-handed component in the radiative decay amplitude. In it was shown that a CP asymmetry in the radiative mode can be generated by right-handed amplitude. In the $`\mathrm{\Lambda }_b\mathrm{\Lambda }\gamma `$ decay was studied and it was shown that the $`\mathrm{\Lambda }`$ polarization is sensitive to the right-handed operator $`\stackrel{~}{O}_7`$.
We comment next on the experimental feasibility of these four methods. Experimentally, each method requires a different analysis, and thus at this stage we can only estimate their relative efficiency. As a benchmark we compare the efficiency of each of these measurement to the efficiency of the $`BK^{}\gamma `$ rate measurement.
First, consider the method based on direct $`BK^{}e^+e^{}`$ decay. Here, the major obstacle is statistics as the rate of the $`BVe^+e^{}`$ decay mode in the $`q^2<1`$ GeV<sup>2</sup> region is smaller by about a factor of $`10^2`$ compared to that of the radiative decay. However, electron pairs produced through virtual photons are most sensitive to the photon polarization (the corresponding efficiency parameter $`\xi =1`$). Moreover, we could hope that the efficiency of the dileptonic mode will be higher than that of the radiative mode. The reconstruction of the dilepton pair emitted near the $`B`$ decay point should be straightforward, as the corresponding tracks are expected to be well separated. In the lab frame, the maximum value of the opening angle between the $`e^+`$ and $`e^{}`$ momenta is $`\mathrm{tan}\theta _{\mathrm{max}}=\sqrt{q^2(2m_e)^2}/|\stackrel{}{q}|`$. For example, at $`q_{\mathrm{max}}^2=0.5`$ GeV<sup>2</sup> the maximum value of this angle is $`\theta _{\mathrm{max}}=15^{}`$. Moreover, at hadron colliders, where the electron pair can be used for triggering, we could hope to get much higher efficiency in the semileptonic mode compared to the radiative mode. Thus, our rough estimates indicate that this decay has an efficiency of the order of few percent compared to the measurement of the $`BK^{}\gamma `$ decay rate.
Next we look at the method using Bethe-Heitler lepton pairs. There are two major drawbacks here. First, the fraction of the photons that are converted is typically of the order of few percent, depending on the detailed matter content of the experimental apparatus. Second, the sensitivity to the photon polarization is not maximal ($`\xi <1`$). On the positive side, we expect that the sensitivity to this mode will be higher compared to that of the $`BK^{}\gamma `$ (where the $`\gamma `$ is identified in the calorimeter) as the energy resolution of the lepton pair is higher and there is less background. The momenta of the conversion electrons produced in the inner layers of the detector are measured very well at CLEO (and the same is expected to be true at BaBar and Belle).
It seems that this kind of measurement is easier to be carried out at $`e^+e^{}`$ machines, as there is much less background. Yet, it is not impossible that also at hadron machines, where the statistics is much larger, this measurement can be done. We may conclude from our rough estimates that the efficiency for the Bethe-Heitler method is also at the few percent level.
The obvious advantage of the method proposed in is that no polarization information is required. However, there are a few factors which offset this advantage. First, it is the fact that time-dependent measurements are necessary. Second, only neutral $`B`$ decays can be used, while in the methods we suggested, also charged $`B`$ (and $`\mathrm{\Lambda }_b`$) can be used. Third, flavor tagging is needed. Last, the final state has to be a CP eigenstate, which gives further suppression of the rate through the chain $`K^0K_S\pi ^0`$. The combined effect of the above factors is a reduction in the efficiency of about a factor of a 100. We can conclude that this method seems somewhat disfavored compared to the methods we described.
Last we estimate the amount of data needed to carry out the suggestion of Ref. . Since this method required $`\mathrm{\Lambda }_b`$ baryons, it is clear that this measurement can be done only at hadron machines. Moreover, $`\mathrm{\Lambda }_b`$ baryons are produced only about $`10\%`$ of the time, and in general are harder to identify than $`B`$ mesons. On the other hand it is relatively easy to collect the polarization information as the $`\mathrm{\Lambda }`$ decay provide it with high efficiency. Again, this very rough estimate suggests that if the radiative decays can be seen at hadron machines, the efficiency for the $`\mathrm{\Lambda }_b\mathrm{\Lambda }\gamma `$ decay with polarization information is at the few percent level compared to the $`BK^{}\gamma `$ rate.
Finally, we comment on the feasibility of the CP asymmetries measurements discussed in Section III. It seems that these measurements are harder to perform since both polarization information and time dependent measurements are needed; thus they suffer from the problem of both our methods and the method of Ref. . Yet, when we try to resolve discrete ambiguities only the sign of $`\mathrm{cos}2\beta `$ is needed. Clearly, the sign of a specific quantity can be determined more easily than its magnitude, and requires less data. Therefore, we could still hope that the large numbers of $`B`$ mesons expected to become available at the hadronic machines would make such measurements feasible.
Clearly, only a detailed experimental analysis can see which method is realistic. According to our estimates, it is possible that all the different analyses discussed above will be carried out.
###### Acknowledgements.
We are grateful to João Silva for helpful comments and in particular for pointing out to us the sensitivity of the time-dependent asymmetry to $`\mathrm{cos}2\beta `$. We thank Helen Quinn for helpful discussions and Stephane Plaszczynski, Soeren Prell, Vivek Sharma and Abner Soffer for useful discussions of the experimental aspects of the methods presented here. Y. G. is supported by the U.S. Department of Energy under contract DE-AC03-76SF00515. The research of D. P. is supported in part by the DOE and by a National Science Foundation Grant no. PHY-9457911.
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# A disk-wind model with correct crossing of all MHD critical surfaces
## 1 Introduction
Astrophysical jets are systematically associated with the presence of an underlying accretion disk, both observationally and theoretically (see Königl & Pudritz 2000 for a recent review). In the case of protostellar objects, accretion disks are resolved by means of infrared and millimeter surveys and interferometric mappings down to scales of a few tens of AU. In the optical and the near infrared, HST high resolution images of disks in several jet sources have also been obtained (Padgett et al. 1999). With an apparent relation found between accretion and ejection in the form of a strong correlation between outflow signatures and accretion diagnostics (see e.g. Cabrit et al. 1990, Cabrit & André 1991, Hartigan et al. 1995), stellar jets seem to be powered by the gravitational energy released in the accretion process.
These facts and considerations have led several authors to develop models of disk winds. The pioneering work of Bardeen & Berger (1978) on a hydrodynamic radially self-similar model of a hot galactic wind was generalized in the seminal paper of Blandford & Payne (1982, henceforth BP82) by including a rotating magnetic field. In particular, in BP82 it was shown that a cold plasma can be launched magneto-centrifugally from a Keplerian disk, similarly to a bead on a wire, provided that the magnetic field lines are sufficiently inclined from the axis. Since then, steady and axisymmetric MHD models, self-similar in the radial direction, have been successfully analyzed and generalized in the literature (see e.g. Contopoulos & Lovelace 1994, henceforth CL94, Li 1995, 1996, Ferreira 1997, Ostriker 1997, Vlahakis & Tsinganos 1998, henceforth VT98, Lery et al. 1999).
A major problem is however still open on the validity of the various classes of radially self-similar solutions analyzed so far. Because, as it is well known since the original work of Weber & Davis (1967) on the rotating magnetized solar wind in the equatorial region, acceptable outflowing solutions must cross smoothly all singularities related to the characteristic speeds of the MHD perturbations, i.e., the poloidal Alfvén velocity and the slow/fast magnetosonic velocities. However, in radially self-similar equations the critical points are not found where the poloidal speed of the flow is equal to the characteristic velocities of these magnetosonic waves. In the cold model of BP82 the “modified” fast magnetosonic critical point (where $`t=1`$ in the BP82 notation) is found downstream of the position where the poloidal velocity of the wind is equal to the fast magnetosonic velocity. Subsequently it has been shown that this is a general property of the axisymmetric steady MHD equations: the singularities of the equations coincide with the positions of the limiting characteristics, or separatrices, within the hyperbolic domain of the governing equations (Bogovalov 1994, Tsinganos et al. 1996). In particular, Bogovalov (1994, 1996) pointed out the key role played by the singularity occurring at the fast magnetosonic separatrix surface (FMSS). Namely, the asymptotic region of the jet is causally disconnected from the base of the flow, only for solutions that cross the critical point at the FMSS. This means that every terminal perturbation or shock does not affect the outflow structure upstream of the position of this critical point. And, Tsinganos et al. (1996) have given several analytical examples where the true singularities of the equations do not coincide with the positions where the governing partial differential MHD equations change character from elliptic to hyperbolic and vice versa. For the sake of simplicity from now on we shall indicate by ‘fast/slow magnetosonic singularity’, or in short ’modified fast/slow’, the critical points at the FMSS/SMSS.
It turns out that in none of the previous models of disk-winds a solution has been found to cross the FMSS. For example, Li (1995, 1996) and Ferreira (1997), starting from the accretion disk, succeeded to cross the slow magnetosonic and the Alfvén ones, but downwind turning points were found where the solutions terminate. Such solutions can be connected to infinity only through a shock, as suggested by Gomez de Castro & Pudritz (1993). However in this case, as the wind velocity is subfast magnetosonic, a temporal evolution of the outflow is expected (Ouyed & Pudritz 1997).
Cylindrically collimated solutions were found by Ostriker (1997) for a cold plasma, integrating the MHD system upstream from infinity and crossing the Alfvén singularity, but always in the subfast magnetosonic regime. On the other hand, it has been shown that in collimated winds oscillations of streamlines are a common feature (Vlahakis & Tsinganos 1997). It thus seems that cylindrically collimated solutions without oscillations correspond to a rather particular choice of parameters that completely suppresses such oscillations. A slight change in these parameters induces the onset of oscillations which increase in amplitude until the configuration is destroyed (Vlahakis 1998). Since the Ostriker (1997) solutions are asymptotically subfast magnetosonic they are likely to be sensitive to perturbations from the external medium, unlike solutions that really satisfy all the criticality conditions. Therefore, such solutions are likely to be structurally and topologically unstable (Vlahakis 1998).
However, it has been shown by Contopoulos (1995) that, in the restricted case of a purely toroidal magnetic field, a smooth crossing of the FMSS is possible. On the other hand in such a case an asymptotically cylindrically collimated configuration is not found; in fact, a new transition to subfast magnetosonic velocities must occur anyway for radially self-similar winds. The only way out is then to match the superfast magnetosonic solution with a shock which is in this case in the physically disconnected domain.
In the present study we extend the analysis of BP82, CL94 and Contopoulos (1995) showing that an exact and simultaneous smooth crossing of all three MHD critical surfaces is possible. In Sec. 2 we define the equations of the hot wind in the framework of a radially self-similar approach and outline the numerical technique. In Sec. 3 we explore the solution topologies in the region around and particularly downstream of the FMSS, where the solution terminates, while in Sec. 4 are shown the features of a few solutions crossing all three critical points with conditions similar to those of BP82. Finally, in Sec. 5 we discuss the possible astrophysical applications of these solutions to stellar jets, and summarize the main implications of our results in comparison with previous ones obtained by other authors.
## 2 Model description
In order to establish notation, in this Section we give a brief derivation of radially self-similar disk-wind models with polytropic thermodynamics. The derivation is along the lines of a systematic method which unifies all self-similar MHD outflows and includes the BP82 model as the simplest case (VT98).
### 2.1 General definitions and self-similar assumption
In steady $`(t=0)`$ and axisymmetric $`(\varphi =0)`$ MHD, the poloidal components of the hydromagnetic field ($`\text{B},\text{V}`$) are defined in terms of the magnetic flux function $`A`$ and mass to magnetic flux function $`\mathrm{\Psi }_A(A)`$ in cylindrical ($`z,\varpi ,\varphi `$) or spherical ($`r,\theta ,\varphi `$) coordinates, as:
$$\text{B}_p=\times \frac{A\widehat{\varphi }}{\varpi },\text{V}_p=\frac{\mathrm{\Psi }_A(A)}{4\pi \rho }\text{B}_p.$$
(1)
The azimuthal components are defined in terms of the total specific angular momentum $`L(A)`$ and of the corotation frequency $`\mathrm{\Omega }(A)`$, which are functions of $`A`$ (Tsinganos 1982):
$$L(A)=\varpi \left(V_\varphi \frac{B_\varphi }{\mathrm{\Psi }_A}\right),\mathrm{\Omega }(A)=\frac{1}{\varpi }\left(V_\varphi M^2\frac{B_\varphi }{\mathrm{\Psi }_A}\right),$$
(2)
and of the poloidal Alfvén number $`M`$ :
$$M=\sqrt{4\pi \rho }\frac{V_p}{B_p}=\frac{\mathrm{\Psi }_A}{\sqrt{4\pi \rho }}.$$
(3)
TransAlfvénic flows require that, when $`M=1`$, $`V_\varphi `$ and $`B_\varphi `$ are finite, i.e.:
$$\frac{L}{\mathrm{\Omega }}=\varpi _\alpha ^2(A)\varpi _{}^2\alpha ,$$
(4)
where $`\varpi _{}`$ is the Alfvén cylindrical radius (the Alfvén lever arm) along the reference field line $`\alpha =1`$, with the dimensionless variable $`\alpha `$ defined as some function of the magnetic flux function $`A`$ which can be reversed to give:
$$A=\frac{B_{}\varpi _{}^2}{2}𝒜\left(\alpha \right).$$
(5)
where $`B_{}`$ is a constant with the dimensions of a magnetic field.
As shown in VT98, all existing classes of radially self-similar MHD solutions can be constructed by making the following two key assumptions:
(i) the Alfvén number $`M`$ is solely a function of $`\theta `$, such that the Alfvén surface is conical:
$$MM(\theta ),$$
(6)
(ii) the cylindrical distance $`\varpi `$ to the polar axis of some fieldline labeled by $`\alpha `$, normalized to its cylindrical distance $`\varpi _\alpha `$ at the Alfvén point is also solely a function of $`\theta `$:
$$G\left(\theta \right)\frac{\varpi }{\varpi _\alpha }.$$
(7)
Following these two assumptions the set of MHD equations is reduced to a system of three ordinary differential equations in $`\theta `$ for $`M(\theta )`$, $`G(\theta )`$ and the $`\theta `$ -dependence of the gas pressure (see VT98 for details).
### 2.2 Polytropic thermodynamics
Depending on the assumptions on the free integrals $`A(\alpha )`$, $`\mathrm{\Psi }_A(\alpha )`$, $`L(\alpha )`$ and $`\mathrm{\Omega }(\alpha )`$, a few classes of radially self-similar solutions exist (see VT98). For only two of these classes a polytropic relationship between the gas pressure and the density is admitted: $`P=Q(\alpha )\rho ^\gamma `$, where $`Q(\alpha )`$ is the specific entropy (the first two cases listed in Table 3 of VT98). In such a case $`A\alpha ^{x/2}`$, $`\mathrm{\Psi }_A\alpha ^{(x3/2)/2}`$, $`\mathrm{\Omega }\alpha ^{3/4}`$, $`L\alpha ^{1/4}`$, and the system of the MHD equations reduces to two first order differential equations for $`M(\theta )`$ and $`G(\theta )`$, supplemented by the Bernoulli integral which also provides the variable $`\psi (\theta )`$, the angle between a particular poloidal fieldline and the cylindrical direction $`\widehat{\varpi }`$ at the spherical angle $`\theta `$. Note that the parameter $`x`$ (with the same notation as in CL94, while in VT98 $`x`$ was denoted by $`F`$) governs the scaling of the magnetic field, while the rotation law is assumed Keplerian. This particular class corresponds to the radially self-similar solutions analyzed in CL94 which contains as a special case the classical BP82 solution with $`x=0.75`$.
The full expressions of $`\mathrm{d}M^2/\mathrm{d}\theta `$, $`\mathrm{d}G^2/\mathrm{d}\theta `$ and $`\psi (\theta `$) are given in the Appendix, Eqs. (A1) - (A3) The expressions for the physical variables become then:
$$\frac{\rho }{\rho _{}}=\alpha ^{x3/2}\frac{1}{M^2},\frac{P}{P_{}}=\alpha ^{x2\gamma \left(x3/2\right)}\left(\frac{\rho }{\rho _{}}\right)^\gamma ,$$
(8)
$$\frac{\text{B}_p}{B_{}}=\alpha ^{\frac{x}{2}1}\frac{1}{G^2}\frac{\mathrm{sin}\theta }{\mathrm{cos}\left(\psi +\theta \right)}\left(\mathrm{sin}\psi \widehat{z}+\mathrm{cos}\psi \widehat{\varpi }\right),$$
$$\frac{\text{V}_p}{V_{}}=\alpha ^{1/4}\frac{M^2}{G^2}\frac{\mathrm{sin}\theta }{\mathrm{cos}\left(\psi +\theta \right)}\left(\mathrm{sin}\psi \widehat{z}+\mathrm{cos}\psi \widehat{\varpi }\right),$$
$$\frac{B_\varphi }{B_{}}=\lambda \alpha ^{\frac{x}{2}1}\frac{1G^2}{G\left(1M^2\right)},$$
(9)
$$\frac{V_\varphi }{V_{}}=\lambda \alpha ^{1/4}\frac{G^2M^2}{G\left(1M^2\right)}.$$
(10)
### 2.3 Parameters
At the Alfvén radius $`\varpi _{}`$ along the reference field line $`\alpha =1`$, we denote by $`P_{}`$ and $`\rho _{}`$ the pressure and density, respectively. The magnitude of the poloidal magnetic field at this Alfvén point is $`B_{}\mathrm{sin}\theta _{}/\mathrm{cos}\left(\psi _{}+\theta _{}\right)`$ while the corresponding poloidal Alfvén speed is $`V_{}\mathrm{sin}\theta _{}/\mathrm{cos}\left(\psi _{}+\theta _{}\right)`$, with $`B_{}=\sqrt{4\pi \rho _{}}V_{}`$.
The expressions of the free integrals defined in Sec. 2.1 can now be written as:
$$A=\frac{B_{}\varpi _{}^2}{x}\alpha ^{x/2},\mathrm{\Psi }_A^2=4\pi \rho _{}\alpha ^{x3/2},$$
(11)
$$\mathrm{\Omega }=\lambda \frac{V_{}}{\varpi _{}}\alpha ^{3/4},L=\lambda V_{}\varpi _{}\alpha ^{1/4},$$
(12)
$$E=V_{}^2ϵ\alpha ^{1/2},V_{}^2=\frac{𝒢}{\varpi _{}\kappa ^2},P_{}=\mu \frac{B_{}^2}{8\pi },$$
(13)
where $`E`$ is the sum of the kinetic, enthalpy, gravitational and Poynting energy flux densities per unit of mass flux density,
$$E(\alpha )=\frac{V^2}{2}+\frac{\gamma }{\gamma 1}\frac{P}{\rho }\frac{𝒢}{r}\frac{\mathrm{\Omega }}{\mathrm{\Psi }_A}r\mathrm{sin}\theta B_\varphi ,$$
(14)
while $`𝒢`$ and $``$ are the gravitational constant and the mass of the central body, respectively.
The solution of the system of Eqs. (A1) - (A3) depends on the six parameters $`x`$, $`\gamma `$, $`\kappa `$, $`\lambda `$, $`ϵ`$ and $`\mu `$, introduced in Eqs. (11) - (13) (but see the discussion in Sect. 2.4.3 for the free parameters of the model). Note that we have used for the parameters a similar but not an identical notation with BP82, since it occurred to us that it is better to choose a different normalization. However, in the following we shall outline for convenience the correspondence between our parameters and those in BP82.
Let us first discuss the physical meaning of the above parameters. First, the exponent $`x`$ is equal to $`3/4`$ in BP82, while in Ferreira (1997) it is related to the ejection index $`\xi =2(x3/4)`$. This index $`\xi `$ is related to the accretion rate and to the mass flux in the wind if also the structure of the disk is assumed radially self-similar (see e.g. Ferreira 1997). Second, we remind that $`\gamma `$ is the usual polytropic index. Next, the constant $`\kappa `$ is the Keplerian speed at radius $`\varpi _{}`$ on the disk, in units of $`V_{}`$, i.e., it is proportional to the ratio of the Keplerian speed to the poloidal flow speed at the Alfvén radial distance, $`V_{\mathrm{p},\mathrm{A}}`$, and is related to the corresponding constant $`\kappa _{\mathrm{BP}}`$ in BP82. Since $`\kappa `$ is also proportional to the mass to magnetic flux ratio, it is often called ’the mass loss parameter’ (Li 1995, Ferreira 1997),
$$\kappa =\sqrt{\frac{𝒢}{\varpi _{}V_{}^2}}=\sqrt{\frac{𝒢}{\varpi _{}V_{\mathrm{p},\mathrm{A}}^2}}\frac{\mathrm{sin}\theta _{}}{\mathrm{cos}\left(\psi _{}+\theta _{}\right)}=\kappa _{\mathrm{BP}}G_o^{3/2}.$$
(15)
The constant $`\lambda `$ is the specific angular momentum of the flow in units of $`V_{}\varpi _{}`$ and is related to the corresponding constant $`\lambda _{\mathrm{BP}}`$ in BP82,
$$\lambda =\frac{L}{V_{}\varpi _{}}=\lambda _{\mathrm{BP}}\kappa \sqrt{G_o}.$$
(16)
The Bernoulli constant $`ϵ`$ is the sum of the enthalpy, kinetic, gravitational and Poynting energy flux densities per unit of mass flux density divided by $`V_{}^2`$ (along $`\alpha =1`$) and is related to the corresponding constant $`ϵ_{\mathrm{BP}}`$ in BP82,
$$ϵ=\frac{E}{V_{}^2}=ϵ_{\mathrm{BP}}\frac{\kappa ^2}{G_o}.$$
(17)
Finally, the constant $`\mu `$ is proportional to the gas entropy,
$$\mu =\mu _{\mathrm{BP}}\left[2G_o^{3\gamma 4}\mathrm{sin}^{2\gamma 2}\psi _o\kappa ^{2\gamma }\right].$$
(18)
In the above expressions, the label $`o`$ indicates the respective values of $`G`$ and $`\psi `$ at the base of the outflow. The correspondence between the parameter $`\xi _o^{}`$ in BP82 and our $`\psi _o`$ is
$$\xi _o^{}=\mathrm{cot}\psi _o.$$
(19)
Note that in BP82 $`\gamma `$ does not appear since the outflow is cold and $`\mu `$, although it is defined, is never used. Also, a similar scaling exists for the parameters used in Li (1995) and Ferreira (1997), although with slightly different notations and a further relation between $`x`$ and $`\kappa `$ (cf. Eq. 28 in Ferreira 1997) due to the connection with a self-similar accretion disk thread by a large scale magnetic field.
### 2.4 Numerical integration
The numerical solution of Eqs. (A1)-(A3) requires the fulfillment of the regularity conditions at the positions of the three singularities (Alfvén and slow/fast modified magnetosonic critical points). This implies that the six parameters of the solution are not all independent. In the following we first shortly summarize the main properties of the critical conditions before we discuss the numerical procedure to obtain the solutions.
#### 2.4.1 Critical Points
It is evident that Eqs. (A2) and (A3) become indeterminate at the Alfvén surface where $`G=M=1`$. The regularity condition at this critical point can be easily found together with the value of the derivative of $`M^2`$ (see Appendix). Furthermore, the denominator of Eq. (A2) vanishes when the meridional component of the velocity $`V_\theta `$ satisfies the quartic (Vlahakis 1998):
$$V_\theta ^4V_\theta ^2\left(C_s^2+V_A^2\right)+C_s^2V_{A,\theta }^2=0,$$
(20)
where $`C_s`$ is the sound speed, and $`V_A`$ and $`V_{A,\theta }`$ the total and meridional components of the Alfvén velocity, respectively.
These singularities are typical ‘X-type’ critical points, and the above equation is the well known dispersion relation for the magnetosonic waves. However it is crucial to see that these singularities appear not when the flow speed, but instead where its meridional component coincides with the meridional component of the slow/fast magnetosonic velocity.
Bogovalov (1994, 1996) and Tsinganos et al. (1996) have emphasized that the singularities in MHD steady flows do not always coincide with the positions where the flow and the magnetosonic velocities coincide, but with the limiting characteristics, i.e., the FMSS and the SMSS. In our case the separatrix is found where $`V_\theta `$ is equal to either one of the triplet of the characteristic speeds ($`V_{s,\theta },V_{A,\theta },V_{f,\theta }`$). This is so because in addition to the azimuthal direction $`\widehat{\varphi }`$ due to the assumed axisymmetry, we have a second symmetry direction, which is the radial direction $`\widehat{r}`$ because of the assumed radial self-similarity. Therefore a compressible slow/fast MHD wave that preserves those two symmetries can only propagate along $`\widehat{\theta }`$ which is perpendicular to both $`\widehat{\varphi }`$ and $`\widehat{r}`$; the speed of propagation of such a wave satisfies exactly the quartic Eq. (20) (for details see Tsinganos et al. 1996).
It is obvious that a physically acceptable solution with low velocity and high density at the base, but high speed and low density asymptotically must smoothly cross at least the SMSS and the Alfvén singularity. Such solutions have been widely analyzed in previous papers and are consistent with the observational data on collimated stellar jets. However it is unescapable that also the fast magnetosonic singularity should be regularly crossed in order to have a steady structure causally disconnected from the asymptotic region, where the jet interacts with the environment (Bogovalov 1994).
#### 2.4.2 Numerical technique for the search of solutions
An inherent difficulty of the problem is due to the fact that the positions of the previous critical points are not known a-priori, but need to be calculated simultaneously and selfconsistently with the sought for solution. At these critical points we do know some relations between various functions, for example, at the Alfvén surface the regularity condition, Eq. (A4), should be satisfied. However, this knowledge alone is not practically enabling us to directly find a solution.
The way we will follow to construct a solution through all critical points is to use the shooting method with successive iterations. By starting the integration from an angle $`\theta =\theta _i`$ we reach a singular point where, e.g., the denominator in $`dM^2/d\theta `$ vanishes, but not the numerator. We then go back and change some parameter and integrate again until it converges, i.e., the denominator and the numerator vanish simultaneously. A similar procedure is followed to cross the other singularity. A rather key point is the choice of the starting position of integration. Most of the previous studies solved the equations by starting from the equator (BP82, CL94), or from infinity (Ostriker 1997). It occurred to us that it is more convenient to integrate the equations starting from the Alfvén critical point, i.e. from the conical surface $`\theta =\theta _{}`$, and move upstream (towards the base) and downstream (towards the external asymptotic region).
For the numerical integration, besides the parameters, we need also to choose the value of the colatitude $`\theta _{}`$ and the value there of the slope of the square of the Alfvén number ($`p_{}=\mathrm{d}M^2/\mathrm{d}\theta |_\theta _{}`$) together with the angle of expansion of the poloidal streamlines ($`\psi _{}`$). Some of these quantities must be tuned to fulfill the singularity conditions at the three critical points. It turned out convenient for the assumed numerical technique to tune the values of $`\lambda `$ and $`p_{}`$ for getting the critical solution.
Hence, we first prescribe the parameters $`\gamma `$, $`x`$, $`\lambda `$ and $`\kappa `$, as well as $`p_{}`$, $`\theta _{}`$ and $`\psi _{}`$ while $`ϵ`$ is deduced from the Bernoulli equation, Eq. (A3), and $`\mu `$ from the regularity condition at the Alfvén point, Eq. (A4). The integration can now start upstream from $`\theta =\theta _{}`$ and the SMSS is encountered, but we cannot pass through it as, e.g., the denominator of $`dM^2/d\theta `$ vanishes there, but not the numerator. We integrate again with different values of $`p_{}`$ until we find the opposite behaviour around the slow magnetosonic singularity (the numerator vanishes but not the denominator). Iteratively, by fine tuning the value of $`p_{}`$, a solution is finally found which pass through the SMSS.
Then we integrate downstream of the Alfvén surface and the FMSS is encountered, but in general it is not crossed. Changing the value of the parameter $`\lambda `$ we integrate upstream again tuning to a new value of $`p_{}`$ until the SMSS is crossed. Then we integrate downstream towards the fast magnetosonic singularity, and repeat all the procedure until we find the right values of $`p_{}`$ and $`\lambda `$ that allow the crossing of the two singularities. At this point the complete solution is obtained by integrating, with the correct values for all the parameters, upstream to the base and downstream towards the asymptotic region.
#### 2.4.3 Selection of the parameters and boundary conditions
In this study, the critical solution depends on the two ‘model’ parameters ($`\gamma `$, $`x`$) and the three independent ‘fieldline’ parameters ($`\kappa `$, $`\theta _{}`$, $`\psi _{}`$). The remaining ones ($`ϵ`$, $`\mu `$, $`p_{}`$, $`\lambda `$) are deduced from the Bernoulli equation and the crossing of the Alfvén, slow magnetosonic and fast magnetosonic singularities, respectively. This is consistent with the analysis of Bogovalov (1997) where it is argued that since the number of equations must be equal to the number of independent boundary conditions, a unique solution can be found if this number of independent boundary conditions equals to the number of outgoing waves generated at the reflection of a plane wave from that boundary. In $`t`$-dependent polytropic MHD there are 7 equations and 7 unknowns: the density, the pressure, the 3 components of the velocity and the 2 components of the magnetic field. There are also 7 waves: the entropy wave and the outwards/inwards propagating slow, Alfvén and fast MHD waves. So, we need 7 parameters with both counts, as expected. Now, if the boundary of the outflow is in the subslow region the number of outgoing waves from this boundary is 4, i.e., the entropy, slow, Alfvén and fast MHD outgoing waves. Subtracting the number of the boundary conditions we are left with 3 independent parameters, precisely $`\kappa `$, $`\theta _{}`$ and $`\psi _{}`$. Note that the polytropic index $`\gamma `$ and $`x`$ should not be included in this count since they are model parameters and do not depend on a particular streamline.
The integration is terminated in the upstream region when $`MM_o<1`$ and $`GG_o<1`$, $`\theta \theta _o\pi /2`$ and $`\psi =\psi _o`$. In all the calculations presented here we were able to follow the solution up to the equator, i.e. $`\theta _o=\pi /2`$. This base should be in principle the disk surface, where our solutions should consistently fit particular boundary conditions. Such an approach has been followed by Ferreira (1997) who looked for inflow/outflow MHD solutions with a consistent matching on the disk surface (see also Li 1995, 1996). This implies some further constraints on the parameters. For instance, if the disk is also self-similar with a large scale magnetic field, a relation between $`\kappa `$ and $`x`$ is expected, i.e., the mass loading in the outflow and the magnetic flux distribution on the disk. In addition, if the outflow carries away all the angular momentum from the disk $`x`$ must be related to $`\lambda `$.
As described above, the procedure to obtain a critical solution is extremely lengthy and rather time consuming. We must in fact approach as close as possible the singularities ($`\mathrm{\Delta }\theta 10^3`$), and this requires the determination of the parameters up to several digits. As we are mainly concerned to analyze the general behaviour of superfast magnetosonic solutions, in the present study we do not investigate the details of the boundaries of the outflow. Therefore we assume that between the base of the wind and the disk surface there is a thin ‘transition’ region that allows the connection of the wind with the disk.
For similar reasons the present analysis has been performed only for a limited set of values of the parameters. We have fixed the ‘fieldline’ parameters $`\theta _{}=59^{}`$, $`\psi _{}=40^{}`$ and $`\kappa =2`$, while two values have been assumed for the ‘model’ parameter $`x`$: 0.75 (model I) and 0.7525 (model II). In this two cases we will assume that the polytropic index $`\gamma <5/3`$, i.e. some amount of heating occurs in the plasma. To make a comparison with a purely magnetocentrifugally driven outflow, we shortly discuss also the very general properties of an adiabatic solution ($`\gamma =5/3`$) assuming $`\theta _{}=60^{}`$, $`\psi _{}=45^{}`$, $`\kappa =3.873`$ and $`x=0.75`$ (i.e., values of $`\theta _{}`$ and $`\psi _{}`$ very close to those used for the nonadiabatic models I and II on purpose have been selected).
In the following two Sections we outline the main properties of the topologies around the FMSS and discuss the structure of the critical solutions.
## 3 Solution Topologies
We present here the topology of two solutions around the fast magnetosonic point, assuming $`\gamma =1.24`$ and $`\gamma =1.23`$ for fixed $`x=0.75`$. The two slightly different values of the polytropic index define the transition between two families of topologies. This drastic change in the topological behavior of the solutions in the neighborhood of the X-type point illustrates the difficulty of exact crossing the fast magnetosonic point. The parameters for the various cases are listed in Tab. 1, while in Fig. 1 we plot the two sets of topologies for the superfast magnetosonic number $`M_{\mathrm{m},\mathrm{f}}(\theta )`$. Note that this plot is obtained from a projection of the solutions from the 3-D space of $`M(\theta )`$, $`G(\theta )`$ and $`\theta `$ to the plane $`M_{\mathrm{m},\mathrm{f}}`$$`\theta `$. This three-dimensional structure of the topology explains why some of the lines obtained by projection are crossing each other (see for another such example Tsinganos & Sauty 1992). This feature of course does not appear in more classical topologies of one-dimensional solutions, e.g., Weber & Davis (1967).
In the first case ($`\gamma =1.24`$) three solutions are plotted in Fig. 1a for different values of $`\lambda `$ and $`\mu `$. The critical solution (solid line in Fig. 1a), moving downstream in the direction of decreasing $`\theta `$ crosses the FMSS at $`\theta _{\mathrm{m},\mathrm{f}}6^{}`$, has a maximum at $`\theta 0.4^{}`$ and then at $`\theta 0.15^{}`$ crosses back the $`M_{\mathrm{m},\mathrm{f}}=1`$ line but with an infinite slope moving towards increasing $`\theta `$. Then, this solution continues marching towards increasing $`\theta `$ and remains always subfast magnetosonic, with $`M_{\mathrm{m},\mathrm{f}}`$ reaching a maximum at $`\theta \theta _{\mathrm{m},\mathrm{f}}`$.
By slightly decreasing $`\lambda `$ the solution crosses the $`M_{\mathrm{m},\mathrm{f}}=1`$ line with infinite slope at $`\theta >\theta _{\mathrm{m},\mathrm{f}}`$ (dashed line in Fig. 1a). Conversely, for a slightly larger value of $`\lambda `$ the solution (dot-dashed line in Fig. 1a) reaches a maximum at $`\theta \theta _{\mathrm{m},\mathrm{f}}`$ remaining subfast magnetosonic (i.e. it behaves like a ‘breeze’ solution) and becomes superfast magnetosonic with diverging slope at $`\theta 0.04^{}`$. Then, this solution remains always in the region $`\theta <\theta _{\mathrm{m},\mathrm{f}}`$, with a spiraling behaviour, i.e., by crossing many times up and down the $`M_{\mathrm{m},\mathrm{f}}=1`$ transition with infinite slope. Note that the solution shown in Fig. 10 of Ferreira (1997) probably belongs to this family of non critical solutions.
For $`\gamma =1.23`$ (Fig. 1b) the topology of the non critical solutions remains the same. The critical solution however shows a different behaviour remaining always in the region $`\theta <\theta _{\mathrm{m},\mathrm{f}}`$ by spiraling around the $`M_{\mathrm{m},\mathrm{f}}=1`$ transition (solid line in Fig. 1b).
The topological structure of our solutions implies that downstream of the FMSS a focal critical point must be present, so that no solution can asymptotically reach $`\theta =0`$ with superfast magnetosonic speeds. This ought to be expected from the construction of this model where we should have $`lim_{\theta 0}V_\theta /V_{\mathrm{f},\theta }=0`$, if we have a cylindrically collimated outflow. In other words, the surface $`M_{\mathrm{m},\mathrm{f}}=1`$ needs to be crossed again with a downstream superfast/subfast magnetosonic transition (see also Contopoulos 1995). At the same time, we should keep in mind that these radially self-similar solutions are not valid to model outflows around the rotational axis, because of their singular behaviour there.
Note that not all solutions with $`M_{\mathrm{m},\mathrm{f}}>1`$ are physically acceptable because they become subfast magnetosonic, crossing the singularity with diverging slope and therefore they are multivalued for the same $`\theta `$. Hence, these solutions could correspond to the terminated solutions in Parker’s terminology for the solar wind with one (Parker 1958), or, multiple critical points (Habbal & Tsinganos 1983). Nevertheless, the present critical solutions are causally disconnected from the inner region of the flow, so that they could be stopped by suitable boundary conditions, e.g. through a shock with the external medium at some angle $`\theta _{min}<\theta _{\mathrm{m},\mathrm{f}}`$ without affecting the structure of the outflow upstream of the FMSS.
It is worth to mention that, from the technical point of view, the main difficulty in obtaining a critical solution is the fact that all solutions (critical ones as well as non critical ones) always reach $`M_{\mathrm{m},\mathrm{f}}=1`$ with infinite slope at some angle $`\theta `$. They become “terminated” at this point, even if they belong to the family of the dot-dashed solution family of Figs. 1. And, both families of non critical solutions almost coincide far from the vicinity of the critical X-point. This is the reason why the crossing of the critical point is so difficult.
## 4 Results
The values of the parameters in the previous Section were chosen such as to illustrate the topology of the solution around the fast critical point. However, they do not correspond to some interesting critical solution from the astrophysical point of view. For example, the fast magnetosonic transition occurs for a rather slow velocity and not far from the Alfvén critical surface. We found that much more interesting results are obtained for a flow closer to isothermal conditions. We then discuss in the following the properties of solutions obtained with $`\gamma =1.05`$ and for two sets of the remaining parameters (models I and II in Tab. 2).
In both cases $`\theta _{}=59^{}`$, $`\psi _{}=40^{}`$ and $`\kappa =2`$, as in the previous topological analysis, with $`x=0.75`$ and $`x=0.7525`$. The remaining parameters are deduced from the requirement to fulfill the criticality conditions and are listed in Tab. 2. We remark that this different choice on the scaling of the magnetic field $`x`$ is important to connect the solution to an accretion disk in the spirit of what has been done by Li (1996) and Ferreira (1997). In such a case, a value of $`x`$ larger than some minimum above the value of BP82, $`x=0.75`$, is necessary to allow ejection (Ferreira & Pelletier 1995). However, it does not mean that our solution fulfills all requirements to connect to such a disk, as we discuss later. The main goal here is to show that the solution is not affected qualitatively by the change in $`x`$ as far as the crossing of critical points is concerned.
In Fig. 2 we plot the various Mach numbers along each field line $`\alpha `$ vs. the vertical height $`z`$ in units of the equatorial cylindrical radius $`\varpi (z=0)`$ of a particular fieldline and for model I. The various critical transitions are indicated, and on the disk surface we find $`G_o0.16`$ and $`M_o0.02`$. The SMSS almost coincides with the point where the flow becomes superslow magnetosonic ($`M_{\mathrm{s},\mathrm{m}}M_\mathrm{s}=1`$), at $`z0.5`$. The Alfvén critical point ($`M=1`$) is crossed at $`z3.5`$ while the wind becomes superfast magnetosonic ($`M_\mathrm{f}=1`$) at $`z20`$. Much farther away is the FMSS, at $`z10^4`$. Downstream of this position the various Alfvén numbers decrease, as expected from the previous topological analysis.
The turning of the solutions is evident in Fig. 3, where the poloidal streamlines together with the characteristics are plotted. They cross all critical surfaces, and for $`\theta <\theta _{\mathrm{m},\mathrm{f}}`$ the fieldlines converge towards the symmetry axis such that the conical region with $`\theta <\theta _{\mathrm{m},\mathrm{f}}`$ is causally disconnected from the rest of the domain. The two families of the characteristics in the hyperbolic domain bounded by the cusp and slow surfaces are better seen in Fig. 4 obtained for the adiabatic case, with a different set of parameters. One family of characteristics (black) is tangent to the SMSS at $`M_{\mathrm{m},\mathrm{s}}=1`$ while the other (grey) crosses it. Similarly, in the hyperbolic domain bounded by the cone where $`M_\mathrm{f}=1`$ one of the family of the characteristics (black) is tangent to the FMSS indicated by $`M_{\mathrm{m},\mathrm{f}}=1`$ while the other (grey) crosses it. We remind that the cusp surface ($`M_\mathrm{c}=1`$) does not coincide with any singularity or typical velocity in the flow.
The components of the outflow speed along a line $`\alpha =const`$ in units of the initial $`z`$-component of the flow speed at the disk, are plotted in Fig. 5. The units are choosen in order to make a direct comparison of this solution with other solutions in the literature (e.g. BP82). Close to the disk level, the escape speed is high, $`V_{\mathrm{esc},\mathrm{o}}440`$, the initial rotational speed is lower, $`V_{\varphi ,\mathrm{o}}=110`$ and of the order of the Keplerian speed, $`V_{\mathrm{Kep}}3V_{\varphi ,\mathrm{o}}`$. The azimuthal speed $`V_\varphi `$ after some increase in the region of corotation, approximately up to the Alfvén critical point at $`z4`$, decays to zero transferring its corresponding kinetic energy to poloidal motion. Thus, the $`z`$\- and $`\varpi `$-components of the poloidal motion grow from their subslow and subescape values at the disk level where $`V_z=1`$ to the high values obtained at the modified fast critical point where $`V_z10^3`$. The poloidal speed exceeds the local escape speed around the Alfvén transition. A comparison of model I and II makes clear that the different values of $`x`$ do not strongly affect the global behaviour of the solutions, even though the boundary conditions of the disk are rather different.
Downstream of the Alfvén transition the azimuthal component of the magnetic field grows to very high values in comparison to the poloidal component (Fig. 6). At the modified fast critical point practically all the magnetic flux is in the azimuthal direction. For example, $`B_\varphi /B_P1`$ at the disk, while $`B_\varphi /B_P60`$ after the modified fast transition for both, models I and II. From Fig. 6 it may be also seen that the flow velocity is largely in the $`z`$-direction with very small components along $`\widehat{\varphi }`$ and $`\widehat{\varpi }`$. In Fig. 6, the main difference between the two models is in the region upstream of the SMSS: for $`x=0.75`$ the angle between the poloidal fieldline and the disk surface is $`\psi _o67^{}`$, while for $`x=0.7525`$ this angle is $`\psi _o56^{}`$. Although these values are not very different, only the second case matches the outflow launching condition for a cold plasma given in BP82. This means that magnetocentrifugal driving is more efficient in model II at the base. However, we note that at the SSMS where the plasma pressure has dropped significantly both solutions can be magneto-centrifugally accelerated. The end result shown in Fig. 5 is that the terminal speed is lower when $`x`$ is larger, i.e., when the ejection index is higher. This result is consistent with Ferreira’s (1997) analysis.
The behaviour of the various components of the conserved total energy $`E`$ vs. $`z`$, plotted in Fig. 7, provides information on the different driving mechanisms that govern the dynamics of the outflow. Upstream of the SMSS and close to the base, most of the energy flux is electromagnetic plus some amount of enthalpy. The kinetic energy of the plasma is negligible. As the slow magnetosonic surface is approached, the kinetic energy sharply increases with a corresponding decrease of the thermal energy. Downstream of the Alfvén surface the Poynting flux rapidly decreases; the poloidal kinetic energy keeps increasing, becoming largely the main component of the energy flux at the position of the FMSS. This behaviour is basically the same for both models I and II. In order words, there is some contribution to the acceleration of thermal origin up to the modified slow critical point after which the pressure drops to a rather constant value while the magnetic pressure maintains considerably higher values up to the Alfvén transition.
We conclude this section by pointing out that the two solutions we have analyzed here correspond to efficient magnetic rotators in the terminology of Bogovalov & Tsinganos (1999), since the ratio of the corotational velocity to the poloidal Alfvén velocity at the Alfvén critical surface (the parameter $`\alpha `$ in their notation) has a value greater than unity ($`2.13`$).
## 5 Discussion
Before discussing the main physical implications of our results, also in connection with those obtained by other authors, we show that the present solutions are suitable to describe the physical properties of astrophysical outflows.
### 5.1 Astrophysical applications
The modeling of a particular astrophysical outflow requires first the calculation of all physical quantities from the non dimensional parameters characterizing the particular model. We will address here this question of calculating some observable quantities of disk-winds associated with protostellar objects from the parameters of our model.
We deduce first the ratios of some characteristic speeds at the disk level, keeping in mind that from the numerical results we have obtained $`M_o0.01`$ and $`G_o0.1`$. We will refer in the following mainly to the solutions with $`x=0.75`$.
First, the ratio of the poloidal Alfvén and Keplerian speeds at the disk level is:
$$\left(\frac{V_{Ap}}{V_{\mathrm{Kep}}}\right)_o=0.316\times \left(\frac{M_o}{0.01}\right)\left(\frac{G_o}{0.1}\right)^{3/2}\frac{1}{\kappa \mathrm{sin}\psi _o}0.178.$$
(21)
The poloidal magnetic field which is essential in the launching of the outflow is anchored in the disk and its energy density is less than the rotational kinetic energy density of the disk. Thus, the field is rather weak to brake the rotation of the plasma at the disk and it is carried passively around by azimuthal rotation.
Second, the ratio of the sound and initial speeds at the disk level is:
$$\left(\frac{C_s}{V_o}\right)_o=70.7\times \left(\gamma \mu \right)^{1/2}\left(\frac{G_o}{0.1}\right)^2\frac{10^4}{M_o^{(\gamma +1)}}94,$$
(22)
where $`V_o=V_z(z=0)`$. The initial ejection speed is negligible in comparison to the thermal speed at the disk, a situation similar to a thermally driven wind.
Next, the ratio of the sound and Keplerian speeds at the disk level is:
$$\left(\frac{C_s}{V_{\mathrm{Kep}}}\right)_o=0.22\times (\gamma \mu )^{1/2}\left(\frac{G_o}{0.1}\right)^{1/2}\frac{1}{\kappa M_o^{(\gamma 1)}}0.314.$$
(23)
We notice that the Keplerian speed is about 3 times higher than the thermal speed at the disk. Thus, thermal effects cannot inhibit the rotation of the disk.
Finally, the ratio of the Keplerian and initial speeds at the disk level is:
$$\left(\frac{V_{\mathrm{Kep}}}{V_o}\right)_o=316\times \kappa \left(\frac{G_o}{0.1}\right)^{3/2}\left(\frac{10^2}{M_o}\right)^2300,$$
(24)
i.e., the initial speed is negligible in comparison to the Keplerian speed.
In our case the flow speed at the fast critical point is about $`10^3`$ the initial speed $`V_o`$. In agreement with the observations we can reasonably assume the terminal speed of the outflow to be $`400`$ km s<sup>-1</sup>, such that its velocity at the base is $`V_o=`$ 0.4 km s<sup>-1</sup>.
In principle, radially self-similar models do not have an intrinsic scale length; however from the previous estimate of the initial speed one allows to calculate the footpoint of the reference fieldline $`\alpha =1`$ on the disk. In units of 10 solar radii this cylindrical distance $`\varpi _o`$ is:
$$\frac{\varpi _o}{10R_{}}=0.19\times \left(\frac{M_o}{0.01}\right)^4\left(\frac{0.1}{G_o}\right)^3\frac{1}{\kappa ^2}\frac{}{_{}}\left(\frac{V_o}{\mathrm{km}\mathrm{s}^1}\right)^2.$$
(25)
Hence, for a one solar mass star we get $`\varpi _o12.5R_{}`$.
It is also interesting to calculate the mass-loss rate $`\dot{}_w`$ in units of $`10^8_{}\mathrm{yr}^1`$ :
$`\begin{array}{c}{\displaystyle \frac{\dot{}_w}{10^8_{}\mathrm{yr}^1}}=0.0386\times \left({\displaystyle \frac{M_o}{0.01}}\right)^2\left({\displaystyle \frac{B_{z,o}}{10G}}\right)^2\times \hfill \\ \\ \left({\displaystyle \frac{\varpi _o}{10R_{}}}\right)^2\left({\displaystyle \frac{V_o}{\mathrm{km}\mathrm{s}^1}}\right)^1f(\alpha _{out},\alpha _{in}),\hfill \end{array}`$ (29)
where
$$f(\alpha _{out},\alpha _{in})=\frac{\alpha _{out}^{x3/4}\alpha _{in}^{x3/4}}{x3/4}\mathrm{if}x0.75,$$
(30)
and
$$f(\alpha _{out},\alpha _{in})=\mathrm{ln}\frac{\alpha _{out}}{\alpha _{in}}\mathrm{if}x=0.75.$$
(31)
By assuming $`\varpi _{in}=\varpi _o`$, $`\varpi _{out}10\varpi _o`$ and $`B_{z,o}=8`$ G we have $`\dot{}_w/(10^8_{}\mathrm{yr}^1)1`$, with a temperature at the base of the flow of:
$$T_{o,in}=3\times 10^5\mu \left(\frac{G_o}{0.1}\right)^4\frac{10^8}{M_o^{2(\gamma +1)}}\left(\frac{V_o}{\mathrm{km}\mathrm{s}^1}\right)^28\times 10^4\mathrm{K},$$
(32)
$$T_{o,out}=T_{o,in}\frac{\varpi _{in}}{\varpi _{out}}8\times 10^3\mathrm{K}.$$
(33)
We remind that $`T_o`$ is not the temperature of the disk as we have assumed a transition layer between the disk surface and the base of the flow (see Sec. 2). This region could be reasonably related to a corona heated by dissipative processes in the plasma (e.g. magnetic reconnection, ohmic heating, etc.; see e.g. Königl & Pudritz 2000).
As the flow corotates roughly up to the Alfvén point (Fig. 5) the specific angular momentum carried by the wind is $`\dot{J}_w=\dot{}_w\mathrm{\Omega }\varpi _\alpha ^2`$ while the angular momentum that has to be extracted locally at the foot point $`\varpi _o`$ of the fieldline in order that the disk accretes at a rate $`\dot{}_a`$ is $`\dot{J}_a=(1/2)\mathrm{\Omega }\varpi _o^2\dot{}_a`$ (Spruit 1996). If the angular momentum carried by the wind is a fraction $`f`$ of $`\dot{J}_a`$ while $`1f`$ is the fraction carried away by viscous stresses, then the ratio of the mass fluxes in the wind and in the accretion flow is
$$\frac{\dot{}_w}{\dot{}_a}=\frac{f}{2}\frac{\varpi _o^2}{\varpi _\alpha ^2}<0.015,$$
taking into account that $`\varpi _\alpha =5.8\varpi _o`$ for model I. It follows that the rate of the outflowing mass is at most of the order of 1$`\%`$ of the rate of the accreted mass; and this is achieved when the wind carries all the angular momentum of the accreted mass. When the outflow carries a smaller fraction of the angular momentum of the disk, the mass loss rate in the wind is an even smaller fraction of the mass loss rate in the wind. In other words, the mass loss rate in the wind is a negligible fraction of the accreted mass, despite that the jet may carry most of the angular momentum of the accreted mass. Similar results are obtained for the case $`x=0.7525`$. Therefore, from the above arguments we may conclude that from our solutions we deduce for the physical parameters values in reasonable agreement with those observed in this class of objects.
Our solution terminates at $`z/\varpi _o2\times 10^4`$, i.e., at $`400`$ Astronomical Units (AU) from the central star. At this position we could argue that there exists a shock matching the solution with the outermost region of the outflow (Gomez de Castro & Pudritz 1993). It is well known that bright knots are observed on scales of thousands AU along most protostellar jets. These configurations are shocks that are interpreted as originated either by fluid instabilities on the jet surface or by temporal variations in the velocity of the outflow (Burrows et al. 1996, Ray 1996, 1998, Micono et al. 1998, Königl & Pudritz 2000). It could be reasonable to associate the terminal shock of our solutions with the inner knots, found at distances down to $`100`$ AU from the star. However we cannot ignore that these knots are non steady configurations and move outwards with velocities $`100÷200`$ km s<sup>-1</sup> (Ray 1996). We could assume that the shock is well upstream of the optical knots: polarimetric radio data on the T Tauri object are consistent with the presence of a shock at $`20÷40`$ AU from the star (Ray et al. 1997). Alternatively the terminal shock could indeed be located approximately at the positions of the inner knots, but there the flow looses both self-similarity and steadiness. However as we are in the superfast magnetosonic regime, the upwind configuration is not affected. Only a much more detailed parametric study will be able to test these two possibilities.
### 5.2 Physical properties of the critical solutions
The solutions of this model, in particular Fig. 7, illustrate nicely the physical process of transferring electromagnetic Poynting energy flux and enthalpy to directed kinetic energy flux of the flow in order to accelerate a disk wind and then form a jet along the symmetry axis of the system. Thus, the analysis of the previous section is interesting in the sense that it reveals the driving mechanisms of the outflow. The poloidal kinetic energy is negligible at the disk level. It then increases rather sharply up to the region of the SMSS and Alfvén surfaces. This increase is at the expense of both, the enthalpy and the electromagnetic Poynting energy flux (see, Fig. 7). The poloidal velocity is directed basically in the radial direction (Figs. 5 and 6), i.e., here part of the random thermal energy together with a part of the electromagnetic energy are mostly transformed to directed wind expansion. Downstream of the Alfvén surface it is mainly the Poynting energy flux that is effectively transformed into kinetic energy directed along the rotational axis, till the FMSS is encountered. After the FMSS, the flow has already reached the maximum speed available from the total energy E, which is also approximately equal to the initial electromagnetic Poynting energy flux. Then, the acceleration asymptotically stops. Despite the fact that most of the acceleration to high speeds is apparently of magnetic origin, the role of the polytropic index and thus of the initial thermal acceleration may not be negligible, in particular in the region before the SMSS. For example, in the case of Fig. 4 where the flow is exactly adiabatic and $`\gamma =5/3`$, the critical solution achieves only a very small axial component of the velocity which is twice the axial velocity on the equatorial plane. In the quasi-isothermal case of models I and II where $`\gamma =1.05`$, the maximum velocity is 1000, higher than the equatorial one (Fig. 5). As a matter of fact, this last case is closer to the one analyzed in Li (1995) and Ferreira (1997) where the gas is isothermal up to the first critical surface and then it is taken to be cold afterwards, wherein the pressure has sufficiently dropped. However, another possibility is that the low terminal speed obtained in the adiabatic case of Fig. 4 could be due to the lower value of the rotation parameter $`\lambda ^2`$ which is $`2.8`$ in the adiabatic case of Fig. 4, as opposed to values $`137`$ and 136 in models I and II and similarly for the case examined in Ferreira (1997).
When the gas has reached a high speed along the $`z`$-axis, its inertia causes it to lag behind the rotation of the field line and the field is wound up, as shown in Fig. 6, resulting to a highly twisted magnetic field. Consequently, the strong curvature force of this predominantly azimuthal magnetic field towards the $`z`$-axis, causes the poloidal field to collimate. Initially the field is flaring away from the rotation axis but the curvature force bends the poloidal field lines toward the rotation axis. The azimuthal velocity peaks around the Alfvén point which is at a height $`z=3.5`$ and a cylindrical distance $`\varpi =5.8`$ times the starting distance $`\varpi _o`$ in model I. Beyond the Alfvén point the rotation drops in accordance to angular momentum conservation and thus the centrifugal force becomes negligible. Then, the strong inwards curvature force of the twisted field, wins, over the weak outwards centrifugal force and gas pressure gradient with the result that the lines are bent and eventually collapse towards the rotation axis.
It is interesting that this feature of the collapse of the outflow towards the rotation axis which appears in cold models (BP82) and models that do not cross the FMSS (Li 1995, Ferreira 1997), is also preserved in the present hot model where also all critical points are crossed. This result seems to indicate the rather dominant role of the magnetic hoop stress in radially self-similar models, contrary to what happens in meridionally self-similar models wherein the structure becomes asymptotically cylindrical (Trussoni et al. 1997, Sauty et al. 1999, Vlahakis & Tsinganos 1999).
It is worth to clarify for a moment the term “disk-wind” that we used in this study. By that term we simply intend to indicate that we describe an outflow from a disk-like structure accreting onto a central gravitational object. Thus, the flow starts at some angle $`\theta _o`$ above or at the equatorial plane of the disk, as opposed to a “stellar” wind flow that starts radially above or at a spherical or quasi-spherical source. Needless to say that a consistent solution of the accreting part of the flow would be required for a consistent solution of the inflow-outflow structure in the case of a disk-wind. However, such a complete undertaking is beyond the scope of the present paper which only intends to emphasize the possibility to construct complete steady self-similar solutions for the wind crossing all critical points.
To make such a connection between the disk and the outflow, in the spirit of BP82, Li (1995) and Ferreira (1997), the first step would be to see how our parameters may fall into the range of parameters considered by those models. For that purpose, in Eqs. (15) - (19) we have made a correspondance between our parameters and those used by BP82. Thus, in the “standard” solution of BP82 the parameters are: $`\kappa _{\mathrm{BP}}=0.03`$, $`\lambda _{\mathrm{BP}}=30`$ and $`\xi _o^{}=1.58`$ corresponding to a launching angle of the jet at the disk $`\psi _o32^{}<60^{}`$. In our case, we find $`\kappa _{\mathrm{BP}}0.13`$, $`\lambda _{\mathrm{BP}}14.57`$, for both, model I and model II. We also have $`\xi _{o}^{}{}_{\mathrm{BP}}{}^{}=\mathrm{cot}\psi _o=0.425`$ ($`\psi _o=67^{}`$) for model I and $`\xi _{o}^{}{}_{\mathrm{BP}}{}^{}=\mathrm{cot}\psi _o=0.675`$ ($`\psi _o=56^{}`$) for model II, in the BP82 notation. We note that the values of $`\kappa _{\mathrm{BP}},\lambda _{\mathrm{BP}}`$ are close in BP82 and the present model. However, the value of the launching angle $`\psi _o`$ is $`>60^{}`$ in our model I because of the additional thermal driving of the outflow at the disk level, contrary to the cold model of BP82 where $`\psi _o32^{}<60^{}`$. In summary, our models I and II occupy in the space of $`\kappa _{\mathrm{BP}}`$ and $`\lambda _{\mathrm{BP}}`$, roughly the same domain as in BP82 (cf. Fig. 2 in BP82). The only difference is in the value of the launching angle $`\psi _o`$ which can be larger in the present hot model, as expected. These values are within the range of the allowed parameters in the ($`\kappa _{\mathrm{BP}}`$, $`\lambda _{\mathrm{BP}}`$) space also in the analysis of Li (1995, cf. Fig. 3) provided that the magnetic diffusivity is of order one. Note also that model II with $`x=0.7525`$ corresponds to an ejection index in the notation of Ferreira (1997) $`\xi =2x3/2=0.005`$.
### 5.3 Summary
In this paper we have extended the classical work of Blandford and Payne (1982), mainly by showing via examples for the first time that a solution passing through all MHD critical points can indeed be constructed.
As is well known, the FMSS plays the role of the MHD signal horizon such that in an outflow crossing this MHD horizon all perturbations which the outflow may encounter are convected downstream by the superfast outflow and so the steady state solution is maintained. In other words, the outflow interior to the FMSS is causally disconnected and protected against any conditions it may encounter in the interstellar or intergalactic medium towards which the jet propagates after it is launched by magnetocentrifugal forces from the surface of an accretion disk.
Unlike other analytical models which produce asymptotically cylindrically collimated outflows (Sauty & Tsinganos 1994, Trussoni et al. 1997, VT98, Sauty et al. 1999, Vlahakis & Tsinganos 1999), this class of radially self-similar models cannot continue to infinity but it has to be stopped downstream of the FMSS and matched via a MHD shock to a subfast outflow that mixes with the interstellar medium (Gomez de Castro & Pudritz 1993). This shock can connect the present solutions to some breeze, subAlfvén or subslow magnetosonic branch perhaps also preserving the self-similarity.
Thus, the main difference here with previous results presented in the literature is that the asymptotic part of the present solutions is causally disconnected from the source and hence any perturbation downstream of the superfast transition cannot affect the whole structure of the steady outflow.
This task of matching the present solutions with a downstream shock however remains a challenge for future studies, together with a (time-consuming) more extended parametric analysis and also a correct matching of the ideal MHD outflow solutions with an inflow in a non-ideal accretion disk (Ferreira 1997).
## Acknowledgments
This research has been supported in part by a bilateral agreement between Greece and France (program Platon) and a NATO collaborative research grant between Greece, Italy and Russia. E.T. acknowledges the hospitality of the Observatoire de Paris and of the Department of Physics of the University of Crete. We wish to thank M. Micono for the information on the last data on protostellar jets and an anonymous referee for his comments which resulted in a better presentation of the paper.
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## Appendix
The two first order differential equations for $`G(\theta )`$, $`M(\theta )`$ governing the present class of solutions are:
$$\frac{dG^2}{d\theta }=\frac{2G^2\mathrm{cos}\psi }{\mathrm{sin}\theta \mathrm{cos}\left(\psi +\theta \right)},$$
(34)
$`\begin{array}{c}{\displaystyle \frac{dM^2}{d\theta }}=2{\displaystyle \frac{\mathrm{sin}\left(\psi +\theta \right)}{\mathrm{cos}\left(\psi +\theta \right)}}\{{\displaystyle \frac{\kappa ^2\mathrm{sin}\theta }{G}}\mu (x2)M^{42\gamma }+\hfill \\ \\ {\displaystyle \frac{M^4}{G^4}}\left(1M^2\right){\displaystyle \frac{\mathrm{cos}\psi \mathrm{sin}\theta }{\mathrm{sin}\left(\psi +\theta \right)}}{\displaystyle \frac{M^4}{G^4}}\left(x2\right){\displaystyle \frac{\mathrm{sin}^2\theta }{\mathrm{cos}^2\left(\psi +\theta \right)}}\hfill \\ \\ \lambda ^2{\displaystyle \frac{M^4}{G^2}}\left(x2\right)\left({\displaystyle \frac{1G^2}{1M^2}}\right)^2+\lambda ^2{\displaystyle \frac{M^2}{G^2}}{\displaystyle \frac{G^4M^2}{1M^2}}\hfill \\ \\ \lambda ^2{\displaystyle \frac{\mathrm{cos}\psi }{\mathrm{sin}\theta \mathrm{sin}\left(\psi +\theta \right)}}{\displaystyle \frac{\left(2M^21\right)G^4M^4}{G^2\left(1M^2\right)}}\}\times \hfill \\ \\ \{\gamma \mu (1M^2)M^{2\gamma }2\lambda ^2{\displaystyle \frac{M^2}{G^2}}\left({\displaystyle \frac{1G^2}{1M^2}}\right)^2+\hfill \\ \\ 2{\displaystyle \frac{M^4\mathrm{sin}^2\theta }{G^4}}(1{\displaystyle \frac{1}{M^2\mathrm{cos}^2\left(\psi +\theta \right)}})\}^1.\hfill \end{array}`$ (46)
In the above two equations $`\psi (\theta )`$ is given by the Bernoulli integral:
$`\begin{array}{c}\psi =\pi \theta \mathrm{arctan}\{{\displaystyle \frac{G^4}{M^4\mathrm{sin}^2\theta }}[2ϵ{\displaystyle \frac{\gamma \mu }{\left(\gamma 1\right)M^{2\left(\gamma 1\right)}}}+\hfill \\ \\ {\displaystyle \frac{2\kappa ^2\mathrm{sin}\theta }{G}}\lambda ^2({\displaystyle \frac{\left(G^2M^2\right)^2}{G^2\left(1M^2\right)^2}}+2{\displaystyle \frac{1G^2}{1M^2}})]1\}^{1/2}.\hfill \end{array}`$ (50)
with the upper sign corresponding to the outflow case ($`V_r>0`$).
On the Alfvén conical surface for $`\theta \theta _{}`$ we have
$$\left(\frac{1G^2}{1M^2}\right)_{}=\frac{2\mathrm{cos}\psi _{}}{p_{}\mathrm{sin}\theta _{}\mathrm{cos}\left(\psi _{}+\theta _{}\right)},$$
where $`p_{}`$ is the slope of the square of the Alfvén number. Then from Eq. (A2) we get the following third degree polynomial for $`p_{}`$:
$`\begin{array}{c}\left(x2\right)\left(4\lambda ^2+p_{}^2\mathrm{sin}^2\theta _{}\right)\mathrm{tan}^2\left(\psi _{}+\theta _{}\right)+\hfill \\ \\ \left(p_{}^3\mathrm{sin}^2\theta _{}+4\lambda ^2p_{}+8\lambda ^2{\displaystyle \frac{\left(x2\right)}{\mathrm{tan}\theta _{}}}\right)\mathrm{tan}\left(\psi _{}+\theta _{}\right)+\hfill \\ \\ \left(x2\right)\left(\mu p_{}^2+p_{}^2\mathrm{sin}^2\theta _{}+4\lambda ^2{\displaystyle \frac{1}{\mathrm{tan}^2\theta _{}}}\right)+\kappa ^2p_{}^2\mathrm{sin}\theta _{}\hfill \\ \\ \lambda ^2p_{}\left(p_{}{\displaystyle \frac{4}{\mathrm{tan}\theta _{}}}\right)=0.\hfill \end{array}`$ (58)
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# Scalar Gluonium and InstantonsHD-TVP-00-1
## I Introduction
Glueballs, as the most immediate manifestation of gluonic self-interactions in the hadron spectrum, represent a key challenge for our understanding of nonperturbative Yang-Mills dynamics. Not surprisingly, therefore, theoretical interest in gluonium dates back to the early days of QCD and has spurred intense research activity ever since. Estimates of glueball properties were obtained in a variety of approaches, ranging from model analyses and steadily improving lattice simulations to QCD sum rule calculations .
The QCD sum-rule technique, in particular, combines the advantages of an analytical approach with a firm and largely model-independent basis in QCD and should therefore be well suited for obtaining qualitative and quantitative insight into the glueball spectrum. During its early development, however, it became clear that this approach encounters conceptual and practical problems in the scalar glueball channel . These problems, which have so far prevented fully consistent and conclusive sum-rule predictions, can be traced to the exceptionally strong coupling of the vacuum to the spin-0 glueball interpolators. The intense vacuum response generates nonperturbative violations of asymptotic freedom starting from unusually small distances $`x0.02`$ fm , and the power corrections of the conventional operator product expansion (OPE) are much too weak to account for this physics . As a consequence, stability and mutual consistency of different sum rules (see below) is partially lacking, and serious difficulties are encountered in reconciling the sum rules with the low-energy theorem which governs the long-distance behavior of the scalar glueball correlator.
In search of the origin for the apparently missing nonperturbative short-distance physics it is natural to expect that it may at least partly be provided by small (or “direct”) instantons . Indeed, as coherent vacuum gluon fields instantons couple particularly strongly to the gluonic interpolators of the $`0^{++}`$ glueball channel. Nevertheless, they are neglected in the perturbatively calculated Wilson coefficients of the conventional OPE. An additional and perhaps more intuitive reason for expecting instanton physics to play a major role in the structure of scalar glueballs derives from their exceptionally small size $`r_G0.2`$ fm, found in lattice and instanton vacuum model calculations. Since $`r_G`$ is much smaller than the confinement scale and the size of heavier glueballs, it is rather unlikely that scalar glueballs are bound predominantly by (iterated) perturbative or by confinement forces. The attractive interactions mediated by instantons provide a suggestive alternative.
Motivated by the above considerations, our main objectives in this paper are to evaluate the direct instanton sector of the scalar glueball correlator by means of an instanton-improved OPE (IOPE) and to analyze the ensuing glueball sum rules. While analytical instanton calculations (at low energies) have long been hampered by insufficient knowledge of the instanton size distribution and notorious infrared problems, this impasse can nowadays be avoided by relying on the results of instanton vacuum model and lattice simulations for the required bulk features of the instanton distribution, i.e., for the average instanton size $`\overline{\rho }(1/3)`$ fm and density $`\overline{n}(1/2)`$ fm<sup>-4</sup>. Despite remaining numerical discrepancies pertaining to the distribution of large-size instantons, these scales provide a solid foundation for our calculations below to which only small instantons of sizes $`\rho `$ 0.5 fm will contribute.
## II IOPE and sum rules
Our study will be based on the correlation function
$$\mathrm{\Pi }\left(q^2\right)=id^4xe^{iqx}0|TO_S\left(x\right)O_S\left(0\right)|0$$
(1)
with the interpolating field
$$O_S=\alpha _sG_{\mu \nu }^aG^{a,\mu \nu }$$
(2)
carrying the quantum numbers of the scalar glueball. The standard OPE of this correlator, including perturbative Wilson coefficients up to $`O\left(\alpha _s\right)`$ and operator contributions up to dimension 8, is known to be (up to polynomials in $`Q^2`$)
$`\mathrm{\Pi }^{\left(OPE\right)}(Q^2)`$ $`=Q^4\mathrm{ln}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right)\left[2\left({\displaystyle \frac{\alpha _s}{\pi }}\right)^2\left[1+{\displaystyle \frac{59}{4}}{\displaystyle \frac{\alpha _s}{\pi }}\right]+b_0\left({\displaystyle \frac{\alpha _s}{\pi }}\right)^3\mathrm{ln}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right)\right]`$ (3)
$`+4\alpha _s\left[1+{\displaystyle \frac{49}{12}}{\displaystyle \frac{\alpha _s}{\pi }}\right]\alpha _sG^2{\displaystyle \frac{\alpha _s^2b_0}{\pi }}\alpha _sG^2\mathrm{ln}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right)`$ (4)
$`+{\displaystyle \frac{1}{Q^2}}\left[8\alpha _s^2gG^3L^{7/11}58\alpha _s^3gG^3L^{7/11}\mathrm{ln}\left({\displaystyle \frac{Q^2}{\mu ^2}}\right)\right]+8\pi \alpha _s{\displaystyle \frac{1}{Q^4}}\alpha _s^2G^4`$ (5)
($`Q^2=q^2`$) where $`b_0=11N_c/32N_f/3`$ is the leading-order contribution to the QCD $`\beta `$ function, $`L\left(Q^2\right)=\mathrm{ln}(Q/\mathrm{\Lambda })/`$ $`\mathrm{ln}\left(\mu _0/\mathrm{\Lambda }\right)`$, and $`\alpha _s`$ is the QCD running coupling at one loop. We will use the parameter values $`\mathrm{\Lambda }=0.12`$ GeV, $`\mu =0.5`$ GeV and the condensate values $`\alpha _sG^2=0.04`$ GeV<sup>4</sup>, $`gG^3=1.5\alpha _sG^2^{3/2}`$ of Ref. . For the four-gluon condensate, finally, we adopt the standard approximation
$$\alpha _s^2G^414(\alpha _sf_{abc}G_{\mu \rho }^bG_\nu ^{\rho c})^2(\alpha _sf_{abc}G_{\mu \nu }^bG_{\rho \lambda }^c)^2\frac{9}{16}\alpha _sG^2^2.$$
(6)
In addition to the perturbative contributions given above, the Wilson coefficients receive nonperturbative contributions from direct instantons which have so far been neglected in the gluonium sum rules (a partial estimate was given in Ref. ). Analogous contributions were found to be important in several nucleon and pion sum rules. As noted there, effects of multi-instanton correlations can be neglected in the short-distance expansion since the relevant distances $`\left|x\right|\left|Q^1\right|0.2`$ fm are much smaller than the average separation $`\overline{R}1`$ fm between instantons in the vacuum. The diluteness of the instanton vacuum distribution, which is a consequence of $`\overline{\rho }/\overline{R}1`$, further reduces the impact of multi-instanton correlations and keeps the separate instantons approximately undeformed.
To leading order in the semiclassical approximation, the instanton contribution to Eq. (1) can thus be calculated by standard techniques from the $`O\left(\mathrm{}^0\right)`$ component of the gluon propagator in the instanton background and reads
$$\mathrm{\Pi }^{\left(I+\overline{I}\right)}\left(Q^2\right)=2^5\pi ^2\overline{n}\overline{\rho }^4Q^4K_2^2\left(Q\overline{\rho }\right).$$
(7)
($`K_2`$ is a McDonald function.) Since the average instanton size $`\overline{\rho }\left(1/3\right)`$ fm is small compared to $`\mathrm{\Lambda }_{QCD}^1`$, $`O\left(\mathrm{}\right)`$ corrections to Eq. (7) are suppressed by the large instanton action $`S_I\left(\overline{\rho }\right)10\mathrm{}`$. Instanton contributions to the Wilson coefficients of power corrections carry additional inverse powers of the relatively large glueball mass scale and are therefore also expected to be small (see Ref. for a more detailed discussion).
Various sum rules can be constructed from the Borel transform of weighted moments of the glueball correlator,
$$_k\left(\tau \right)=\widehat{B}\left[\left(Q^2\right)^k\mathrm{\Pi }\left(Q^2\right)\right].$$
(8)
Typically, one considers $`k\{1,0,1,2\}`$. The corresponding expressions for $`_k^{\left(OPE\right)}`$ (together with the explicit form of the Borel operator $`\widehat{B}`$) are given in Ref. . The instanton contributions $`_k^{\left(I+\overline{I}\right)}`$ are obtained recursively, via
$$_k^{\left(I+\overline{I}\right)}\left(\tau \right)=\left(\frac{}{\tau }\right)^{k+1}_1^{\left(I+\overline{I}\right)}\left(\tau \right)\text{ }\left(k1\right)$$
(9)
from $`_1^{\left(I+\overline{I}\right)}`$, which can be calculated in closed form \[$`x\overline{\rho }^2/\left(2\tau \right)`$\]:
$$_1^{\left(I+\overline{I}\right)}\left(\tau \right)=2^6\pi ^2\overline{n}x^2e^x\left[\left(1+x\right)K_0\left(x\right)+\left(2+x+\frac{2}{x}\right)K_1\left(x\right)\right]+2^7\pi ^2\overline{n}.$$
(10)
The sum $`_k^{\left(OPE\right)}+_k^{\left(I+\overline{I}\right)}`$ constitutes the IOPE. Note that we have removed the constant subtraction term $`\mathrm{\Pi }^{\left(I+\overline{I}\right)}\left(0\right)=2^7\pi ^2\overline{n}`$ in Eq. (10) because it originates from soft instanton contributions which do not belong to the OPE coefficients. Double-counting of soft instanton physics is thereby excluded since the instanton contributions (9) do not contain powers $`\tau ^n`$ with $`n>\left(k+3\right)`$.
In order to write down the sum rules, we have to match the IOPE expressions to their “phenomenological” counterparts, which are derived from the twice subtracted dispersion relation
$$\mathrm{\Pi }^{\left(phen\right)}\left(Q^2\right)=\mathrm{\Pi }^{\left(phen\right)}\left(0\right)\mathrm{\Pi }^{\left(phen\right)^{}}\left(0\right)Q^2+\frac{\left(Q^2\right)^2}{\pi }_0^{\mathrm{}}𝑑s\frac{Im\mathrm{\Pi }^{\left(phen\right)}\left(s\right)}{s^2\left(s+Q^2\right)}$$
(11)
by parametrizing the spectral function in terms of a pole contribution and an effective continuum. Following standard procedure, the latter is obtained from the dispersive cut of the theoretical side and starts at an effective threshold $`s_0`$. Thus we have
$$Im\mathrm{\Pi }^{\left(phen\right)}\left(s\right)=\pi f_G^2m_G^4\delta \left(sm_G^2\right)+\left[Im\mathrm{\Pi }^{\left(OPE\right)}\left(s\right)+Im\mathrm{\Pi }^{\left(I+\overline{I}\right)}\left(s\right)\right]\theta \left(ss_0\right).$$
(12)
As a consequence of the exceptional size of the instanton contributions to the scalar glueball correlator, their contributions to the continuum are an indispensable part of Eq. (12) and will turn out to play an essential role in the subsequent analysis. Explicitly, we find
$$Im\mathrm{\Pi }^{\left(I+\overline{I}\right)}\left(s\right)=2^4\pi ^4\overline{n}\overline{\rho }^4s^2J_2\left(\sqrt{s}\overline{\rho }\right)Y_2\left(\sqrt{s}\overline{\rho }\right)$$
(13)
where the $`J_2`$ ($`Y_2`$) are Bessel (Neumann) functions. (A more detailed multipole analysis, allowing for neighboring and mixed quarkonium resonances, will be relegated to Ref. .)
By equating the phenomenological Borel moments, obtained from Eqs. (8) and (11), to the corresponding IOPE expressions one finally obtains the sum rules
$$\frac{_k(\tau ,s_0)}{m_G^{2+2k}}=f_G^2m_G^2e^{\tau m_G^2}$$
(14)
with
$$_k(\tau ,s_0)=\underset{X=OPE,I+\overline{I}}{}\left[_k^{\left(X\right)}\left(\tau \right)_k^{\left(Xcont\right)}(\tau ,s_0)\right]+\delta _{k,1}\mathrm{\Pi }^{\left(phen\right)}\left(0\right)$$
(15)
and
$$_k^{\left(Xcont\right)}(\tau ,s_0)=\frac{1}{\pi }_{s_0}^{\mathrm{}}𝑑ss^kIm\mathrm{\Pi }^{\left(X\right)}\left(s\right)e^{s\tau }.$$
(16)
Note that the higher moments weight the higher-mass region of the spectral function more strongly and thus receive enhanced contributions from the relatively heavy (see below) glueball pole. The subtraction constant $`\mathrm{\Pi }^{\left(phen\right)}\left(0\right)`$ in the $`_1`$ sum rule (regularized by removing the high-momentum contributions) can be related to the gluon condensate by the low-energy theorem (LET)
$$\mathrm{\Pi }\left(0\right)=\frac{32\pi }{b_0}\alpha G^2.$$
(17)
This relation provides an important consistency check for the sum-rule results, as we will discuss below.
## III Sum rule analysis
The quantitative analysis of the sum rules amounts to determining those values of the hadronic parameters in Eq. (12) for which both sides of Eq. (14) optimally match in the fiducial Borel domain. Towards large $`\tau `$ this domain is bounded by keeping the contribution of the highest-dimensional operator ($`\alpha _s^2G^4`$) to $`_k`$ below 10% and requiring multi-instanton contributions to be negligible. The latter requirement will be (conservatively) implemented by demanding $`\tau 1`$ GeV<sup>-2</sup>. Towards small $`\tau `$ we prescribe that the continuum contributions do not exceed 50% of the $`_k`$.
The standard optimization procedure followed in previous analyses determined only the glueball mass $`m_G`$ and coupling $`f_G`$ from matching the sum rules, while the threshold $`s_0`$ had to be found by other means (e.g. by finite-energy sum rules or stability criteria ). The IOPE sum rules turn out to be stable enough, however, to determine $`s_0`$ together with the resonance parameters $`m_G`$ and $`f_G`$ from the same sum rule. This is the procedure which we will adopt below.
We start by analyzing the $`_0`$ sum rule \[i.e., Eq. (14) with $`k=0`$\]. Figure 1 shows both sides of the optimized sum rule and separately the three components (OPE with subtracted OPE continuum, instanton contribution, and instanton continuum) which make up its left-hand side. The matching between both sides of the sum rule is almost perfect over the whole fiducial region. Comparing standard OPE and instanton (including continuum) contributions shows that the latter are about 5 times larger. Thus, the instanton contributions strongly dominate over the whole fiducial region and increase the predictions for $`f_G^2`$ by about a factor of 5, resulting in $`f_G=1.14`$ GeV. A similarly strong enhancement of $`f_G`$ was found in instanton vacuum model calculations . The prediction for the glueball mass, $`m_G=1.40`$ GeV, on the other hand, differs surprisingly little from what is obtained without the instanton part. The continuum threshold becomes $`s_0=5.1`$ GeV<sup>2</sup>.
Figure 1 furthermore reveals that the instanton contributions to the unitarity cut are indispensable for generating an exponential $`\tau `$ behavior below $`\tau 0.8`$ GeV<sup>-2</sup>, and thus for an acceptable fit to the pole contribution. (For larger values of $`\tau `$ the continuum contributions are practically negligible and a fit to the instanton contribution alone would become possible, although mostly outside of the fiducial domain and with about 20% smaller values for $`m_G`$ and $`f_G^2`$.) It should also be noted that the hard nonperturbative instanton physics begins to enter $`_0`$ at much smaller $`\tau `$ than the soft condensate contributions, thereby confirming the existence of an exceptionally large mass scale in the scalar glueball channel .
The analysis of the remaining sum rules (which all show a high degree of stability) confirms the above observations about the role of the instanton contributions. As an additional example, we plot in Fig. 2 the separate contributions to the $`_2`$ sum rule and the fit of both sides, which is again excellent. Since the instanton contribution is somewhat less pronounced than in $`_0`$, we find a smaller value for the coupling: $`f_G=1.01`$ GeV. The result for the glueball mass increases to $`m_G=1.53`$ GeV, and the threshold $`s_0=4.89`$ GeV<sup>2</sup> is slightly reduced. The results of the $`_1`$ and $`_1`$ sum rules confirm the tendency of lower moments to predict somewhat smaller masses and somewhat larger couplings \[while maintaining consistency with the low-energy theorem (17)\]. The predictions of the higher moments should be more reliable, however, because they receive stronger pole contributions.
The $`_1`$ sum rule has played both conceptually and practically a special role since it contains the subtraction constant $`\mathrm{\Pi }\left(0\right)`$ which dominates the power corrections of the conventional OPE (in the fiducial region). Attempts to fit the resulting, almost flat $`\tau `$ behavior to the exponential pole contribution inevitably generate very small pole masses, at least half an order of magnitude smaller than those predicted by the other sum rules. This well-known inconsistency (and the need to abandon the $`_1`$ sum rule in practice) is largely resolved by the massive instanton contributions. Their strong decay yields excellent fits and pole masses of the same order as those obtained from the higher moments. Still, the $`_1`$ sum rule remains probably least suited for quantitative predictions since it is most sensitive to the inaccurately known value of the gluon condensate and least sensitive to the pole contribution. The mutual agreement of all four IOPE sum rules (in the typical range of uncertainty for sum rule results) and their consistency with the low-energy theorem<sup>*</sup><sup>*</sup>*We have checked that the zero-momentum correlator $`\mathrm{\Pi }^{\left(reg\right)}\left(0\right)`$, obtained from the UV-regularized, unsubtracted dispersion relation with the spectral density (12) (where $`m_G`$, $`f_G`$, and $`s_0`$ have the predicted values) satisfies the low-energy theorem (17) for all four sum rules in the range of uncertainty introduced by the inaccurately known value of the gluon condensate ., however, is reassuring and of considerable conceptual importance.
For a quantitative consistency check between the predictions of different sum rules, we have evaluated the $`\tau `$-dependent mass function
$$m_G^{(1,2)}\left(\tau \right)\sqrt{\frac{_2(\tau ,s_0)}{_1(\tau ,s_0)}}.$$
(18)
Figure 3 shows that it deviates less than 2% from the constant $`m_G`$ over the whole fiducial region, indicating a high degree of compatibility between the sum rules.
We have also found that the instanton contributions to the $`_k`$ by themselves can generate stable sum rules. Their approximately exponential $`\tau `$ behavior matches very well to the pole term, although with about 20% smaller glueball masses than those obtained from the full sum rules. This indicates that instantons alone can (over-) bind the scalar glueball, in agreement with the findings of instanton vacuum models (which also show a tendency towards smaller glueball masses).
## IV Discussion and conclusions
We evaluated and analyzed the instanton contributions both to the OPE of the scalar glueball correlator (or, more precisely, to the Wilson coefficient of the unit operator) and to the continuum part of its phenomenological spectral-function model, and we solved the corresponding QCD sum rules. The previously neglected instanton contributions turn out to be dominant and render the IOPE sum rules the first overall consistent set in the scalar glueball channel.
In particular, the IOPE resolves two long-standing flaws of the earlier sum rules: the mutual inconsistency between different Borel moments and the inconsistency with the low-energy theorem for the zero-momentum correlator. Even the previously deficient and usually discarded lowest-moment ($`_1`$) sum rule becomes consistent both with those from the higher moments and with the low-energy theorem. Any evidence for a low-lying ($`m1`$ GeV) gluonium state (or a state strongly coupled to gluonic interpolators), sometimes argued for on the basis of this sum rule , is thereby rendered obsolete.
The most dramatic phenomenological impact of the direct instanton contributions is associated with $`f_G^2`$, the residuum of the glueball pole. Due to the exceptional size of the instanton contributions, its value increases by about a factor of 5. Taking the quantitative predictions of the $`_2`$ sum rule to be the most reliable ones, we obtain $`m_G=1.53\pm 0.2`$ GeV (in accord with recent lattice results ) and $`f_G=1.01\pm 0.25`$ GeV, where the errors are estimated from the uncertainties of the input parameters and the spread between the individual sum rules. Potential ramifications for experimental glueball searches will be considered in a forthcoming publication.
All four IOPE sum rules show an unprecedented degree of stability and allow for a simultaneous 3-parameter fit to the glueball mass, its coupling, and the continuum threshold. The stability region extends far beyond the fiducial $`\tau `$ interval and renders, as a side effect, the IOPE sum rule results almost insensitive to the precise boundaries of the fiducial domain. Most importantly, however, the high stability indicates that the IOPE provides a rather complete description of the short-distance glueball correlator.
A crucial contribution to the IOPE sum rules arises from the discontinuity of the instanton terms in the extended continuum part of the spectral functions. (The rough estimate of instanton contributions to the $`_0`$ sum rule in Ref. missed this contribution.) In addition to substantially improving the overall consistency and stability of the sum rules, the richer structure on the phenomenological sides also sheds new light on the spectral content of the scalar glueball correlator. For once, the quantitative sum-rule analysis reveals that the instanton contributions, together with the weaker perturbative terms, counterbalance the pole contribution. This leads to an improved description of the correlator towards low momenta and thereby reconciles the sum rules with the low-energy theorem (17), a stringent consistency checkOn the lattice, the feasibility of this check is compromised by finite-size effects. in the $`Q0`$ limit.
A remarkable interplay between perturbative and nonperturbative physics can also be seen in the opposite limit, i.e. at short distances. As a consequence of the improved continuum description, the instanton contributions to the correlator remain effective at small $`\tau `$ and stay finite even for $`\tau 0`$ (in contrast to their contributions to the IOPE in the same limit, which suffer the expected suppression associated with funnelling a sizeable momentum through the coherent instanton field). This indicates that small-instanton physics accounts not only for much of the ground-state contribution but also for part of the higher-lying spectral strength (in the sense of a generalized quark-hadron duality) in the scalar glueball correlator. Thus, the spectral distribution favored by the IOPE sum rules seems to imply a rather prominent role for instanton-induced effects in excited glueball states (or in multiparticle states with strong coupling to the energy-momentum tensor). The phenomenological impact of these results, which do not depend on details of the IOPE and should therefore be rather robust, deserves further study.
In contrast to previously studied IOPE sum rules for quark-based correlators , those for the scalar glueball are the first where (i) the instanton contributions do not enter via topological quark zero-modes and (ii) the sum rules reach a satisfactory (though not excellent) level of consistency even without any perturbative and soft contributions, i.e., with the instanton terms alone. The latter result explains why the instanton liquid model yields scalar glueball properties similar to those obtained above , and can be traced to both the exceptional strength and the particular shape (mainly the curvature) of the instanton contributions. In combination, those properties produce an approximately exponential $`\tau `$ dependence which extends far beyond the fiducial region and fairly well matches the ground-state signal without any perturbative and condensate contributions.
The predominance and approximate self-sufficiency of the instanton contributions has a suggestive physical interpretation: it indicates that instantons may generate the bulk of the attractive forces which bind the scalar glueball. Moreover, and in contrast to the instanton liquid model of Ref. , the IOPE allows to consistently compare the instanton contributions with those of the remaining soft and perturbative fields, and to thereby judge their relative importance. At present, the IOPE seems to be the only controlled and analytical framework in which this can be achieved. The combined effect of the soft and perturbative contributions turns out to be repulsive and increases, consistently in all sum rules, the glueball mass by about 20%.
Finally, closer inspection of the IOPE sum rules reveals another, quite striking instanton effect: the scales of the predicted $`0^{++}`$ glueball properties turn out to be approximately set by the bulk features of the instanton size distribution. Indeed, neglecting the standard OPE contributions one finds that the glueball parameters scale asIt is worth noting that the arguments leading to these scaling relations take advantage of the analytical character of the IOPE sum rules. The scaling would be more difficult to uncover in numerical approaches (as in lattice simulations, where in addition the instanton distribution parameters cannot easily be varied).
$`m_G`$ $``$ $`\overline{\rho }^1,`$ (19)
$`f_G^2`$ $``$ $`\overline{n}\overline{\rho }^2.`$ (20)
In the case of proportionality between the glueball size $`r_G`$ and its Compton wavelength, one would obtain another scaling relation
$$r_G\overline{\rho },$$
(21)
which could explain the small values $`r_G`$ 0.2 fm found on the lattice .
Conceptually, the main virtue of the above scaling relations lies in establishing an explicit link between fundamental vacuum and hadron properties. Although strong interdependences between QCD vacuum and hadron structure are expected on general grounds, such scaling relations seem to have not been encountered previously. They could be of practical use e.g. for the test of instanton vacuum models, to provide constraints for glueball model building, or generalized to finite temperature and baryon density, where the scales of the instanton distribution change.
This work was supported by Deutsche Forschungsgemeinschaft under Habilitation Grant Fo 156/2-1.
## V Figure captions
1. Fig. 1: The right-hand side $`f_G^2m_G^2\mathrm{exp}\left(m_G^2\tau \right)`$ of the $`_0`$ sum rule (dotted), compared with the optimized left-hand side $`_0(\tau ,s_0)/m_G^2`$ (solid line) and its three components (all in units of GeV<sup>4</sup>): the conventional OPE $`_0^{\left(OPE\right)}(\tau ,s_0)/m_G^2`$ (dash-double-dotted), the instanton contribution $`_0^{\left(I+\overline{I}\right)}\left(\tau \right)/m_G^2`$ (dashed), and the instanton continuum part $`_0^{\left(Icont\right)}(\tau ,s_0)/m_G^2`$ (dash-dotted).
2. Fig. 2: Same as Fig. 1 for the $`_2`$ sum rule.
3. Fig. 3: The square root of the ratio $`_2\left(\tau \right)/_1\left(\tau \right)`$. The weak $`\tau `$ dependence confirms the high consistency between the $`_1\left(\tau \right)`$ and $`_2\left(\tau \right)`$ sum rules.
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# Finite conformal modules over 𝑁={2,3,4} superconformal algebras
## 1. Introduction
Superconformal algebras have been playing an important role in the study of string theory and conformal field theory, which have been the subject of intensive study since the seminal paper . Superconformal algebras may be viewed as natural super-extensions of the Virasoro algebra and their roots in physics literature can be traced at least back to as early as the 70’s . A mathematically rigorous definition of a superconformal algebra is as follows. It is a simple Lie superalgebra $`𝔤`$ over the complex numbers $``$ spanned by the modes of a finite family $`𝔉`$ of mutually local fields satisfying the following two axioms :
* $`𝔉`$ contains the Virasoro field,
* the coefficients of the operator product expansions of members from $`𝔉`$ are linear combinations of members from $`𝔉`$ and their derivatives.
A Lie superalgebra $`𝔤`$ satisfying the second axiom only is referred to as a *formal distribution Lie superalgebra* in .
In order to facilitate the study of formal distribution Lie superalgebras the notion of a *conformal superalgebra* was introduced in (see Section 2). It proves to be an effective tool for this purpose.
A natural class of representations of formal distribution Lie superalgebras to study is the class of *conformal modules* . A conformal module is a pair consisting of a $`𝔤`$-module $`V`$ and a family $``$ of fields whose modes span $`V`$ such that members from $`𝔉`$ and $``$ are mutually local. Just as the study of formal distribution Lie superalgebras reduces to the study of conformal superalgebras, the study of conformal modules is essentially reduced to the study of modules over the corresponding conformal superalgebras.
The study of modules over the conformal superalgebra can further be reduced to the study of modules over the *extended annihilation subalgebra*, which is a semidirect sum of the subalgebra of positive modes of the corresponding formal distribution Lie superalgebra and a one-dimensional derivation. It is in this language that the problem of classifying finite irreducible conformal modules over the Virasoro, $`N=1`$ (Neveu-Schwarz) and the current superalgebra was solved in .
The problem of classifying conformal modules over other superconformal algebras, which is the main theme of the present paper, turns out to be more subtle. The main purpose here is to give a classification of finite irreducible conformal modules over the $`N=2`$, $`N=3`$ and the two $`N=4`$ superconformal algebras.
We first construct finite Verma-type conformal modules for a general superconformal algebra and prove that every finite irreducible conformal module is a homomorphic image of such a module. As a consequence we obtain a bijection between finite irreducible conformal modules of a superconformal algebra and finite-dimensional irreducible modules of a certain finite-dimensional reductive Lie (super)algebra (Corollary 3.1).
We then study these Verma-type modules in detail for the four members of the family of superconformal algebras mentioned above. It turns out that, unlike for the Virasoro and the $`N=1`$ (Neveu-Schwarz) superconformal algebras, the Verma-type modules for these superconformal algebras are in general reducible, and thus we need to analyze their submodules. This is accomplished by finding explicit formulas for all singular vectors inside such a module and then show that the submodule generated by these singular vectors is maximal (in all but two cases). We also find an explicit basis for this maximal submodule, which then enables us to give a quite explicit description of all finite irreducible conformal modules over these superconformal algebras.
This paper is organized as follows. In Section 2 basic facts of formal distribution Lie superalgebras, conformal superalgebras and extended annihilation subalgebras are recalled. Section 3 is devoted to the study of a class of modules over a certain class of Lie superalgebras that include the annihilation subalgebra of every superconformal algebra. This class of modules gives rise to finite Verma-type conformal modules of superconformal algebras. The results of Section 3 are then used in Section 4, Section 5, Section 6 and Section 7, where finite irreducible conformal modules over the $`N=2`$, $`N=3`$, the “small” $`N=4`$ and the “big” $`N=4`$ superconformal algebra, respectively, are classified.
In this paper all vector spaces, (super)algebras and tensor products are over taken over the complex numbers $``$.
## 2. Preliminaries
In this section we review some of the basic facts on formal distribution Lie (super)algebras and conformal modules that will be used later on. The material here is taken from , and , and the reader is referred to these articles for more details.
### 2.1. Formal Distribution Lie Superalgebras
Recall that a formal distribution or a *field* with coefficients in a Lie superalgebra $`𝔤=𝔤_{\overline{0}}+𝔤_{\overline{1}}`$ is a formal series of the form:
$$a(z)=\underset{n}{}a_{[n]}z^{n1},$$
where $`a_{[n]}𝔤`$ and $`z`$ is an indeterminate.
Two formal distributions $`a(z)`$ and $`b(z)`$ with coefficients in $`𝔤`$ are said to be mutually *local* if there exists $`N_+`$ such that
(2.1)
$$(zw)^N[a(z),b(w)]=0.$$
Let $`\delta (zw)=z^1_n(\frac{z}{w})^n`$ be the formal delta function. Then (2.1) may be written as
(2.2)
$$[a(z),b(w)]=\underset{j=0}{\overset{N1}{}}(a_{(j)}b)(w)_w^{(j)}\delta (zw),$$
(here $`_w^{(j)}`$ stands for $`\frac{1}{j!}\frac{^j}{w^j}`$) for some uniquely determined formal distributions $`(a_{(j)}b)(w)`$, and thus defines a $``$-bilinear product $`_{(j)}`$ for each $`j_+`$ on the space of all formal distributions with coefficients in $`𝔤`$. Also $`_za(z)=_n(a)_{[n]}z^{n1}`$, where $`(a)_{[n]}=na_{[n1]}`$, and hence the space of all formal distributions is also a (left) $`[_z]`$-module.
A Lie superalgebra $`𝔤`$ is called a formal distribution Lie superalgebra, if there exists a family $`𝔉`$ of mutually local formal distributions whose coefficients span $`𝔤`$. We will write $`(𝔤,𝔉)`$ for such a Lie superalgebra.
Given a formal distribution Lie superalgebra $`(𝔤,𝔉)`$, we may include $`𝔉`$ in the minimal family $`\overline{𝔉}`$ of mutally local distributions which is closed under $`_z`$ and all products $`_{(j)}`$. Then $`\overline{𝔉}`$ is a conformal superalgebra, i.e. it is a left $`_2`$-graded $`[]`$-module $`R`$ with a $``$-bilinear product $`a_{(n)}b`$ for each $`n_+`$ such that the following axioms hold ($`a,b,cR;m,n_+`$ and $`^{(j)}=\frac{1}{j!}^j`$) (cf. , ):
* $`a_{(n)}b=0`$, for $`n>>0`$,
* $`(a)_{(n)}b=na_{(n1)}b`$,
* $`a_{(n)}b=(1)^{p(a)p(b)}_{j=0}^{\mathrm{}}(1)^{j+n+1}^{(j)}(b_{(n+j)}a)`$,
* $`a_{(m)}(b_{(n)}c)=_{j=0}^{\mathrm{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(a_{(j)}b)_{(m+nj)}c+(1)^{p(a)p(b)}b_{(n)}(a_{(m)}c)`$.
It is convenient to write the products of $`a,bR`$ in the generating series form
$$a_\lambda b=\underset{n=0}{\overset{\mathrm{}}{}}a_{(n)}b\frac{\lambda ^n}{n!},$$
where $`\lambda `$ is a formal indeterminate. Such an expression lies in $`R[\lambda ]`$.
Conversely, if a conformal superalgebra $`R=_{iI}[]a^i`$ is free $`[]`$-module, we may associate to $`R`$ a formal distribution Lie superalgebra $`(𝔤(R),𝔉(R))`$ with Lie superalgebra $`𝔤(R)`$ spanned by $``$-basis $`a_{[m]}^i`$ ($`iI`$, $`m`$) and fields $`𝔉(R)=\{a^i(z)=_na_{[n]}^iz^{n1}\}_{iI}`$ with bracket (cf. (2.2)):
$$[a^i(z),a^j(w)]=\underset{k_+}{}(a_{(k)}^ia^j)(w)_w^{(k)}\delta (zw),$$
so that $`\overline{𝔉(R)}=R`$, giving rise to commutation relations ($`m,n`$; $`i,jI`$)
(2.3)
$$[a_{[m]}^i,a_{[n]}^j]=\underset{k_+}{}\left(\genfrac{}{}{0pt}{}{m}{k}\right)(a_{(k)}^ia^j)_{[m+nk]}.$$
It follows that the Lie superalgebra $`𝔤`$ of a formal distribution Lie superalgebra $`(𝔤,𝔉)`$ is isomorphic to $`𝔤(\overline{𝔉})`$ divided by an *irregular* ideal, that is an ideal which does not contain every $`a_{[n]}`$ for some non-zero element $`a\overline{𝔉}`$.
###### Example 2.1.
The (centerless) Virasoro algebra $`𝔙`$ has a basis $`L_n`$ ($`n`$) with commutation relations
$$[L_m,L_n]=(mn)L_{m+n}.$$
It is spanned by the coefficients of the field $`L(z)=_nL_nz^{n2}`$ satisfying
(2.4)
$$[L(z),L(w)]=_wL(w)\delta (zw)+2L(w)_w\delta (zw).$$
The conformal algebra associated to the Virasoro algebra, is the Virasoro conformal algebra $`R(𝔙)=[]L`$ with products $`L_\lambda L=(+2\lambda )L`$.
###### Example 2.2.
Let $`𝔤`$ be a finite-dimensional Lie (super)algebra. Let $`\stackrel{~}{𝔤}=𝔤[t,t^1]`$ denote the corresponding *current algebra* with bracket
$$[af(t),bg(t)]=[a,b]f(t)g(t),a,b𝔤;f(t),g(t)[t,t^1].$$
For each $`a𝔤`$ define a field $`a(z)=_n(at^n)z^{n1}`$. Then $`\stackrel{~}{𝔤}`$ is spanned by the coefficients of $`a(z)`$ satisfying
(2.5)
$$[a(z),b(w)]=[a,b](w)\delta (zw).$$
The conformal (super)algebra associated to the current algebra is the current conformal algebra $`R(\stackrel{~}{𝔤})=[]𝔤`$ with products $`a_\lambda b=[a,b]`$, $`a,b𝔤`$.
###### Example 2.3.
The semidirect sum $`𝔙\stackrel{~}{𝔤}`$ is another example of a formal distribution Lie (super)algebra. The collection of fields is $`\{L(z),a(z)|a𝔤\}`$ and we have in addition to (2.4) and (2.5)
(2.6)
$$[L(z),a(w)]=_wa(w)\delta (zw)+a(w)_w\delta (zw).$$
The conformal algebra associated to the semidirect sum of the Virasoro algebra and the current algebra is $`R(𝔙\stackrel{~}{𝔤})=R(𝔙)R(\stackrel{~}{𝔤})`$. For $`a𝔤`$ we have $`L_\lambda a=(+\lambda )a`$.
### 2.2. Conformal Modules
Let $`(𝔤,𝔉)`$ be a formal distribution Lie superalgebra. Let $`V`$ be a $`𝔤`$-module such that $`V`$ is spanned over $``$ by the coefficients of a family $``$ of fields. If all $`a(z)𝔉`$ are local with respect to all $`v(z)`$, then the pair $`(V,)`$ is called a conformal module over $`(𝔤,𝔉)`$.
Now the family $``$ of a conformal module $`(V,)`$ over $`(𝔤,𝔉)`$ similarly can be included in a larger family $`\overline{}`$, which is still local with respect to the fields from $`\overline{𝔉}`$, and invariant under $``$ and $`a_{(j)}`$, for all $`a\overline{𝔉}`$ and $`j_+`$. It can be shown that for $`a,b\overline{𝔉}`$ and $`v\overline{}`$ ($`m,n_+`$) one has
$$[a_{(m)},b_{(n)}]v=\underset{j=0}{\overset{m}{}}\left(\genfrac{}{}{0pt}{}{m}{j}\right)(a_{(j)}b)_{(m+nj)}v,(a)_{(n)}v=[,a_{(n)}]v=na_{(n1)}v.$$
Thus it follows that any conformal module $`(V,)`$ over a formal distribution Lie superalgebra $`(𝔤,𝔉)`$ gives rise to a module $`M=\overline{}`$ over the *conformal superalgebra* $`R=\overline{𝔉}`$, defined as follows. It is a (left) $`_2`$-graded $`[]`$-module equipped with a family of $``$-linear maps $`aa_{(n)}^M`$ of $`R`$ to $`\mathrm{End}_{}M`$, for each $`n_+`$, such that the following properties hold for $`a,bR`$ and $`m,n_+`$:
* $`a_{(n)}^Mv=0`$, for $`vM`$ and $`n>>0`$,
* $`[a_{(m)}^M,b_{(n)}^M]=_{j=0}^m\left(\genfrac{}{}{0pt}{}{m}{j}\right)(a_{(j)}b)_{(m+nj)}^M`$,
* $`(a)_{(n)}^M=[,a_{(n)}^M]=na_{(n1)}^M`$.
Again it is convenient to write the action of an element $`aR`$ on an element $`vM`$ in the form of a generating series in $`V[\lambda ]`$
$$a_\lambda v:=\underset{n=0}{\overset{\mathrm{}}{}}a_{(n)}v\frac{\lambda ^n}{n!}.$$
Conversely, suppose that a conformal superalgebra $`R=_{iI}[]a^i`$ is a free $`[]`$-module and consider the associated formal distribution Lie superalgebra $`(𝔤(R),𝔉(R))`$. Let $`M`$ be a module over the conformal superalgebra $`R`$ and suppose that $`M`$ is a free $`[]`$-module with $`[]`$-basis $`\{v^\alpha \}_{\alpha J}`$. This gives rise to a conformal module $`V(M)`$ over $`𝔤(R)`$ with fields $`=\{v^\alpha (z)=_nv_{[n]}^\alpha z^{n1}|\alpha J\}`$ and $``$-basis $`v_{[n]}^\alpha `$, defined by:
$$a^i(z)v^\alpha (w)=\underset{j_+}{}(a_{(j)}^iv^\alpha )(w)_w^{(j)}\delta (zw).$$
A conformal module $`(V,)`$ (respectively module $`M`$) over a formal distribution Lie superalgebra $`(𝔤,𝔉)`$ (respectively over a conformal superalgebra $`R`$) is called finite, if $`\overline{}`$ (respectively $`M`$) is a finitely generated $`[]`$-module. A conformal module $`(V,)`$ over $`(𝔤,𝔉)`$ is called irreducible, if there is no non-trivial invariant subspace which contains all $`v_{[n]}`$, $`n`$, for some non-zero $`v\overline{}`$. An invariant subspace that does not contain all $`v_{[n]}`$, for some non-zero $`v`$, is called an *irregular submodule* and conformal modules that differ by an irregular submodule are called referred to as *equivalent* in . Clearly a conformal module is irreducible if and only if the associated module $`\overline{}`$ over the conformal superalgebra $`\overline{𝔉}`$ is irreducible.
###### Remark 2.1.
It follows from (M2) that an eigenvector $`vM`$ of the linear operator $``$ is an $`R`$-invariant, i.e. $`a_{(n)}v=0`$, for all $`n0`$. Thus a finite irreducible module over a conformal superalgebra $`R`$ is either free over $`[]`$ or else it is one-dimensional over $``$.
Suppose that $`(𝔤,𝔉)`$ is a formal distribution Lie superalgebra such that $`𝔤(\overline{𝔉})𝔤`$. Our discussion implies that any irreducible conformal module $`(V,)`$ over $`(𝔤,𝔉)`$ is a quotient of an irreducible conformal module of the form $`V(M)`$ divided by an irregular submodule, where $`M`$ is an irreducible module over the conformal superalgebra $`\overline{𝔉}`$. Hence in particular if $`V(M)`$ is irreducible as a $`𝔤`$-module for every irreducible $`M`$, then every finite irreducible conformal modules over $`(𝔤,𝔉)`$ isomorphic to $`V(M)`$, for some finite irreducible $`\overline{𝔉}`$-module $`M`$.
###### Example 2.4.
The Virasoro algebra $`𝔙`$ may be identified with the Lie algebra of regular vector fields on $`^\times `$, where $`L_n=t^{n+1}\frac{d}{dt}`$, $`n`$. For $`\alpha ,\mathrm{\Delta }`$ let
$$F_𝔙(\alpha ,\mathrm{\Delta })=[t,t^1]e^{\alpha t}dt^{1\mathrm{\Delta }}.$$
The Lie algebra $`𝔙`$ acts on the space $`F_𝔙(\alpha ,\mathrm{\Delta })`$ in a natural way:
$$(f(t)\frac{}{t})g(t)dt^{1\mathrm{\Delta }}=(f(t)g^{}(t)+(1\mathrm{\Delta })g(t)f^{}(t))dt^{1\mathrm{\Delta }},$$
where $`f(t)[t,t^1]`$ and $`g(t)[t,t^1]e^{\alpha t}`$. Letting $`v_{[n]}=t^ne^{\alpha t}dt^{1\mathrm{\Delta }}`$ and $`v(z)=_nv_{[n]}z^{n1}`$ this action is equivalent to
$$L(z)v(w)=(_w+\alpha )v(w)\delta (zw)+\mathrm{\Delta }v(w)_w\delta (zw).$$
Hence we have constructed a two-parameter family of conformal modules over $`𝔙`$. This gives a family of $`R(𝔙)`$-modules $`[]v_\mathrm{\Delta }`$ with products $`L_\lambda v_\mathrm{\Delta }=(\alpha ++\mathrm{\Delta }\lambda )v_\mathrm{\Delta }`$. This module is irreducible if and only if $`\mathrm{\Delta }0`$, in which case it will be denoted by $`L_𝔙(\alpha ,\mathrm{\Delta })`$. We set $`L_𝔙(\alpha ,0)`$ to be the one-dimensional (over $``$) $`R(𝔙)`$-module on which $``$ acts as the scalar $`\alpha `$.
###### Example 2.5.
Let $`𝔤`$ be a finite-dimensional simple Lie algebra and $`U^\mathrm{\Lambda }`$ the finite-dimensional irreducible module of highest weight $`\mathrm{\Lambda }`$. Then $`F_{\stackrel{~}{𝔤}}(\mathrm{\Lambda })=U^\mathrm{\Lambda }[t,t^1]`$ is naturally a module over $`\stackrel{~}{𝔤}`$ with action given by
(2.7)
$$(af(t))(ug(t))=auf(t)g(t),a𝔤,uU^\mathrm{\Lambda };f(t),g(t)[t,t^1].$$
For each vector $`uU^\mathrm{\Lambda }`$ define $`u(z)=_n(ut^n)z^{n1}`$ so that (2.7) is equivalent to
$$a(z)u(w)=au(w)\delta (zw),$$
and hence $`F_{\stackrel{~}{𝔤}}(\mathrm{\Lambda })`$ is conformal. This gives a family of $`R(\stackrel{~}{𝔤})`$-modules, which is irreducible if and only if $`\mathrm{\Lambda }0`$, in which case it will be denoted by $`L_{\stackrel{~}{𝔤}}(\mathrm{\Lambda })`$. By $`L_{\stackrel{~}{𝔤}}(0)`$ we will mean the trivial $`R(\stackrel{~}{𝔤})`$-module. Similarly one defines the one-dimensional module $`L_{\stackrel{~}{𝔤}}(\alpha ,0)`$.
###### Example 2.6.
$`\stackrel{~}{𝔤}`$ acts on $`F_{𝔙\stackrel{~}{𝔤}}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })=U^\mathrm{\Lambda }F_𝔙(\alpha ,\mathrm{\Delta })`$ similarly as in Example 2.5. However, on $`F_{𝔙\stackrel{~}{𝔤}}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })`$ we have also an action of $`𝔙`$, thus making it into a module over $`𝔙\stackrel{~}{𝔤}`$. This module defines an $`R(𝔙\stackrel{~}{𝔤})`$-module which is irreducible if and only if $`(\mathrm{\Delta },\mathrm{\Lambda })(0,0)`$, and in which case it will be denoted by $`L_{𝔙\stackrel{~}{𝔤}}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })`$. By $`L_{𝔙\stackrel{~}{𝔤}}(\alpha ,0,0)`$ we will mean the one-dimensional module on which $``$ acts a the scalar $`\alpha `$.
The following theorem was proved in .
###### Theorem 2.1.
Let $`𝔤`$ stand for a finite-dimensional simple Lie algebra. Any finite irreducible module over the conformal algebras $`R(𝔙)`$, $`R(\stackrel{~}{𝔤})`$ and $`R(𝔙\stackrel{~}{𝔤})`$ are as follows:
* $`L_𝔙(\alpha ,\mathrm{\Delta })`$,
* $`L_{\stackrel{~}{𝔤}}(\mathrm{\Lambda })`$ and $`L_{\stackrel{~}{𝔤}}(\alpha ,0)`$,
* $`L_{𝔙\stackrel{~}{𝔤}}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })`$.
###### Remark 2.2.
We note that a similar statement as Theorem 2.1 part (iii) holds even if $`𝔤`$ is replaced by the $`1`$-dimensional Lie algebra $`a`$. In this case $`U^\mathrm{\Lambda }=u`$ with $`au=\mathrm{\Lambda }u`$, $`\mathrm{\Lambda }`$. Also part (ii) remains true for all but three series of finite-dimensional simple Lie superalgebras.
### 2.3. Extended Annihilation Subalgebras
Given a formal distribution Lie superalgebra $`(𝔤,𝔉)`$ we let $`𝔤_+`$ denote the $``$-span of all $`a_{[n]}`$, where $`n0`$ and $`a𝔉`$. Due to (2.3) $`𝔤_+`$ is closed under the bracket and hence form a subalgebra of $`𝔤`$, which we will call the annihilation algebra of $`(𝔤,𝔉)`$. Let $``$ be the derivation of $`𝔤_+`$ defined by $`[,a_{[n]}]=na_{[n1]}`$, and consider the semi-direct sum of $`𝔤^+=𝔤_+`$. Then $`𝔤^+`$ is called the *extended annihilated algebra* of $`(𝔤,𝔉)`$. The following proposition, which follows by comparing (M1) with (2.3), is important for the theory of conformal modules.
###### Proposition 2.1.
Let $`R`$ be a conformal superalgebra and $`(𝔤(R),R(𝔉))`$ be its associated formal distribution Lie superalgebra with extended annihilation algebra $`𝔤(R)^+`$. Then a module over the conformal superalgebra $`R`$ is precisely a $`𝔤(R)^+`$-module $`M`$ satisfying $`a_{[n]}v=0`$, for each $`vM`$, $`aR`$ and $`n>>0`$.
###### Remark 2.3.
Let $`R`$ be a conformal superalgebra with $`[]`$-basis $`\{a^i|iI\}`$ and $`M`$ a free $`[]`$-module with basis $`\{v^j|jJ\}`$. Given $`a_{(n)}^iv^jM`$ for all $`iI`$, $`jJ`$, $`n_+`$, which is $`0`$ for $`n>>0`$, condition (M2) uniquely extends the action of $`a_{(n)}^i`$ to all of $`M`$. If in addition (M1) holds, then $`M`$ is an $`R`$-module. Hence the action of an $`R`$-module $`M`$ is completely determined by the action of a $`[]`$-basis of $`R`$ on a $`[]`$-basis of $`M`$.
###### Example 2.7.
In the case of the Virasoro algebra $`𝔙`$ the annihilation algebra $`𝔙_+`$ is spanned by elements $`L_n`$, $`n1`$. In the case of the current algebra $`\stackrel{~}{𝔤}_+`$ is spanned by $`at^n`$, where $`a𝔤`$ and $`n0`$, while in the case of $`𝔙\stackrel{~}{𝔤}`$ it is $`𝔙_+\stackrel{~}{𝔤}_+`$.
The problem of classifying conformal modules over $`(𝔤,𝔉)`$ is thus reduced to the problem of classifying a class of modules over $`𝔤(\overline{𝔉})^+`$. It is clear that in all our examples one has $`𝔤(\overline{𝔉})=𝔤`$, and thus we are to study modules over $`𝔤^+`$. Now if in addition there exists an element $`L_1`$ in $`𝔤_+`$ such that $`L_1`$ is central in $`𝔤^+`$, then every irreducible representation of $`𝔤^+`$ is an irreducible representation of $`𝔤_+`$, on which $`(L_1)`$ acts as a scalar $`\alpha `$. In the case of the $`𝔙`$ and $`𝔙\stackrel{~}{𝔤}`$ and the $`N=2,3,4`$ superconformal superalgebras, which we will define later, such an $`L_1`$ always exists so that we only need to consider representations of $`𝔤_+`$. The irreducible representations of $`𝔙_+`$, and $`𝔙_+\stackrel{~}{𝔤}_+`$ that give rise to those in Theorem 2.1 are denoted by $`L_{𝔙_+}(\mathrm{\Delta })`$ and $`L_{𝔙_+\stackrel{~}{𝔤}_+}(\mathrm{\Delta },\mathrm{\Lambda })`$, respectively. The corresponding actions are clear and can be found in .
## 3. Finite Verma-type Conformal Modules
Let $``$ be a Lie superalgebra over $``$ with a distinguished element $``$ and a descending sequence of subspaces $`=_1_0_1_2\mathrm{}_n\mathrm{}`$, such that $`[,_k]=_{k1}`$, for all $`k>0`$. Let $`W`$ be an $``$-module, which is finitely generated over $`[]`$, such that for all $`wW`$ there exists a non-negative integer $`k`$ (depending on $`w`$) with $`_kw=0`$. For $`m2`$ set $`W_m=\{wW|_{m+1}w=0\}`$ and let $`M`$ be the minimal non-negative integer such that $`W_M0`$.
###### Lemma 3.1.
Suppose that $`M0`$. Then $`[]W_M=[]W_M`$ and hence $`[]W_MW_M=W_M`$. In particular $`W_M`$ is a finite-dimensional vector space.
Let $`𝔤`$ be a Lie superalgebra satisfying the following three conditions.
* $`𝔤`$ is $``$-graded of finite depth $`d`$, i.e. $`𝔤=_{jd}𝔤_j`$ with $`[𝔤_i,𝔤_j]𝔤_{i+j}`$.
* There exists a semisimple element $`z𝔤_0`$ such that it centralizer in $`𝔤`$ is contained in $`𝔤_0`$.
* There exits an element $`𝔤_d`$ such that $`[,𝔤_i]=𝔤_{id}`$, for $`i0`$.
###### Remark 3.1.
If $`𝔤`$ contains the grading operator with respect to its gradation, then condition (L2) is automatic.
Examples of Lie superalgebras satisfying (L1)–(L3) are provided by annihilation subalgebras of superconformal algebras, which we will describe in more detail.
Let $`t`$ be an even indeterminate and $`\xi _1,\mathrm{},\xi _N`$ be $`N`$ odd indeterminate. Denote by $`\mathrm{\Lambda }(N)`$ the Grassmann superalgebra in the indeterminates $`\xi _1,\mathrm{},\xi _N`$ and set $`\mathrm{\Lambda }(1,N):=[t,t^1]\mathrm{\Lambda }(N)`$. Let $`W(1,N)`$ be the derivation superalgebra of $`\mathrm{\Lambda }(1,N)`$, then $`W(1,N)`$ is a formal distribution Lie superalgebra . Letting $`\frac{}{t}`$ and $`\frac{}{\xi _i}`$, for $`i=1,\mathrm{},N`$, be the usual differential operators, every element in $`DW(1,N)`$ can be written as
$$D=a_0\frac{}{t}+\underset{i=1}{\overset{N}{}}a_i\frac{}{\xi _i},a_0,a_i,\mathrm{},a_N\mathrm{\Lambda }(1,N).$$
The *standard gradation* of $`W(1,N)`$ is obtained by setting the degree of $`t`$ and $`\xi _i`$ to be $`1`$. Its annihilation subalgebra is $`W(1,N)_+=_{j1}(W(1,N))_j`$. $`W(1,N)_+`$ in this gradation contains its grading operator given by $`z=t\frac{}{t}+_{i=1}^N\xi _i\frac{}{\xi _i}`$ so that (L2) is satisfied. Also choosing $``$ to be $`\frac{}{t}`$ it follows that (L3) is also satisfied so that $`W(1,N)`$ is a Lie superalgebra of the type above. Note that $`W(1,N)_0gl(1,N)`$.
The subalgebra of divergence zero vector fields in $`W(1,N)`$ contains an ideal of codimension $`1`$. This ideal is its derived algebra and is the superconformal algebra $`S(1,N)`$ . The standard gradation of $`W(1,N)_+`$ induces a gradation on the annihilation subalgebra $`S(1,N)_+`$ of $`S(1,N)`$. Choosing $`z=t\frac{}{t}+\frac{1}{N}_{i=1}^N\xi _i\frac{}{\xi _i}`$ along with $`=\frac{}{t}`$ it follows that $`S(1,N)_+`$ in this gradation also satisfies (L1)–L(3). Observe that $`S(1,N)_0sl(1,N)`$ and also that the “small” $`N=4`$ superconformal algebra (to be defined in Section 6) is isomorphic to $`S(1,2)`$ .
The contact superalgebra $`K(1,N)`$ is the subalgebra of $`W(1,N)`$ defined by
$$K(1,N):=\{DW(1,N)|D\omega =f_D\omega ,\mathrm{for}\mathrm{some}f_D\mathrm{\Lambda }(1,N)\},$$
where $`\omega :=dt_{i=1}^N\xi _id\xi _i`$ is the standard contact form. Here the action of $`D`$ on $`\omega `$ is the usual action of vector fields on differential forms.
The map from $`\mathrm{\Lambda }(1,N)`$ to $`K(1,N)`$ given by to
$$f2f\frac{}{t}+(1)^{p(f)}\underset{i=1}{\overset{N}{}}(\xi _i\frac{f}{t}+\frac{f}{\xi _i})(\xi _i\frac{}{t}+\frac{}{\xi _i})$$
is a bijection and hence it allows us to identify $`K(1,N)`$ with the polynomial superalgebra $`\mathrm{\Lambda }(1,N)`$. The Lie bracket in $`\mathrm{\Lambda }(1,N)`$, also called the contact bracket, then reads for homogeneous elements $`f,g\mathrm{\Lambda }(1,N)`$:
$$[f,g]=(2E)f\frac{g}{t}\frac{f}{t}(2E)g+(1)^{p(f)}\underset{i=1}{\overset{N}{}}\frac{f}{\xi _i}\frac{g}{\xi _i},$$
where $`E=_{i=1}^N\xi _i\frac{}{\xi _i}`$ is the Euler operator.
When $`N`$ is even it is sometimes more convenient to make the change of basis $`\xi _j^+=\frac{1}{\sqrt{2}}(\xi _j+i\xi _{j+\frac{N}{2}})`$ and $`\xi _j^{}=\frac{1}{\sqrt{2}}(\xi _ji\xi _{j+\frac{N}{2}})`$, for $`j=1,\mathrm{},\frac{N}{2}`$ and $`i=\sqrt{1},`$ so that the contact bracket takes the split form:
$$[f,g]=(2E)f\frac{g}{t}\frac{f}{t}(2E)g+(1)^{p(f)}\underset{i=1}{\overset{\frac{N}{2}}{}}(\frac{f}{\xi _i^+}\frac{g}{\xi _i^{}}+\frac{f}{\xi _i^{}}\frac{g}{\xi _i^+}),$$
where $`E`$ again is the Euler operator $`_{i=1}^{\frac{N}{2}}(\xi _i^+\frac{}{\xi _i^+}+\xi _i^{}\frac{}{\xi _i^{}})`$.
The contact superalgebra $`K(1,N)`$ is a formal distribution Lie superalgebra with fields defined as follows: Let $`I=\{i_1,\mathrm{},i_k\}`$ be an ordered subset of $`\{1,\mathrm{},N\}`$, and denote by $`\xi _I`$ the monomial $`\xi _{i_1}\mathrm{}\xi _{i_k}`$. Each such monomial gives rise to a field $`\xi _I(z)=_j\xi _It^jz^{j1}`$. Evidently the span of the coefficients of all such $`\xi _I(z)`$ is $`K(1,N)`$. Furthermore it is easy to check that these fields are mutually local and form a formal distribution Lie superalgebra. This Lie superalgebra becomes $``$-graded by putting the degree of $`\xi _It^n`$ to $`2n+k2`$. Obviously $`t`$ is the grading operator of this gradation. This gradation of $`K(1,N)`$ is usually referred to as its *standard gradation*.
The annihilation subalgebra $`K(1,N)_+`$ of $`K(1,N)`$ is spanned by the basis elements $`\xi _It^n`$, where $`n0`$ and $`I`$ runs over all subsets of $`\{i_1,\mathrm{},i_k\}`$ ordered in (strictly) increasing order. The $``$-gradation from $`K(1,N)`$ induces a gradation on $`K(1,N)_+`$ making it a $``$-graded Lie superalgebra of depth $`2`$ so that $`K(1,N)_+=_{j=2}^{\mathrm{}}(K(1,N)_+)_j`$ satisfies (L1) and (L2). In this gradation it is easy to check that $`[1,(K(1,N)_+)_j]=(K(1,N)_+)_{j2}`$ for all $`j0`$, so that $`K(1,N)_+`$ also satisfies condition (L3). It is easy to see that the annihilation subalgebra of the small $`N=4`$ superconformal algebra, which we define in Section 6, also satisfies conditions (L1)–(L3). Note that $`K(1,N)_0cso_N`$, the direct sum of the Lie algebra $`so_N`$ and the one-dimensional Lie algebra.
Finally it follows from the description of the exceptional superconformal algebra $`CK_6`$ as a subalgebra of $`K(1,6)`$ in that its annihilation subalgebra $`(CK_6)_+=_{j2}(CK_6)_j`$ is a Lie superalgebra satisfying (L1)–(L3) with $`(CK_6)_0cso_6`$.
The modules over the annihilation subalgebras that are equivalent to modules over the corresponding conformal superalgebras are then $`𝔤`$-modules $`V`$ satisfying the following conditions.
* For all $`vV`$ there exists an integer $`k_0d`$ (depending on $`v`$) such that $`𝔤_kv=0`$, for all $`kk_0`$.
* $`V`$ is finitely generated over $`[]`$.
We shall call $`𝔤`$-modules satisfying these two properties *finite*. Let $`V`$ be a finite irreducible $`𝔤`$-module. For $`nd1`$ set $`V_n=\{vV|𝔤_jv=0,j>n\}`$. Let $`N`$ be the minimal integer such that $`V_N0`$. Such an $`N`$ exists by (V1).
###### Lemma 3.2.
If $`N0`$, then $`V_N`$ is a finite-dimensional vector space over $``$.
###### Proof.
We let $`=𝔤`$ and put $`_j=_{ijd}𝔤_i`$ so that we have a filtration of subspaces
$$_0_1_2\mathrm{}_n\mathrm{},$$
with $`[,_i]=_{i1}`$, for all $`i0`$ by (L3). Let $`W_m:=\{vV|_{m+1}v=0\}`$ and let $`M`$ be the minimal integer such that $`W_M0`$. Since $`N0`$ implies that $`M0`$, this setting puts us in the situation of Lemma 3.1, from which we conclude that $`W_M`$ is a finite-dimensional vector space over $``$. Of course $`V_NW_M`$ and hence it follows that $`V_N`$ is finite-dimensional as well. ∎
We obtain the following description of finite irreducible $`𝔤`$-modules.
###### Theorem 3.1.
Let $`𝔤=_{jd}𝔤_j`$ be a Lie superalgebra satisfying conditions (L1)–(L3) and $`V`$ a finite irreducible $`𝔤`$-module. There exists a finite-dimensional irreducible $`𝔤_0`$-module $`U_0`$, extended trivially to an $`_0(=_{j0}𝔤_j)`$-module, and a $`𝔤`$-epimorphism $`\phi :\mathrm{Ind}__0^𝔤U_0V`$.
###### Proof.
We will continue to use the notation defined earlier. First we show that $`N0`$. Suppose that $`N>0`$. It is easy to see that $`V_N`$ is invariant under $`_0`$. Now there exits a basis $`\{x_1,\mathrm{},x_m\}`$ of $`𝔤_N`$ together with non-zero complex number $`\lambda _1,\mathrm{},\lambda _m`$ such that $`[z,x_i]=\lambda _ix`$, where $`z`$ is the element of (L2). Since $`V_N`$ is a finite-dimensional vector space it follows in particular that $`x_i`$ acts nilpotently on $`V_N`$ for all $`1im`$. But $`[𝔤_N,𝔤_N]_{jN+1}𝔤_j`$ and so the action of the $`x_i`$’s on $`V_N`$ commutes. Therefore there exits a non-zero $`vV_N`$ such that $`𝔤_Nv=0`$. But in this case $`V_{N1}0`$, which contradicts the minimality of $`N`$. Thus $`N0`$.
In the case when $`N=0`$, there exists an epimorphism of $`𝔤`$-modules $`\mathrm{Ind}__0^𝔤V_0V`$, with $`V_0`$ finite-dimensional due to Lemma 3.2. By irreducibility of $`V`$ it follows that $`V_0=U_0`$ is an irreducible $`𝔤_0`$-module. Now if $`N<0`$, then there exists a non-zero vector $`v`$ invariant under the action of $`𝔤_j`$, for $`j0`$. Again we have an epimorphism of $`𝔤`$-modules $`\mathrm{Ind}__0^𝔤vV`$. ∎
As a corollary of Theorem 3.1 we obtain the following.
###### Corollary 3.1.
There exists a bijection between finite irreducible conformal modules of the superconformal algebra $`𝔤`$ and finite-dimensional irreducible representations of the Lie (super)algebra $`𝔤_0`$, where
* $`𝔤=K(1,N)`$ and $`𝔤_0=cso_N`$,
* $`𝔤=W(1,N)`$ and $`𝔤_0=gl(1,N)`$,
* $`𝔤=S(1,N)`$ and $`𝔤_0=sl(1,N)`$,
* $`𝔤=CK_6`$ and $`𝔤_0=cso_6`$.
###### Proof.
By Theorem 3.1 every finite irreducible $`𝔤`$-module is a homomorphic image of $`\mathrm{Ind}__0^𝔤U_0`$. Now the usual argument for highest weight representations implies that given a finite-dimensional irreducible $`𝔤_0`$-module $`U_0`$ the $`𝔤`$-module $`\mathrm{Ind}__0^𝔤U_0`$ contains a unique maximal submodule, from which the bijection then follows. ∎
###### Remark 3.2.
It is usual to put a half-integer gradation on $`K(1,N)`$ when thinking of it as a superconformal algebra. The grading operator of $`K(1,N)`$ with respect to this gradation is then $`\frac{t}{2}`$ rather than $`t`$. In this gradation one has $`K(1,N)_+=_{j1}𝔤_j`$, where $`j\frac{1}{2}`$. Theorem 3.1 of course remains valid after making some obvious changes regarding gradation. For a Lie superalgebra $`𝔤=_{j1}𝔤_j`$ with $`j\frac{1}{2}`$, we will make it a convention to write $`𝔤_{}`$ for the subalgebra $`_{j<0}𝔤_j`$.
## 4. Finite irreducible Modules over the $`N=2`$ conformal superalgebra
The $`N=2`$ superconformal algebra is the formal distribution Lie superalgebra $`K(1,2)`$. Letting $`\xi ^+,\xi ^{}`$ denote the two odd indeterminates (so that we are using the split contact form) this algebra is generated by the following four fields: $`L(z)=_n\frac{t^{n+1}}{2}z^{n2}`$, $`G^\pm (z)=_{r\frac{1}{2}+}\xi ^\pm t^{r+\frac{1}{2}}z^{r\frac{3}{2}}`$ and $`J(z)=_n\xi ^{}\xi ^+t^nz^{n1}`$. Its corresponding conformal superalgebra is then generated freely over $`[]`$ by $`\{L,J,G^\pm \}`$ with products:
$`L_\lambda L=(+2\lambda )L,L_\lambda J=(+\lambda )J,L_\lambda G^\pm =(+{\displaystyle \frac{3}{2}}\lambda )G^\pm ,`$
$`J_\lambda G^\pm =\pm G^\pm ,G_\lambda ^+G^{}=(+2\lambda )J+2L.`$
Letting $`L_n=\frac{t^{n+1}}{2}`$, $`G_r^\pm =\xi ^\pm t^{r+\frac{1}{2}}`$ and $`J_n=\xi ^{}\xi ^+t^n`$ with $`n`$, $`r\frac{1}{2}+`$, the non-zero brackets in $`K(1,2)`$ are ($`m,n`$ and $`r,s\frac{1}{2}+`$):
$`[L_m,L_n]=(mn)L_{m+n},[L_m,G_r^\pm ]=({\displaystyle \frac{m}{2}}r)G_{m+r}^\pm ,[L_m,J_n]=nJ_{n+m},`$
$`[J_m,G_r^\pm ]=\pm G_{m+r}^\pm ,[G_r^+,G_s^{}]=2L_{r+s}+(rs)J_{r+s}.`$
The annihilation subalgebra $`𝔤=K(1,2)_+`$ is then spanned by $`L_m`$, $`J_n`$ and $`G_r^\pm `$, where $`m1`$, $`n0`$ and $`r\frac{1}{2}`$. Note that letting $`𝔤_j`$ be the span of $`X_j`$, where $`X=L,J,G^\pm `$, equips $`𝔤=_{j1}𝔤_j`$, $`j\frac{1}{2}`$, with a (consistent) $`\frac{1}{2}`$-gradation. We denote $`L_1`$ by $``$ from now on.
Let $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, $`\mathrm{\Delta },\mathrm{\Lambda }`$, be the one-dimensional module over the abelian Lie algebra $`𝔤_0=L_0+J_0`$, determined by
$$L_0v_{\mathrm{\Delta },\mathrm{\Lambda }}=\mathrm{\Delta }v_{\mathrm{\Delta },\mathrm{\Lambda }},J_0v_{\mathrm{\Delta },\mathrm{\Lambda }}=\mathrm{\Lambda }v_{\mathrm{\Delta },\mathrm{\Lambda }}.$$
We may extend $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ to a module over $`_0=_{j0}𝔤_j`$ by setting $`𝔤_jv_{\mathrm{\Delta },\mathrm{\Lambda }}=0`$, for $`j>0`$. Let $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda }):=\mathrm{Ind}__0^𝔤v_{\mathrm{\Delta },\mathrm{\Lambda }}`$. We denote by $`N`$ the unique maximal submodule of $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$. The quotient $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })/N`$ is the irreducible highest weight module $`L_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ of highest weight $`(\mathrm{\Delta },\mathrm{\Lambda })`$. By Theorem 3.1 $`L_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ for $`\mathrm{\Delta },\mathrm{\Lambda }`$ form a complete list of finite irreducible $`K(1,2)_+`$-modules. Our next objective is to give a more explicit description of $`N`$ and hence of $`L_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$.
It is clear that $`^kv_{\mathrm{\Delta },\mathrm{\Lambda }}`$, $`^kG_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, $`^kG_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ and $`^kG_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, $`k0`$, is a basis consisting of $`(L_0,J_0)`$-weight vectors for $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ of $`(L_0,J_0)`$-weights $`(\mathrm{\Delta }+k,\mathrm{\Lambda })`$, $`(\mathrm{\Delta }+k+\frac{1}{2},\mathrm{\Lambda }+1)`$, $`(\mathrm{\Delta }+k+\frac{1}{2},\mathrm{\Lambda }1)`$ and $`(\mathrm{\Delta }+k+1,\mathrm{\Lambda })`$, respectively. A non-zero $`(L_0,J_0)`$-weight vector $`vM_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ is called a *singular vector* if $`𝔤_jv=0`$, for all $`j>0`$. We call a singular vector *proper* if it is not a scalar multiple of the highest weight vector $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$. Obviously $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ is irreducible if and only if $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ contains no proper singular vector. We now analyze singular vectors inside $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$.
###### Lemma 4.1.
Let $`k1`$ and suppose that $`w=\alpha ^kv_{\mathrm{\Delta },\mathrm{\Lambda }}+\beta ^{k1}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ is a singular vector of $`(L_0,J_0)`$-weight $`(\mathrm{\Delta }+k,\mathrm{\Lambda })`$ in $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$, where $`\alpha ,\beta `$. Then $`k=1`$. Furthermore any proper singular vector of this form is a scalar multiple of either $`G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, in which case $`\mathrm{\Delta }=\frac{1}{2}`$ and $`\mathrm{\Lambda }=1`$, or $`(2+G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{})v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, in which case $`\mathrm{\Delta }=\frac{1}{2}`$ and $`\mathrm{\Lambda }=1`$.
###### Proof.
Note that $`w`$ is singular if and only if $`J_1w=G_{\frac{1}{2}}^\pm w=0`$. We compute
(4.1) $`G_{\frac{1}{2}}^+w`$ $`=(\alpha k\beta (2\mathrm{\Delta }+\mathrm{\Lambda }))^{k1}G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}=0,`$
(4.2) $`G_{\frac{1}{2}}^{}w`$ $`=(\alpha k+\beta (2\mathrm{\Delta }\mathrm{\Lambda }+2k))^{k1}G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}=0,`$
(4.3) $`J_1w`$ $`=(\alpha \mathrm{\Lambda }k+\beta (2\mathrm{\Delta }+\mathrm{\Lambda }))^{k1}v_{\mathrm{\Delta },\mathrm{\Lambda }}+\beta (k1)\mathrm{\Lambda }^{k2}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}=0.`$
But then $`\beta 0`$, since otherwise (4.1) would imply that $`k=0`$. However, $`\beta 0`$ together with (4.1) and (4.2) implies that
(4.4)
$$2\mathrm{\Delta }+k=0.$$
Now (4.3) gives
(4.5)
$$\alpha \mathrm{\Lambda }k+\beta (2\mathrm{\Delta }+\mathrm{\Lambda })=0,\beta (k1)\mathrm{\Lambda }=0.$$
Now if $`k>1`$, then (4.5) gives $`\mathrm{\Lambda }=0`$ and $`\mathrm{\Delta }=0`$. But then $`k=0`$ by (4.1). Hence $`k=1`$ so that by (4.4) we have $`\mathrm{\Delta }=\frac{1}{2}`$.
Now if $`\alpha 0`$, we have from (4.1) and (4.3) $`\alpha (1+\mathrm{\Lambda })=0`$ and hence $`\mathrm{\Lambda }=1`$. The first equation of (4.5) then implies that $`\alpha +2\beta =0`$.
On the other hand if $`\alpha =0`$, the first equation of (4.5) gives $`\mathrm{\Lambda }=1`$. ∎
###### Lemma 4.2.
Let $`k_+`$.
* If $`^kG_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ is a singular vector of $`(L_0,J_0)`$-weight $`(\mathrm{\Delta }+k+\frac{1}{2},\mathrm{\Lambda }+1)`$, then $`k=0`$ and $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$. Furthermore in this case $`G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ is a singular vector.
* If $`^kG_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ is a singular vector of $`(L_0,J_0)`$-weight $`(\mathrm{\Delta }+k+\frac{1}{2},\mathrm{\Lambda }1)`$, then $`k=0`$ and $`2\mathrm{\Delta }+\mathrm{\Lambda }=0`$. Furthermore in this case $`G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ is a singular vector.
###### Proof.
The lemma follows immediately from the following two equations:
$`G_{\frac{1}{2}}^{}^kG_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ $`=(2\mathrm{\Delta }\mathrm{\Lambda }+2k)^kv_{\mathrm{\Delta },\mathrm{\Lambda }}+k^{k1}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}=0,`$
$`G_{\frac{1}{2}}^+^kG_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ $`=(2\mathrm{\Delta }+\mathrm{\Lambda })^kv_{\mathrm{\Delta },\mathrm{\Lambda }}+k^{k1}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}=0.`$
Thus Lemma 4.1 and Lemma 4.2 prove the following.
###### Proposition 4.1.
Any proper singular vector in $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ is a scalar multiple of
* $`G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, in which case we have $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$. In the particular case of $`\mathrm{\Delta }=\frac{1}{2}`$ and $`\mathrm{\Lambda }=1`$ we have in addition $`G_{\frac{1}{2}}^{}G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}`$.
* $`G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, in which case we have $`2\mathrm{\Delta }+\mathrm{\Lambda }=0`$. In the particular case of $`\mathrm{\Delta }=\frac{1}{2}`$ and $`\mathrm{\Lambda }=1`$ we have in addition $`G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$.
Let $`N`$ be the subspace of $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ given by
$`N`$ $`=[]G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}+[]G_{\frac{1}{2}}^{}G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }},\mathrm{if}2\mathrm{\Delta }\mathrm{\Lambda }=0\mathrm{and}\mathrm{\Lambda }0,`$
$`N`$ $`=[]G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}+[]G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }},\mathrm{if}2\mathrm{\Delta }+\mathrm{\Lambda }=0\mathrm{and}\mathrm{\Lambda }0.`$
It follows from Proposition 4.1 that in either case $`N`$ is a submodule of $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$.
###### Theorem 4.1.
The modules $`L_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$, for $`\mathrm{\Delta },\mathrm{\Lambda }`$, form a complete list of non-isomorphic finite (over $`[]`$) irreducible $`K(1,2)_+`$-modules. Furthermore $`L_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ as a $`[]`$-module has rank
* $`4`$, in the case $`2\mathrm{\Delta }\pm \mathrm{\Lambda }0`$,
* $`2`$, in the case $`2\mathrm{\Delta }\pm \mathrm{\Lambda }=0`$ and $`2\mathrm{\Delta }\mathrm{\Lambda }0`$,
* $`0`$, in the case $`\mathrm{\Delta }=\mathrm{\Lambda }=0`$.
###### Proof.
If $`2\mathrm{\Delta }+\mathrm{\Lambda }0`$ and $`2\mathrm{\Delta }\mathrm{\Lambda }0`$, then by Proposition 4.1 $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ contains no proper singular vector and hence is irreducible.
Suppose that $`2\mathrm{\Delta }+\mathrm{\Lambda }=0`$ and $`2\mathrm{\Delta }\mathrm{\Lambda }0`$. In this case consider the submodule of $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ generated by the singular vector $`G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$. This module is precisely $`N`$ above and hence $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })/N`$ is freely generated over $`[]`$ by $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ and $`G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}`$. We claim that $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })/N`$ is irreducible. The even part of $`K(1,2)_+`$ is isomorphic to the semi-direct sum of $`𝔙_+`$ (generated by $`L_n`$) and $`\stackrel{~}{𝔤}_+`$ (generated by $`J_n`$), where $`𝔤`$ is the one-dimensional Lie algebra. We first consider $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })/N`$ as a module over the $`𝔙_+\stackrel{~}{𝔤}_+`$. The vectors $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ and $`G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ have $`(L_0,J_0)`$-weights $`(\mathrm{\Delta },\mathrm{\Lambda })`$ and $`(\mathrm{\Delta }+\frac{1}{2},\mathrm{\Lambda }+1)`$, respectively, and furthermore are both annihilated by $`L_n`$ and $`J_n`$, for $`n1`$. Now since $`2\mathrm{\Delta }+\mathrm{\Lambda }=0`$ and $`2\mathrm{\Delta }\mathrm{\Lambda }0`$, we have $`(\mathrm{\Delta },\mathrm{\Lambda })(0,0)`$ and $`(\mathrm{\Delta }+\frac{1}{2},\mathrm{\Lambda }+1)(0,0)`$. From this it follows that $`M_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })/N`$ as a module over $`𝔙_+\stackrel{~}{𝔤}_+`$ is a direct sum of two non-isomorphic irreducible modules, namely $`[]v_{\mathrm{\Delta },\mathrm{\Lambda }}L_{𝔙_+\stackrel{~}{𝔤}_+}(\mathrm{\Delta },\mathrm{\Lambda })`$ and $`[]G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}L_{𝔙_+\stackrel{~}{𝔤}_+}(\mathrm{\Delta }+\frac{1}{2},\mathrm{\Lambda }+1)`$ (see Section 2 for notation). But we have
$$G_{\frac{1}{2}}^{}G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}=(2\mathrm{\Delta }\mathrm{\Lambda })v_{\mathrm{\Delta },\mathrm{\Lambda }}0,$$
which implies that as a $`K(1,2)_+`$-module $`L_{𝔑_+^2}(\mathrm{\Delta },\mathrm{\Lambda })`$ is irreducible.
The case when $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$ and $`2\mathrm{\Delta }+\mathrm{\Lambda }0`$ is completely analogous and we leave it to the reader.
Finally in the case when $`\mathrm{\Delta }=\mathrm{\Lambda }=0`$, both $`G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ and $`G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ are proper singular vectors. Now the submodule in $`M_{𝔑_+^2}(0,0)`$ generated by these two vectors contains $`[G_{\frac{1}{2}}^+,G_{\frac{1}{2}}^{}]v_{\mathrm{\Delta },\mathrm{\Lambda }}=2v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, and hence has codimension $`1`$ over $``$. So the resulting quotient is the trivial module. ∎
It follows that every finite irreducible module over the $`N=2`$ conformal superalgebra is of the form $`L_{𝔑^2}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })`$, where $`\alpha ,\mathrm{\Delta },\mathrm{\Lambda }`$. We will write down explicit formulas for the action of the conformal superalgebra on such irreducible modules in the generating series form. Since we have already explained in Section 2 how such formulas can be obtained in general, we will omit the proofs.
In the case when $`2\mathrm{\Delta }\pm \mathrm{\Lambda }0`$ the module $`L_{𝔑^2}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })`$ is generated freely over $`[]`$ by two even vectors $`v,v^+`$ and two odd vectors $`v^+,v^{}`$. We have the following action on the generators:
$`L_\lambda v=(+\alpha +\mathrm{\Delta }\lambda )v,L_\lambda v^\pm =(+\alpha +(\mathrm{\Delta }+{\displaystyle \frac{1}{2}})\lambda )v^\pm ,`$
$`L_\lambda v^+=(+\alpha +(\mathrm{\Delta }+1)\lambda )v^++(\mathrm{\Delta }+{\displaystyle \frac{\mathrm{\Lambda }}{2}})\lambda ^2v,`$
$`J_\lambda v=\mathrm{\Lambda }v,J_\lambda v^\pm =(\mathrm{\Lambda }\pm 1)v^\pm ,J_\lambda v^+=\mathrm{\Lambda }v^++(2\mathrm{\Delta }+\mathrm{\Lambda })\lambda v,`$
$`G_\lambda ^\pm v=v^\pm ,G_\lambda ^+v^+=G_\lambda ^{}v^{}=0,G_\lambda ^+v^{}=v^++(2\mathrm{\Delta }+\mathrm{\Lambda })\lambda v,`$
$`G_\lambda ^+v^+=\lambda (2\mathrm{\Delta }+\mathrm{\Lambda })v^+,G_\lambda ^{}v^+=(2+2\alpha +\lambda (2\mathrm{\Delta }\mathrm{\Lambda }))vv^+,`$
$`G_\lambda ^{}v^+=(2+2\alpha +(2\mathrm{\Delta }+2\mathrm{\Lambda })\lambda )v^{}.`$
In the case when $`2\mathrm{\Delta }+\mathrm{\Lambda }=0`$ but $`2\mathrm{\Delta }\mathrm{\Lambda }0`$ the module $`L_{𝔑^2}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })`$ is generated freely over $`[]`$ by one even vector $`v`$ and one odd vector $`v^+`$. The action is then given by
$`L_\lambda v=(+\alpha +\mathrm{\Delta }\lambda )v,L_\lambda v^+=(+\alpha +(\mathrm{\Delta }+{\displaystyle \frac{1}{2}})\lambda )v^+,`$
$`J_\lambda v=2\mathrm{\Delta }v,J_\lambda v^+=(2\mathrm{\Delta }+1)v^+,G_\lambda ^+v=v^+,G_\lambda ^+v^+=0,`$
$`G_\lambda ^{}v=0,G_\lambda ^{}v^+=(2+2\alpha +4\mathrm{\Delta }\lambda )v.`$
In the case $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$ but $`2\mathrm{\Delta }+\mathrm{\Lambda }0`$ the module $`L_{𝔑^2}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })`$ is generated freely over $`[]`$ by one even vector $`v`$ and one odd vector $`v^{}`$ with action:
$`L_\lambda v=(+\alpha +\mathrm{\Delta }\lambda )v,L_\lambda v^{}=(+\alpha +(\mathrm{\Delta }+{\displaystyle \frac{1}{2}})\lambda )v^{},`$
$`J_\lambda v=2\mathrm{\Delta }v,J_\lambda v^{}=(2\mathrm{\Delta }1)v^{},G_\lambda ^+v=0,`$
$`G_\lambda ^+v^{}=(2+2\alpha +4\mathrm{\Delta }\lambda )v,G_\lambda ^{}v=v^{},G_\lambda ^{}v^{}=0.`$
Finally $`L_{𝔑^2}(\alpha ,0,0)`$ is the one-dimensional trivial module on which $``$ acts as the scalar $`\alpha `$.
###### Remark 4.1.
We note that the formulas above are obtained by first putting $`v=v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, $`v^\pm =G_{\frac{1}{2}}^\pm v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ and $`v^+=G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ and then compute the action of the operators $`L_n`$, $`J_m`$ and $`G_r^\pm `$, for $`n1`$, $`m0`$ and $`r\frac{1}{2}`$ on these vector. Translation into the language of conformal modules is an easy task using these formulas and we will omit this. Of course the parity of the vectors $`v,v^\pm ,v^+`$ in all the examples above can be reversed. Finally we note that the adjoint module is isomorphic to $`L_{𝔑^2}(0,1,0)`$.
## 5. Finite irreducible Modules over the $`N=3`$ conformal superalgebra
The $`N=3`$ superconformal algebra is the formal distribution Lie superalgebra $`K(1,3)`$. Letting $`\xi _1,\xi _2,\xi _3`$ be the three odd indeterminates $`K(1,3)`$ is spanned over $``$ by the following basis elements ($`n`$ and $`r\frac{1}{2}+`$):
$`L_n={\displaystyle \frac{t^{n+1}}{2}},H_n=2i\xi _1\xi _2t^n,E_n=(\xi _1\xi _3i\xi _2\xi _3)t^n,F_n=(\xi _1\xi _3i\xi _2\xi _3)t^n,`$
$`\mathrm{\Psi }_r=\xi _1\xi _2\xi _3t^{r\frac{1}{2}}h_r=2i\xi _3t^{r+\frac{1}{2}},e_r=(i\xi _1\xi _2)t^{r+\frac{1}{2}},f_r=(i\xi _1+\xi _2)t^{r+\frac{1}{2}}.`$
Let $`\{H,E,F\}`$ denote the standard basis of the Lie algebra $`sl_2`$ and $`\{h,e,f\}`$ denote the standard basis of its adjoint module. Furthermore we let $`(|)`$ denote the non-degenerate invariant symmetric bilinear form on $`sl_2`$ with $`(H|H)=2`$. Keeping this notation in mind the commutation relations of $`K(1,3)`$ are then given as follows (where $`X,Y=H,E,F`$ and $`x,y=h,e,f`$):
$`[L_m,L_n]=(mn)L_{m+n},[L_m,X_n]=nX_{m+n},[L_m,x_r]=({\displaystyle \frac{m}{2}}r)x_{m+r},`$
$`[L_m,\mathrm{\Psi }_r]=({\displaystyle \frac{m}{2}}r)\mathrm{\Psi }_{m+r},[X_m,Y_n]=[X,Y]_{m+n},[X_m,\mathrm{\Psi }_r]=0,`$
$`[X_m,y_r]=[X,y]_{m+r}+2m(X|Y)\mathrm{\Psi }_{m+r},[x_r,\mathrm{\Psi }_s]=X_{r+s},[\mathrm{\Psi }_r,\mathrm{\Psi }_s]=0,`$
$`[x_r,y_s]=(rs)[X,Y]_{r+s}4(X|Y)L_{r+s},`$
where $`m,n`$ and $`r,s\frac{1}{2}+`$. Above we have written $`[X,y]`$ for the action of $`X`$ on $`y`$. The eight formal distributions generating this algebra are given by $`L(z)=_nL_nz^{n2}`$, $`X(z)=_nX_nz^{n1}`$, $`x(z)=_{r\frac{1}{2}+}x_rz^{r\frac{3}{2}}`$ and $`\mathrm{\Psi }(z)=_{r\frac{1}{2}+}\mathrm{\Psi }_rz^{r\frac{1}{2}}`$. The corresponding operator product expansions of these fields are easily derived from (2.3), and so we will omit them.
The annihilation subalgebra $`K(1,3)_+`$ is equipped with a $`\frac{1}{2}`$-gradation of depth $`1`$, i.e. $`K(1,3)_+=𝔤=_{j1}𝔤_j`$, $`j\frac{1}{2}`$, and its $`0`$-th graded component $`𝔤_0`$ is isomorphic to a copy of $`gl_2sl_2L_0`$, with $`H_0`$, $`E_0`$ and $`F_0`$ providing the standard basis for the copy of $`sl_2`$.
Let $`U^{\mathrm{\Delta },\mathrm{\Lambda }}`$ be the finite-dimensional irreducible $`sl_2`$-module of highest weight $`\mathrm{\Lambda }_+`$ on which $`L_0`$ acts as the scalar $`\mathrm{\Delta }`$. We let $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ be a highest weight vector in $`U^{\mathrm{\Delta },\mathrm{\Lambda }}`$. We extend $`U^{\mathrm{\Delta },\mathrm{\Lambda }}`$ to a module over the subalgebra $`_0=_{j0}𝔤_j`$ in a trivial way and call this $`_0`$-module also $`U^{\mathrm{\Delta },\mathrm{\Lambda }}`$. By Theorem 3.1 every finite irreducible $`𝔤`$-module is a homomorphic image of $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })=\mathrm{Ind}__0^𝔤U^{\mathrm{\Delta },\mathrm{\Lambda }}`$ and furthermore $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ has a unique maximal submodule $`N`$, whose irreducible quotient we denote by $`L_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$.
Note that $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ as a module over $`sl_2`$ is a direct sum of infinitely many copies of finite-dimensional irreducible representations. Since $``$ commutes with $`E_0`$, the $`E_0`$-invariants $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$ is a $`[]`$-submodule of $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$, and hence is a free $`[]`$-module. We can write down explicitly formulas for a $`[]`$-basis of $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$. In the case when $`\mathrm{\Lambda }2`$ the following is a $`[]`$-basis:
$`a_1=v_{\mathrm{\Delta },\mathrm{\Lambda }},a_2=e_{\frac{1}{2}}v_{\mathrm{\Delta },\mathrm{\Lambda }},a_3=(\mathrm{\Lambda }h_{\frac{1}{2}}+2e_{\frac{1}{2}}F_0)v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_4=((\mathrm{\Lambda }1)(\mathrm{\Lambda }f_{\frac{1}{2}}h_{\frac{1}{2}}F_0)e_{\frac{1}{2}}F_0^2)v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_5=e_{\frac{1}{2}}h_{\frac{1}{2}}v_{\mathrm{\Delta },\mathrm{\Lambda }},a_6=(\mathrm{\Lambda }e_{\frac{1}{2}}f_{\frac{1}{2}}e_{\frac{1}{2}}h_{\frac{1}{2}}F_0)v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_7=((\mathrm{\Lambda }1)(\mathrm{\Lambda }h_{\frac{1}{2}}f_{\frac{1}{2}}+4F_0+2e_{\frac{1}{2}}f_{\frac{1}{2}}F_0)e_{\frac{1}{2}}h_{\frac{1}{2}}F_0^2)v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_8=(e_{\frac{1}{2}}h_{\frac{1}{2}}f_{\frac{1}{2}}2h_{\frac{1}{2}})v_{\mathrm{\Delta },\mathrm{\Lambda }}.`$
The cases $`\mathrm{\Lambda }=0,1`$ are similar. Namely, when $`\mathrm{\Lambda }=1`$ we have $`a_4=a_7=0`$, and the remaining $`6`$ vectors form a $`[]`$-basis. Finally, in the case when $`\mathrm{\Lambda }=0`$, the terms $`a_3=a_4=a_6=a_7=0`$, so that $`M_{𝔑_+^3}(\mathrm{\Delta },0)^{E_0}`$ has rank $`4`$ over $`[]`$. (Actually the vectors $`a_i`$ depend on $`\mathrm{\Lambda }`$, so it would be more appropriate to write something like $`a_i^\mathrm{\Lambda }`$ instead of just $`a_i`$. However, from the context it will always be clear what $`\mathrm{\Lambda }`$ is, so that it is safe to adopt the simpler notation of $`a_i`$.) We denote the coefficient of $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ in the expression $`a_i`$ by $`u_i^\mathrm{\Lambda }`$ so that we have $`a_i=u_i^\mathrm{\Lambda }v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, for $`i=1,\mathrm{},8`$. For example $`u_1^\mathrm{\Lambda }=1`$, while $`u_2^\mathrm{\Lambda }=e_{\frac{1}{2}}`$ etc. We note that finding all vectors in $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$ above amounts essentially to decomposing tensor products of irreducible representations of $`sl_2`$ and then finding the corresponding highest weight vectors of the irreducible components.
Similarly we call a non-zero $`(L_0,H_0)`$-weight vector $`v`$ in $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ *singular* if $`vM_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$ and $`𝔤_jv=0`$, for all $`j>0`$. As before a singular vector is called *proper* if it is not a scalar multiple of $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$. Evidently $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ is irreducible if and only if $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ contains no proper singular vector. Our first objective is to classify singular vectors inside $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$.
###### Proposition 5.1.
Any proper singular vector in $`M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ is of the form ($`\alpha `$ with $`\alpha 0`$)
* $`\alpha a_2`$, if $`4\mathrm{\Delta }\mathrm{\Lambda }=0`$,
* $`\alpha a_4`$, if $`4\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }2`$,
* $`\alpha a_6`$, if $`4\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }=1`$.
###### Remark 5.1.
The proof of the proposition is a straightforward, albeit a tedious, calculation. We will not give the details here, but instead just point out that a weight vector $`vM_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$ is singular if and only if $`f_{\frac{1}{2}}`$ and $`\mathrm{\Psi }_{\frac{1}{2}}`$ annihilates $`v`$. This fact simplifies the calculation significantly.
From Proposition 5.1 one obtains immediately the following.
###### Corollary 5.1.
Suppose that $`(\mathrm{\Delta },\mathrm{\Lambda })`$ does not satisfy either $`4\mathrm{\Delta }\mathrm{\Lambda }=0`$ or $`4\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }1`$. Then $`L_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })=M_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ is an irreducible $`K(1,3)_+`$-module of rank $`8\mathrm{\Lambda }+8`$ over $`[]`$.
###### Proposition 5.2.
Suppose that $`4\mathrm{\Delta }\mathrm{\Lambda }=0`$. Then $`L_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ is a free $`[]`$-module of rank $`4\mathrm{\Lambda }`$.
###### Proof.
By Proposition 5.1 $`a_2`$ is a singular vector in $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })`$ of $`H_0`$-weight $`\mathrm{\Lambda }+2`$. Consider $`N`$, the $`𝔤`$-submodule generated by $`a_2`$. Then we have $`N=U(𝔤_{})V_2`$, where $`V_2`$ is the irreducible $`sl_2`$-submodule generated by $`a_2`$. Note that the map $`v_{\frac{\mathrm{\Lambda }}{4}+\frac{1}{2},\mathrm{\Lambda }+2}a_2`$ extends uniquely to an epimorphism of $`K(1,3)_+`$-modules from $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4}+\frac{1}{2},\mathrm{\Lambda }+2)`$ to $`N`$. In particular it is an $`sl_2`$-module epimorphism. Now both modules are completely reducible $`sl_2`$-modules and hence this map sends $`E_0`$-invariants onto $`E_0`$-invariants. Since $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4}+\frac{1}{2},\mathrm{\Lambda }+2)^{E_0}`$ is generated over $`[]`$ by $`\{u_i^{\mathrm{\Lambda }+2}v_{\frac{\mathrm{\Lambda }}{4}+\frac{1}{2},\mathrm{\Lambda }+2}|1i8\}`$, it follows that $`N^{E_0}`$ is generated over $`[]`$ by $`\{u_i^{\mathrm{\Lambda }+2}a_2|1i8\}`$. Now $`N^{E_0}`$ is a $`[]`$-submodule of $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })`$, since $`[,E_0]=0`$. Thus it is a free $`[]`$-submodule generated by $`\{u_i^{\mathrm{\Lambda }+2}a_2|1i8\}`$. We compute
$`u_1^{\mathrm{\Lambda }+2}a_2=a_2,u_2^{\mathrm{\Lambda }+2}a_2=0,u_3^{\mathrm{\Lambda }+2}a_2=(\mathrm{\Lambda }+4)a_5,`$
$`u_4^{\mathrm{\Lambda }+2}a_2=(\mathrm{\Lambda }+3)a_64(\mathrm{\Lambda }+1)(\mathrm{\Lambda }+3)a_1,u_5^{\mathrm{\Lambda }+2}a_2=0,`$
$`u_6^{\mathrm{\Lambda }+2}a_2=4(\mathrm{\Lambda }+3)a_2,u_7^{\mathrm{\Lambda }+2}a_2=(\mathrm{\Lambda }+3)(\mathrm{\Lambda }+2)a_8+2(\mathrm{\Lambda }+3)a_3,`$
$`u_8^{\mathrm{\Lambda }+2}a_2=2a_5.`$
By inspection it is clear that the following is a set of $`[]`$-generators for $`N^{E_0}`$.
$$S^\mathrm{\Lambda }=\{a_2,a_5,a_6+4(\mathrm{\Lambda }+1)a_1,a_8+\frac{2}{\mathrm{\Lambda }+2}a_3\}.$$
First consider the case when $`\mathrm{\Lambda }2`$. It follows from the description of $`S^\lambda `$ above that $`\{a_1,a_3,a_4,a_7\}`$ is a $`[]`$-basis for the $`E_0`$-invariants of the quotient $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$. Since $`a_1`$ and $`a_3`$ both have $`H_0`$-weight $`\mathrm{\Lambda }`$, they generate two copies of the irreducible $`sl_2`$-module of dimension $`\mathrm{\Lambda }+1`$. On the other hand $`a_4`$ and $`a_7`$ both have weight $`\mathrm{\Lambda }2`$, and so they generate two copies of the irreducible $`sl_2`$-module of dimension $`\mathrm{\Lambda }1`$. Thus $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$ is a free $`[]`$-module of rank $`2(\mathrm{\Lambda }+1)+2(\mathrm{\Lambda }1)=4\mathrm{\Lambda }`$. So in order to complete the proof it remains to show that $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$ is irreducible.
Note that $`L_n`$, $`n1`$, together with $`E_0,H_0,F_0`$ generate a copy of $`𝔙_+sl_2`$ and so we may consider $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$ as a module over $`𝔙_+sl_2`$. By parity consideration $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$ is a direct sum of two $`(𝔙_+sl_2)`$-modules, namely $`(M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N)_{\overline{0}}=[]V_1+[]V_7`$ and $`(M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N)_{\overline{1}}=[]V_3+[]V_4`$, where $`V_i`$ is the irreducible $`sl_2`$-module generated by $`a_i`$. It is subject to a direct verification that $`L_n`$, for $`n1`$, annihilates the vectors $`a_1,a_3,a_4,a_7`$ (in fact one only needs to check that $`L_1a_7=0`$, others being trivial) and hence $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$ as a $`𝔙_+sl_2`$-module is a direct sum the following four non-isomorphic irreducible modules: $`[]V_3L_{𝔙_+}(\frac{\mathrm{\Lambda }}{4}+\frac{1}{2})U^\mathrm{\Lambda }`$, $`[]V_4L_{𝔙_+}(\frac{\mathrm{\Lambda }}{4}+\frac{1}{2})U^{\mathrm{\Lambda }2}`$, $`[]V_1L_{𝔙_+}(\frac{\mathrm{\Lambda }}{4})U^\mathrm{\Lambda }`$ and $`[]V_7L_{𝔙_+}(\frac{\mathrm{\Lambda }}{4}+1)U^{\mathrm{\Lambda }2}`$, where we denote by $`U^\mu `$ the irreducible $`sl_2`$-module of highest weight $`\mu `$. Now we compute
(5.1) $`\mathrm{\Psi }_{\frac{1}{2}}a_3=\mathrm{\Lambda }(\mathrm{\Lambda }+2)a_1,f_{\frac{1}{2}}a_4=(2\mathrm{\Lambda }+2)F_0^2a_1`$
$`E_1a_7=2\mathrm{\Lambda }(\mathrm{\Lambda }1)(2\mathrm{\Lambda }+2)a_1,`$
from which it follows that we may go from each irreducible $`𝔙_+sl_2`$-component of $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$ to the irreducible component containing the highest weight vectors, and hence the module $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$ is irreducible.
Now if $`\mathrm{\Lambda }=1`$ the vectors $`a_4`$ and $`a_7`$ are both zero. Therefore the quotient $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N=[]V_1[]V_3`$. But then the first identity in (5.1) shows that $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$ is irreducible. The rank of $`L_{𝔑_+^3}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })`$ is then $`2(\mathrm{\Lambda }+1)=4\mathrm{\Lambda }`$ in the case $`\mathrm{\Lambda }=1`$.
Finally when $`\mathrm{\Lambda }=0`$, the vectors $`a_3,a_4,a_6,a_7=0`$, so that $`S^\mathrm{\Lambda }`$ reduces to $`\{a_2,a_5,a_1,a_8\}`$. Therefore $`M_{𝔑_+^3}(0,0)/N=a_1`$ is the trivial module and so has rank $`0`$. ∎
###### Proposition 5.3.
Suppose that $`4\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }1`$. Then $`L_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ is a free $`[]`$-module of rank $`4\mathrm{\Lambda }+8`$.
###### Proof.
By Proposition 5.1 $`a_4`$ is a singular vector of $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })`$ of $`H_0`$-weight $`\mathrm{\Lambda }2`$. Let $`N`$ denote the $`𝔤`$-submodule generated by $`a_4`$.
Consider first the case $`\mathrm{\Lambda }4`$. As in the proof of Proposition 5.1 $`N^{E_0}`$ is a free $`[]`$-module generated over $`[]`$ by $`\{u_i^{\mathrm{\Lambda }2}a_4|1i8\}`$. We compute
$`u_1^{\mathrm{\Lambda }2}a_4=a_4,u_2^{\mathrm{\Lambda }2}a_4=(\mathrm{\Lambda }1)a_6,u_3^{\mathrm{\Lambda }2}a_4=(\mathrm{\Lambda }2)a_7,u_4^{\mathrm{\Lambda }2}a_4=0,`$
$`u_5^{\mathrm{\Lambda }2}a_4=\mathrm{\Lambda }(\mathrm{\Lambda }1)a_8+2(\mathrm{\Lambda }1)a_3,u_6^{\mathrm{\Lambda }2}a_4=0,u_7^{\mathrm{\Lambda }2}a_4=0,`$
$`u_8^{\mathrm{\Lambda }2}a_4=2a_7.`$
This implies that the set $`S^\mathrm{\Lambda }=\{a_4,(\mathrm{\Lambda }1)a_6,(\mathrm{\Lambda }2)a_7,\mathrm{\Lambda }(\mathrm{\Lambda }1)a_8+2(\mathrm{\Lambda }1)a_3,a_7\}`$ generates $`N^{E_0}`$ over $`[]`$ and so $`\{a_1,a_2,a_3,a_5\}`$ is a $`[]`$-basis for $`(M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N)^{E_0}`$ in the case when $`\mathrm{\Lambda }4`$.
Next consider the case $`\mathrm{\Lambda }=3`$. In this case, letting $`N`$ be as before, $`N^{E_0}`$ is generated over $`[]`$ by $`\{u_i^{\mathrm{\Lambda }2}a_4|1i8,i4,7\}`$. Hence it follows from the above formulas that again $`\{a_1,a_2,a_3,a_5\}`$ is a $`[]`$-basis for $`(M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N)^{E_0}`$.
In the case when $`\mathrm{\Lambda }=2`$ we let $`N^{}`$ denote the module generated by $`a_4`$. It follows that the vectors $`\{u_1^{\mathrm{\Lambda }2}a_4,u_2^{\mathrm{\Lambda }2}a_4,u_5^{\mathrm{\Lambda }2}a_4,u_8^{\mathrm{\Lambda }2}a_4\}`$ generate $`N_{}^{}{}_{}{}^{E_0}`$ over $`[]`$ so that $`S^\mathrm{\Lambda }=\{a_4,a_6,a_8+a_3,a_7\}`$ generate $`N^{E_0}`$. Hence $`(M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N^{})^{E_0}`$ contains in addition a one-dimensional (over $``$) subspace spanned by $`a_7`$. However, $`a_7=0`$ in $`(M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N^{})`$ and hence it is a $`𝔤`$-invariant by Remark 2.1. In this case we set $`N=N^{}+a_7`$ so that the quotient module $`(M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N)^{E_0}`$ is again generated over $`[]`$ by $`\{a_1,a_2,a_3,a_5\}`$.
Now $`a_1`$ and $`a_3`$ have $`H_0`$-weight $`\mathrm{\Lambda }`$, while $`a_2`$ and $`a_5`$ have $`H_0`$-weight $`\mathrm{\Lambda }+2`$. Thus $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N`$ has rank $`2(\mathrm{\Lambda }+1)+2(\mathrm{\Lambda }+3)=4\mathrm{\Lambda }+8`$ over $`[]`$. So it remains to show that $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N`$ is irreducible.
Again we consider $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N`$ as a module over $`𝔙_+sl_2`$. It is easy to check that $`L_n`$, $`n1`$, annihilates $`a_1,a_2,a_3,a_5`$. (Again one really only needs to check that $`L_1a_5=0`$.) Thus it follows that in the case of $`\mathrm{\Lambda }3`$ that $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N`$ is a direct sum of the following four non-isomorphic irreducible $`𝔙_+sl_2`$-modules: $`[]V_1L_{𝔙_+}(\frac{\mathrm{\Lambda }+2}{4})U^\mathrm{\Lambda }`$, $`[]V_2L_{𝔙_+}(\frac{\mathrm{\Lambda }}{4})U^{\mathrm{\Lambda }+2}`$, $`[]V_3L_{𝔙_+}(\frac{\mathrm{\Lambda }}{4})U^\mathrm{\Lambda }`$ and $`[]V_5L_{𝔙_+}(\frac{\mathrm{\Lambda }2}{4})U^{\mathrm{\Lambda }+2}`$, where as before $`U^\mu `$ stands for the irreducible $`sl_2`$-module of highest weight $`\mu `$ and $`V_i`$ is the $`sl_2`$-submodule generated by the vector $`a_i`$. Now we compute
(5.2)
$$f_{\frac{1}{2}}a_2=2(\mathrm{\Lambda }+1)a_1,\mathrm{\Psi }_{\frac{1}{2}}a_3=\mathrm{\Lambda }(\mathrm{\Lambda }+2)a_1,F_1a_5=4(\mathrm{\Lambda }+1)a_1,$$
from which again it follows that we may go from any irreducible $`𝔙_+sl_2`$-component of $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N`$ to the component containing the highest weight vectors, and hence $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N`$ is irreducible.
As for the case $`\mathrm{\Lambda }=2`$ we have $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N`$ as a $`𝔙_+sl_2`$-module is also a direct sum of the $`[]V_1[]V_2[]V_3[]V_5`$. The first three modules, as in the case of $`\mathrm{\Lambda }3`$ are irreducible. However, $`[]V_5`$ contains a unique irreducible submodule generated by the vector $`a_5`$, which is isomorphic to $`L_{𝔙_+}(1)U^{\mathrm{\Lambda }+2}`$. But then (5.2) together with the fact that
$$F_2a_1=24a_1$$
shows that $`M_{𝔑_+^3}(\frac{\mathrm{\Lambda }+2}{4},\mathrm{\Lambda })/N`$ is irreducible in this case as well.
In the case when $`\mathrm{\Lambda }=1`$ we have by Proposition 5.1 that $`a_6`$ is the unique (up to a scalar) singular vector inside $`M_{𝔑_+^3}(\frac{3}{4},1)`$. Let $`N`$ denote the $`𝔤`$-submodule generated by $`a_6`$. Since $`a_6`$ has $`H_0`$-weight $`1`$, $`N^{E_0}`$ is the free $`[]`$-module generated by $`\{u_i^\mathrm{\Lambda }a_6|1i8,i4,7\}`$. We have
$`u_1^\mathrm{\Lambda }a_6=a_6,u_2^\mathrm{\Lambda }a_6=0,u_3^\mathrm{\Lambda }a_6=3a_86a_3,`$
$`u_5^\mathrm{\Lambda }a_6=0,u_6^\mathrm{\Lambda }a_6=8a_6,u_8^\mathrm{\Lambda }a_6=2a_84^2a_3,`$
from which it follows that $`N^{E_0}`$ is generated over $`[]`$ by $`S^\mathrm{\Lambda }=\{a_6,a_8+2a_3\}`$. Since $`a_4=a_7=0`$ in this situation, we see that $`(M_{𝔑_+^3}(\frac{3}{4},1)/N)^{E_0}`$ is generated over $`[]`$ by the vectors $`a_1,a_2,a_3,a_5`$, just as in the case $`\mathrm{\Lambda }2`$. Now the exact same argument as in the $`\mathrm{\Lambda }2`$ case shows that $`M_{𝔑_+^3}(\frac{3}{4},1)/N`$ is irreducible and has rank $`4\mathrm{\Lambda }+8`$ over $`[]`$. ∎
We summarize the work in this section in the following theorem.
###### Theorem 5.1.
The modules $`L_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$, for $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }_+`$, form a complete list of non-isomorphic finite (over $`[]`$) irreducible $`K(1,3)_+`$-modules. Furthermore $`L_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })`$ as a $`[]`$-module has rank
* $`4\mathrm{\Lambda }`$, in the case $`4\mathrm{\Delta }\mathrm{\Lambda }=0`$,
* $`4\mathrm{\Lambda }+8`$, the case $`4\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }1`$.
* $`8\mathrm{\Lambda }+8`$, in all other cases.
Furthermore the $`[]`$-rank of $`L_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })_{\overline{0}}`$ equals the $`[]`$-rank of $`L_{𝔑_+^3}(\mathrm{\Delta },\mathrm{\Lambda })_{\overline{1}}`$ in all cases.
###### Remark 5.2.
Translating the above theorem back into the languages of modules over conformal superalgebras and of conformal modules is now a straightforward task. We thus have proved that all finite irreducible modules over the $`N=3`$ conformal superalgebra are of the form $`L_{𝔑^3}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })`$, where $`\alpha ,\mathrm{\Delta }`$ and $`\mathrm{\Lambda }_+`$. The definition of these modules and also the action of the $`N=3`$ conformal superalgebra on them are quite easy to obtain from our explicit description of a $`[]`$-basis of these modules. To do so would however take up quite a significant portion of space, and thus we leave this task to the interested reader. We only remark that the adjoint module is isomorphic to $`L_{𝔑^3}(0,\frac{1}{2},0)`$.
## 6. Finite irreducible Modules over the “small” $`N=4`$ conformal superalgebra
The “small” $`N=4`$ superconformal algebra is the following subalgebra of $`K(1,4)`$: Let $`\xi _1,\xi _2,\xi _3,\xi _4`$ denote four odd indeterminates generating the Grassmann superalgebra $`\mathrm{\Lambda }(4)`$. For a monomial $`\xi _I`$ in $`\mathrm{\Lambda }(4)`$ we let $`\xi _I^{}`$ be its Hodge dual, i.e. the unique monomial in $`\mathrm{\Lambda }(4)`$ such that $`\xi _I\xi _I^{}=\xi _1\xi _2\xi _3\xi _4`$. Then the small $`N=4`$ superconformal algebra is isomorphic to any of the following two subalgebras in $`K(1,4)`$ spanned by the following basis elements ($`n`$, $`r\frac{1}{2}+`$, $`\beta ^2=1`$) :
$`L_n^\beta ={\displaystyle \frac{1}{2}}(t^{n+1}+\beta n(n+1)\xi _1\xi _2\xi _3\xi _4t^{n1}),`$
$`H_n^\beta =i(\xi _1\xi _2\beta \xi _3\xi _4)t^n,`$
$`E_n^\beta ={\displaystyle \frac{1}{2}}(\xi _1\xi _3\beta \xi _2\xi _4i\xi _2\xi _3+i\beta \xi _1\xi _4)t^n,`$
$`F_n^\beta ={\displaystyle \frac{1}{2}}(\xi _1\xi _3+\beta \xi _2\xi _4i\xi _2\xi _3+i\beta \xi _1\xi _4)t^n,`$
$`G_r^{+\beta }={\displaystyle \frac{1}{\sqrt{2}}}((\xi _3+i\xi _4)t^{r+\frac{1}{2}}\beta (r+{\displaystyle \frac{1}{2}})(\xi _3^{}+i\xi _4^{})t^{r\frac{1}{2}}),`$
$`G_r^{++\beta }={\displaystyle \frac{1}{\sqrt{2}}}((\xi _1+i\xi _2)t^{r+\frac{1}{2}}\beta (r+{\displaystyle \frac{1}{2}})(\xi _1^{}+i\xi _2^{})t^{r\frac{1}{2}}),`$
$`G_r^{+\beta }={\displaystyle \frac{1}{\sqrt{2}}}((\xi _3i\xi _4)t^{r+\frac{1}{2}}\beta (r+{\displaystyle \frac{1}{2}})(\xi _3^{}i\xi _4^{})t^{r\frac{1}{2}}),`$
$`G_r^\beta ={\displaystyle \frac{1}{\sqrt{2}}}((i\xi _2\xi _1)t^{r+\frac{1}{2}}\beta (r+{\displaystyle \frac{1}{2}})(i\xi _2^{}\xi _1^{})t^{r\frac{1}{2}}).`$
As before let $`\{H,E,F\}`$ denote the standard basis of the Lie algebra $`sl_2`$ and $`\{G^{++},G^+\}`$ denote the standard basis of its standard module, i.e. $`HG^{++}=G^{++}`$, $`HG^+=G^+`$, $`EG^{++}=FG^+=0`$, $`FG^{++}=G^+`$ and $`EG^+=G^{++}`$. Likewise $`\{G^+,G^{}\}`$ also denotes a copy of the standard basis of the standard $`sl_2`$-module with actions $`HG^+=G^+`$, $`HG^{}=G^{}`$, $`EG^+=FG^{}=0`$, $`FG^+=G^{}`$ and $`EG^{}=G^+`$. With this notation in mind the commutation relations are then given as follows (where $`X,Y=H,E,F`$ and $`x,y=G^{++},G^+,G^+,G^{}`$):
$`[L_m^\beta ,L_n^\beta ]=(mn)L_{m+n}^\beta ,[L_m^\beta ,X_n^\beta ]=nX_{m+n}^\beta ,[L_m^\beta ,x_r^\beta ]=({\displaystyle \frac{m}{2}}r)x_{m+r}^\beta ,`$
$`[X_m^\beta ,Y_n^\beta ]=[X,Y]_{m+n}^\beta ,[X_m^\beta ,y_r^\beta ]=(Xy)_{m+r}^\beta ,[x_r^\beta ,x_s^\beta ]=0,`$
$`[G_r^{++\beta },G_s^{+\beta }]=(rs)(1+\beta )E_{r+s}^\beta ,[G_r^{++\beta },G_s^{+\beta }]=(rs)(1\beta )E_{r+s}^\beta ,`$
$`[G_r^{++\beta },G_s^\beta ]=(rs)H_{r+s}^\beta 2L_{r+s}^\beta ,[G_r^{+\beta },G_s^{+\beta }]=(rs)\beta H_{r+s}^\beta +2L_{r+s}^\beta ,`$
$`[G_r^{+\beta },G_s^\beta ]=(rs)(1\beta )F_{r+s}^\beta ,[G_r^{+\beta },G_s^\beta ]=(rs)(1+\beta )F_{r+s}^\beta ,`$
where $`m,n`$ and $`r,s\frac{1}{2}+`$. The eight formal distributions generating this algebra are given by $`L^\beta (z)=_nL_n^\beta z^{n2}`$, $`X^\beta (z)=_nX_n^\beta z^{n1}`$, $`x^\beta (z)=_{r\frac{1}{2}+}x_r^\beta z^{r\frac{3}{2}}`$. The operator product expansions of these fields are easily derived using (2.3).
We will denote the “small” $`N=4`$ superconformal algebra simply by $`SK(1,4)`$ and assume for the rest of this section that we have chosen its realization as the subalgebra of $`K(1,4)`$ with $`\beta =1`$ for future computational purposes. For simplicity we will drop the superscript $`\beta `$ and write $`L_n`$ for $`L_n^\beta `$ etc. when we mean $`\beta =1`$.
The annihilation subalgebra $`𝔤=SK(1,4)_+`$ of $`SK(1,4)`$ is equipped with a $`\frac{1}{2}`$-gradation of depth $`1`$, i.e. $`𝔤=_{j1}𝔤_j`$, $`j\frac{1}{2}`$, and its $`0`$-th graded component $`𝔤_0`$ is isomorphic to a copy of $`gl_2sl_2L_0`$, with $`H_0`$, $`E_0`$ and $`F_0`$ providing the standard basis of the copy of $`sl_2`$. Again we let $`U^{\mathrm{\Delta },\mathrm{\Lambda }}`$ be the finite-dimensional irreducible $`sl_2`$-module of highest weight $`\mathrm{\Lambda }_+`$ on which $`L_0`$ acts as the scalar $`\mathrm{\Delta }`$ and $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ be a highest weight vector in $`U^{\mathrm{\Delta },\mathrm{\Lambda }}`$. As in the case of $`K(1,3)_+`$, we may extend $`U^{\mathrm{\Delta },\mathrm{\Lambda }}`$ to a module over the subalgebra $`_0=_{j0}𝔤_j`$ trivially and call this $`_0`$-module also $`U^{\mathrm{\Delta },\mathrm{\Lambda }}`$. Again Theorem 3.1 tells us that every finite irreducible $`𝔤`$-module is the quotient of $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })=\mathrm{Ind}__0^𝔤U^{\mathrm{\Delta },\mathrm{\Lambda }}`$ by its unique maximal submodule, for some $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }_+`$. We denote the unique irreducible quotient by $`L_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ so that every finite irreducible $`SK(1,4)_+`$-module is of the form $`L_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$, for $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }_+`$.
Now $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ is completely reducible as a module over $`sl_2=H_0+E_0+F_0`$, and the subspace of $`E_0`$-invariants $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$ is a free $`[]`$-submodule of $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ due to $`[E_0,]=0`$. We write down explicit formulas for a $`[]`$-basis of $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$, which in the case when $`\mathrm{\Lambda }2`$ takes the following form:
$`a_1=v_{\mathrm{\Delta },\mathrm{\Lambda }},a_2=G_{\frac{1}{2}}^{++}v_{\mathrm{\Delta },\mathrm{\Lambda }},a_3=G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }},a_4=(\mathrm{\Lambda }G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{++}F_0)v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_5=(\mathrm{\Lambda }G_{\frac{1}{2}}^{}+G_{\frac{1}{2}}^+F_0)v_{\mathrm{\Delta },\mathrm{\Lambda }},a_6=G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{++}v_{\mathrm{\Delta },\mathrm{\Lambda }},a_7=G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_8=G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }},a_9=(G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^{})v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_{10}=(\mathrm{\Lambda }G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^{}+G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+F_0)v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_{11}=((\mathrm{\Lambda }1)(\mathrm{\Lambda }G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}+G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+F_0+G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^{}F_0)G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+F_0^2)v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_{12}=G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda }},a_{13}=G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_{14}=G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{++}(\mathrm{\Lambda }G_{\frac{1}{2}}^{}+G_{\frac{1}{2}}^+F_0)v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_{15}=(\mathrm{\Lambda }G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}+G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}F_0)v_{\mathrm{\Delta },\mathrm{\Lambda }},`$
$`a_{16}=G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}v_{\mathrm{\Delta },\mathrm{\Lambda }}.`$
Now in the case when $`\mathrm{\Lambda }=1`$ we have $`a_{11}=0`$ so that the remaining $`15`$ vectors form a $`[]`$-basis for $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$, while in the case when $`\mathrm{\Lambda }=0`$ we have $`a_4=a_5=a_{10}=a_{11}=a_{14}=a_{15}=0`$, so that $`M_{𝔑_+^4}(\mathrm{\Delta },0)^{E_0}`$ has rank $`10`$ over $`[]`$. As in Section 5 we denote the coefficient of $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ in the expression $`a_i`$ by $`u_i^\mathrm{\Lambda }`$ so that we have $`a_i=u_i^\mathrm{\Lambda }v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, for $`i=1,\mathrm{},16`$.
*Singular vectors* are then defined to be non-zero $`(L_0,H_0)`$-weight vectors $`vM_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$ with $`𝔤_jv=0`$, for all $`j>0`$. Similarly we define *proper singular vectors*. Our approach is analogous to the one of Section 5, that is first to analyze singular vectors inside $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$. This is given by the following proposition, whose proof is again a straightforward calculation, which admittedly is rather tedious.
###### Proposition 6.1.
A complete list of proper singular vectors inside $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ is given by:
* $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$.
+ $`\alpha a_2+\beta a_3`$, $`(\alpha ,\beta )(0,0)`$,
+ $`\alpha a_8`$, $`\alpha 0`$.
* $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }2`$.
+ $`\alpha a_4+\beta a_5`$, $`(\alpha ,\beta )(0,0)`$,
+ $`\alpha a_{11}`$, $`\alpha 0`$.
* $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }=1`$.
+ $`\alpha a_4+\beta a_5`$, $`(\alpha ,\beta )(0,0)`$,
+ $`\alpha a_{14}+\beta (a_{15}2a_5)`$, $`(\alpha ,\beta )(0,0)`$,
+ $`\alpha (a_{16}2a_{10})`$, $`\alpha 0`$.
* $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }=0`$.
+ $`\alpha a_6+\beta a_7+\gamma (a_92a_1)`$, $`(\alpha ,\beta ,\gamma )(0,0,0)`$,
+ $`\alpha a_{13}+\beta (a_{12}+2a_2)`$, $`(\alpha ,\beta )(0,0)`$.
###### Remark 6.1.
We note that in order to check that a weight vector $`vM_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })^{E_0}`$ is singular, it is enough to check that $`v`$ is annihilated by $`F_1`$, $`G_{\frac{1}{2}}^+`$ and $`G_{\frac{1}{2}}^{}`$.
###### Corollary 6.1.
Suppose that $`(\mathrm{\Delta },\mathrm{\Lambda })`$ does not satisfy either $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$ or $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$. Then $`L_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })=M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ is an irreducible $`SK(1,4)_+`$-module of rank $`16\mathrm{\Lambda }+16`$ over $`[]`$.
###### Proposition 6.2.
Suppose that $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$. Then $`L_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ is a free $`[]`$-module of rank $`4\mathrm{\Lambda }`$.
###### Proof.
By Proposition 6.1 $`a_2`$ and $`a_3`$ are singular vectors in $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })`$. Consider $`N_2`$ and $`N_3`$, the $`𝔤`$-submodules generated by $`a_2`$ and $`a_3`$, respectively, and let $`N=N_2+N_3`$. Then we have $`N_2=U(𝔤_{})V_2`$ and $`N_3=U(𝔤_{})V_3`$, where $`V_2`$ and $`V_3`$ are the irreducible $`sl_2`$-submodules generated by $`a_2`$ and $`a_3`$, respectively. Let’s first compute $`N_2^{E_0}`$. Since the $`H_0`$-weight of $`a_2`$ is $`\mathrm{\Lambda }+1`$, we know that $`N_2^{E_0}`$ is a free $`[]`$-module generated over $`[]`$ by $`\{u_i^{\mathrm{\Lambda }+1}a_2|1i16\}`$. We have
$`u_1^{\mathrm{\Lambda }+1}a_2=a_2,u_2^{\mathrm{\Lambda }+1}a_2=0,u_3^{\mathrm{\Lambda }+1}a_2=a_8,u_4^{\mathrm{\Lambda }+1}a_2=(\mathrm{\Lambda }+2)a_6,`$
$`u_5^{\mathrm{\Lambda }+1}a_2=a_9a_{10}+2(\mathrm{\Lambda }+2)a_1,u_6^{\mathrm{\Lambda }+1}a_2=0,u_7^{\mathrm{\Lambda }+1}a_2=a_{13}2a_3,`$
$`u_8^{\mathrm{\Lambda }+1}a_2=0,u_9^{\mathrm{\Lambda }+1}a_2=a_{12}+2a_2,u_{10}^{\mathrm{\Lambda }+1}a_2=a_{12}+2(\mathrm{\Lambda }+2)a_2,`$
$`u_{11}^{\mathrm{\Lambda }+1}a_2=2(\mathrm{\Lambda }+2)a_4(\mathrm{\Lambda }+2)a_{14},u_{12}^{\mathrm{\Lambda }+1}a_2=0,u_{13}^{\mathrm{\Lambda }+1}a_2=2a_8,`$
$`u_{14}^{\mathrm{\Lambda }+1}a_2=2(\mathrm{\Lambda }+2)a_6,u_{15}^{\mathrm{\Lambda }+1}a_2=(\mathrm{\Lambda }+2)a_{16}+2(\mathrm{\Lambda }+1)a_92a_{10},`$
$`u_{16}^{\mathrm{\Lambda }+1}a_2=2a_{12}.`$
Next we find $`[]`$-generators of $`N_3^{E_0}`$. Similarly $`\{u_i^{\mathrm{\Lambda }+1}a_3|1i16\}`$ generates $`N_3^{E_0}`$ over $`[]`$:
$`u_1^{\mathrm{\Lambda }+1}a_3=a_3,u_2^{\mathrm{\Lambda }+1}a_3=a_8,u_3^{\mathrm{\Lambda }+1}a_3=0,u_4^{\mathrm{\Lambda }+1}a_3=(\mathrm{\Lambda }+1)a_9a_{10},`$
$`u_5^{\mathrm{\Lambda }+1}a_3=(\mathrm{\Lambda }+2)a_7,u_6^{\mathrm{\Lambda }+1}a_3=a_{12},u_7^{\mathrm{\Lambda }+1}a_3=0,u_8^{\mathrm{\Lambda }+1}a_3=0,`$
$`u_9^{\mathrm{\Lambda }+1}a_3=a_{13},u_{10}^{\mathrm{\Lambda }+1}a_3=(\mathrm{\Lambda }+2)a_{13},u_{11}^{\mathrm{\Lambda }+1}a_3=(\mathrm{\Lambda }+2)a_{15},`$
$`u_{12}^{\mathrm{\Lambda }+1}a_3=a_{16},u_{13}^{\mathrm{\Lambda }+1}a_3=0,u_{14}^{\mathrm{\Lambda }+1}a_3=(\mathrm{\Lambda }+2)a_{16},u_{15}^{\mathrm{\Lambda }+1}a_3=0,`$
$`u_{16}^{\mathrm{\Lambda }+1}a_3=0.`$
It follows that $`N^{E_0}`$ is generated over $`[]`$ by the following set
$$S^\mathrm{\Lambda }=\{a_2,a_3,a_6,a_7,a_8,a_92a_1,a_{10}2(\mathrm{\Lambda }+1)a_1,a_{12},a_{13},a_{14}2a_4,a_{15},a_{16}\}.$$
Suppose that $`\mathrm{\Lambda }2`$. From the description of $`S^\mathrm{\Lambda }`$ we see that $`\{a_1,a_4,a_5,a_{11}\}`$ is a $`[]`$-basis for the $`E_0`$-invariants of the quotient $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N`$. Since $`a_1`$ has $`H_0`$-weight $`\mathrm{\Lambda }`$, it generates a copy of the irreducible $`sl_2`$-module of dimension $`\mathrm{\Lambda }+1`$. Now $`a_4`$ and $`a_5`$ both have weight $`\mathrm{\Lambda }1`$, and so they generate two copies of the irreducible $`sl_2`$-module of dimension $`\mathrm{\Lambda }`$. Finally $`a_{11}`$ has weight $`\mathrm{\Lambda }2`$, and so it generates a copy of the irreducible $`sl_2`$-module of dimension $`\mathrm{\Lambda }1`$. Thus $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N`$ is a free $`[]`$-module of rank $`(\mathrm{\Lambda }+1)+2\mathrm{\Lambda }+(\mathrm{\Lambda }1)=4\mathrm{\Lambda }`$. So we need to show that $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N`$ is irreducible.
As in Section 5 $`L_n`$, $`n1`$, together with $`E_0,H_0,F_0`$ generate a copy of $`(𝔙_+sl_2)`$, which thus allow us to study the $`(𝔙_+sl_2)`$-module structure of $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N`$. By parity consideration $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N`$ is a direct sum of two modules, namely $`(M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N)_{\overline{0}}=[]V_1+[]V_{11}`$ and $`(M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N)_{\overline{1}}=[]V_4+[]V_5`$, where $`V_i`$ is the irreducible $`sl_2`$-module generated by $`a_i`$. We can easily check that $`L_n`$, for $`n1`$, annihilates the vectors $`a_1,a_4,a_5,a_{11}`$. (Again the only non-trivial part is to check that $`L_1a_{11}=0`$.) Thus $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N`$ as a $`𝔙_+sl_2`$-module is a direct sum of the following four irreducible modules: $`[]V_1L_{𝔙_+}(\frac{\mathrm{\Lambda }}{2})U^\mathrm{\Lambda }`$, $`[]V_4L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }1}`$, $`[]V_5L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }1}`$ and $`[]V_{11}L_{𝔙_+}(\frac{\mathrm{\Lambda }+2}{2})U^{\mathrm{\Lambda }2}`$, where as usual $`U^\mu `$ is the irreducible $`sl_2`$-module of highest weight $`\mu `$. Note that, contrary to the $`K(1,3)_+`$ case, the odd part here is a sum of two isomorphic modules. To conclude that $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N`$ is irreducible, we show again that one may go from each irreducible $`𝔙_+sl_2`$-component to the irreducible component containing the $`𝔤`$-highest weight vectors. But this follows from the following computation.
(6.1) $`G_{\frac{1}{2}}^{++}(\alpha a_4+\beta a_5)=\beta \mathrm{\Lambda }(2\mathrm{\Lambda }+2)a_1,\alpha ,\beta ,`$
(6.2) $`G_{\frac{1}{2}}^{}a_4=(2\mathrm{\Lambda }+2)F_0a_1,E_1a_{11}=\mathrm{\Lambda }(\mathrm{\Lambda }1)(2\mathrm{\Lambda }+2)a_1.`$
Now if $`\mathrm{\Lambda }=1`$ the vector $`a_{11}`$ is zero. Therefore the quotient $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N=[]V_1[]V_4[]V_5`$. But then (6.1) and the first identity in (6.2) show that $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })/N`$ is irreducible. The rank of $`L_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{4},\mathrm{\Lambda })`$ is then $`(\mathrm{\Lambda }+1)+2\mathrm{\Lambda }`$, which equals to $`4\mathrm{\Lambda }`$, in the case $`\mathrm{\Lambda }=1`$.
Finally when $`\mathrm{\Lambda }=0`$, the vectors $`a_4=a_5=a_{10}=a_{11}=a_{14}=a_{15}=0`$ so that $`S^\mathrm{\Lambda }`$ reduces to $`\{a_2,a_3,a_6,a_7,a_8,a_9,a_1,a_{12},a_{13},a_{16}\}`$. Hence $`M_{𝔑_+^4}(0,0)/N=a_1`$ is the trivial module and so has rank $`0`$. ∎
###### Proposition 6.3.
Suppose that $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$. Then $`L_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ is a free $`[]`$-module of rank $`4\mathrm{\Lambda }+8`$.
###### Proof.
By Proposition 6.1 $`a_4`$ and $`a_5`$ are singular vectors of $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })`$ in the case $`\mathrm{\Lambda }1`$.
Assume first that $`\mathrm{\Lambda }3`$. Let $`N_4`$ and $`N_5`$ be the $`𝔤`$-submodules generated by $`a_4`$ and $`a_5`$, respectively. We form the $`𝔤`$-submodule $`N=N_4+N_5`$ and consider $`N^{E_0}`$. The set $`\{u_i^{\mathrm{\Lambda }1}a_4|1i16\}`$ is a set of $`[]`$-generators for $`N_4^{E_0}`$, since $`a_4`$ has $`H_0`$-weight $`\mathrm{\Lambda }1`$. We have
$`u_1^{\mathrm{\Lambda }1}a_4=a_4,u_2^{\mathrm{\Lambda }1}a_4=\mathrm{\Lambda }a_6,u_3^{\mathrm{\Lambda }1}a_4=2\mathrm{\Lambda }a_1\mathrm{\Lambda }a_9+a_{10},u_4^{\mathrm{\Lambda }1}a_4=0,`$
$`u_5^{\mathrm{\Lambda }1}a_4=a_{11},u_6^{\mathrm{\Lambda }1}a_4=0,u_7^{\mathrm{\Lambda }1}a_4=a_{15}+2a_5,u_8^{\mathrm{\Lambda }1}a_4=2\mathrm{\Lambda }a_2+\mathrm{\Lambda }a_{12},`$
$`u_9^{\mathrm{\Lambda }1}a_4=a_{14}+2a_4,u_{10}^{\mathrm{\Lambda }1}a_4=(\mathrm{\Lambda }1)a_{14},u_{11}^{\mathrm{\Lambda }1}a_4=0,u_{12}^{\mathrm{\Lambda }1}a_4=2\mathrm{\Lambda }a_6,`$
$`u_{13}^{\mathrm{\Lambda }1}a_4=\mathrm{\Lambda }a_{16}+2a_{10},u_{14}^{\mathrm{\Lambda }1}a_4=0,u_{15}^{\mathrm{\Lambda }1}a_4=2a_{11},u_{16}^{\mathrm{\Lambda }1}a_4=2a_{14}.`$
Similarly, the following is a set of $`[]`$-generators for $`N_5^{E_0}`$.
$`u_1^{\mathrm{\Lambda }1}a_5=a_5,u_2^{\mathrm{\Lambda }1}a_5=a_{10},u_3^{\mathrm{\Lambda }1}a_5=a_7,u_4^{\mathrm{\Lambda }1}a_5=a_{11},u_5^{\mathrm{\Lambda }1}a_5=0,`$
$`u_6^{\mathrm{\Lambda }1}a_5=a_{14},u_7^{\mathrm{\Lambda }1}a_5=0,u_8^{\mathrm{\Lambda }1}a_5=\mathrm{\Lambda }a_{13},u_9^{\mathrm{\Lambda }1}a_5=a_{15},u_{10}^{\mathrm{\Lambda }1}a_5=0,`$
$`u_{11}^{\mathrm{\Lambda }1}a_5=0,u_{12}^{\mathrm{\Lambda }1}a_5=a_{16},u_{13}^{\mathrm{\Lambda }1}a_5=0,u_{14}^{\mathrm{\Lambda }1}a_5=0,u_{15}^{\mathrm{\Lambda }1}a_5=0,`$
$`u_{16}^{\mathrm{\Lambda }1}a_5=0.`$
Therefore
$$S^\mathrm{\Lambda }=\{a_4,a_5,a_6,a_7,a_92a_1,a_{10},a_{11},a_{12}+2a_2,a_{13},a_{14},a_{15},a_{16}\}$$
is a set of $`[]`$-generators for $`N^{E_0}`$, which implies that $`\{a_1,a_2,a_3,a_8\}`$ is a $`[]`$-basis for $`(M_{𝔑_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })/N)^{E_0}`$ in the case when $`\mathrm{\Lambda }3`$.
In the case when $`\mathrm{\Lambda }=2`$ the set $`\{u_i^{\mathrm{\Lambda }1}a_4,u_i^{\mathrm{\Lambda }1}a_5|1i16,i11\}`$ generate $`N^{E_0}`$. But $`u_{11}^{\mathrm{\Lambda }1}a_4=u_{11}^{\mathrm{\Lambda }1}a_5=0`$, and hence $`\{a_1,a_2,a_3,a_8\}`$ is also a $`[]`$-basis for $`(M_{𝔑_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })/N)^{E_0}`$ in this case as well.
In the case when $`\mathrm{\Lambda }=1`$ we note that $`a_{11}=0`$ and $`\{u_i^{\mathrm{\Lambda }1}a_4,u_i^{\mathrm{\Lambda }1}a_5|1i16,i4,5,10,11,14,15\}`$ generates $`N^{E_0}`$ over $`[]`$. From the formulas above one sees that a set of $`[]`$-generators for $`N^{E_0}`$ is given by the set $`S^\mathrm{\Lambda }`$ above, but with $`a_{11}`$ removed. Hence the quotient module is again generated freely over $`[]`$ by $`\{a_1,a_2,a_3,a_8\}`$.
Hence in the case when $`\mathrm{\Lambda }1`$ the quotient module $`(M_{𝔑_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })/N)^{E_0}`$ is generated freely over $`[]`$ by $`\{a_1,a_2,a_3,a_8\}`$. Now $`a_1`$ has $`H_0`$-weight $`\mathrm{\Lambda }`$, $`a_2`$ and $`a_3`$ both have $`H_0`$-weight $`\mathrm{\Lambda }+1`$, and $`a_8`$ has $`H_0`$-weight $`\mathrm{\Lambda }+2`$. Therefore $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })/N`$ has rank $`(\mathrm{\Lambda }+1)+2(\mathrm{\Lambda }+2)+(\mathrm{\Lambda }+3)=4\mathrm{\Lambda }+8`$ over $`[]`$. So it remains to show that $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })/N`$ is irreducible.
We again study $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })/N`$ as a $`𝔙_+sl_2`$-module. It is easy to check that $`L_n`$, $`n1`$, annihilates $`a_1,a_2,a_3,a_8`$ and hence $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })/N`$ is a direct sum of the following four irreducible $`𝔙_+sl_2`$-modules: $`[]V_1L_{𝔙_+}(\frac{\mathrm{\Lambda }+2}{2})U^\mathrm{\Lambda }`$, $`[]V_2L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }+1}`$, $`[]V_3L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }+1}`$ and $`[]V_8L_{𝔙_+}(\frac{\mathrm{\Lambda }}{2})U^{\mathrm{\Lambda }+2}`$, where $`V_i`$ is the $`sl_2`$-submodule generated by the vector $`a_i`$. Again $`[]V_2[]V_3`$ as $`𝔙_+sl_2`$-modules. Now we compute
$`G_{\frac{1}{2}}^{}(\alpha a_2+\beta a_3)=2\alpha (\mathrm{\Lambda }+1)a_1,\alpha ,\beta ,`$
$`G_{\frac{1}{2}}^+a_3=2(\mathrm{\Lambda }+1)a_1,F_1a_8=2(\mathrm{\Lambda }+1)a_1.`$
Therefore $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })/N`$ is irreducible.
Now consider the case of $`\mathrm{\Lambda }=0`$. By Proposition 6.1 $`a_6`$, $`a_7`$ and $`a_92a_1`$ are singular vectors inside $`M_{𝔑_+^4}(1,0)`$. Let $`N_6`$, $`N_7`$ and $`N_9`$ be the $`𝔤`$-submodules generated by $`a_6`$, $`a_7`$ and $`a_92a_1`$, respectively, and put $`N=N_6+N_7+N_9`$. We note that $`a_6`$, $`a_7`$ and $`a_92a_1`$ have $`H_0`$-weight $`0`$, hence $`N_6^{E_0}`$ is generated over $`[]`$ by $`\{u_i^\mathrm{\Lambda }a_6|1i16,i4,5,10,11,14,15\}`$ and similarly for $`N_7^{E_0}`$ and $`N_9^{E_0}`$. We first compute a set of $`[]`$-generators for $`N_6^{E_0}`$.
$`u_1^\mathrm{\Lambda }a_6=a_6,u_2^\mathrm{\Lambda }a_6=0,u_3^\mathrm{\Lambda }a_6=2a_2+a_{12},u_6^\mathrm{\Lambda }a_6=0,`$
$`u_7^\mathrm{\Lambda }a_6=4^2a_12a_9+a_{16},u_8^\mathrm{\Lambda }a_6=0,u_9^\mathrm{\Lambda }a_6=4a_6,`$
$`u_{12}^\mathrm{\Lambda }a_6=0,u_{13}^\mathrm{\Lambda }a_6=4^2a_2+2a_{12},u_{16}^\mathrm{\Lambda }a_6=4^2a_6.`$
A set of $`[]`$-generators for $`N_7^{E_0}`$ is given as follows.
$`u_1^\mathrm{\Lambda }a_7=a_7,u_2^\mathrm{\Lambda }a_7=a_{13},u_3^\mathrm{\Lambda }a_7=0,u_6^\mathrm{\Lambda }a_7=a_{16},u_7^\mathrm{\Lambda }a_7=0,`$
$`u_8^\mathrm{\Lambda }a_7=0,u_9^\mathrm{\Lambda }a_7=0,u_{12}^\mathrm{\Lambda }a_7=0,u_{13}^\mathrm{\Lambda }a_7=0,u_{16}^\mathrm{\Lambda }a_7=0.`$
Finally we have the following set of $`[]`$-generators for $`N_9^{E_0}`$.
$`u_1^\mathrm{\Lambda }(a_92a_1)=a_92a_1,u_2^\mathrm{\Lambda }(a_92a_1)=a_{12}2a_2,`$
$`u_3^\mathrm{\Lambda }(a_92a_1)=a_{13},u_6^\mathrm{\Lambda }(a_92a_1)=2a_6,u_7^\mathrm{\Lambda }(a_92a_1)=2a_7,`$
$`u_8^\mathrm{\Lambda }(a_92a_1)=0,u_9^\mathrm{\Lambda }(a_92a_1)=2a_{16},u_{12}^\mathrm{\Lambda }(a_92a_1)=0,`$
$`u_{13}^\mathrm{\Lambda }(a_92a_1)=2a_{13},u_{16}^\mathrm{\Lambda }(a_92a_1)=2a_{16}.`$
From this it follows that $`\{a_6,a_7,a_92a_1,a_{12}+2a_2,a_{13},a_{16}\}`$ generate $`N^{E_0}`$ over $`[]`$. But $`a_4=a_5=a_{10}=a_{11}=a_{14}=a_{15}=0`$, and thus $`(M_{𝔑_+^4}(1,0)/N)^{E_0}`$ is generated over $`[]`$ by the vectors $`a_1`$, $`a_2`$, $`a_3`$ and $`a_8`$, which takes us back to the case when $`\mathrm{\Lambda }1`$, except that here $`[]V_8`$ is not irreducible. It contains a unique irreducible submodule isomorphic to $`L_{𝔙_+}(1)U^2`$ generated by $`a_8`$. But then the above calculation plus the fact that
$$F_2a_8=4(\mathrm{\Lambda }+1)a_1$$
show that $`M_{𝔑_+^4}(1,0)/N`$ is irreducible of rank $`8`$. ∎
###### Theorem 6.1.
The modules $`L_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$, for $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda }_+`$, form a complete list of non-isomorphic finite (over $`[]`$) irreducible $`SK(1,4)_+`$-modules. Furthermore $`L_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ as a $`[]`$-module has rank
* $`4\mathrm{\Lambda }`$, in the case $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$,
* $`4\mathrm{\Lambda }+8`$, in the case $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$.
* $`16\mathrm{\Lambda }+16`$, in all other cases.
Furthermore the $`[]`$-rank of $`L_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })_{\overline{0}}`$ equals the $`[]`$-rank of $`L_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })_{\overline{1}}`$ in all cases.
###### Remark 6.2.
Translating the above theorem into the languages of modules over conformal algebras and of conformal modules is again straightforward. We therefore obtain that all finite irreducible modules over the “small” $`N=4`$ conformal superalgebra are of the form $`L_{𝔑^4}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda })`$, where $`\alpha ,\mathrm{\Delta }`$ and $`\mathrm{\Lambda }_+`$. The definition of these modules and also the action of the conformal superalgebra on them are easily gotten from our explicit description of a $`[]`$-basis in this section and hence omitted, as to reproduce them would take up quite a significant portion of space. Again we only note that the adjoint module is isomorphic to $`L_{𝔑^4}(0,1,2)`$.
## 7. Finite irreducible Modules over the “big” $`N=4`$ conformal superalgebra
In this section we give a classification of finite irreducible conformal modules over the contact superalgebra $`K(1,4)`$, also known as the “big” $`N=4`$ superconformal algebra. Our approach is based on our results obtained in Section 6.
Recall from Section 6 that $`L_n^\beta `$, $`X_n^\beta `$ and $`x_r^\beta `$, where $`X=H,E,F`$, $`x=G^{++}`$, $`G^+`$,$`G^+`$,$`G^{}`$, $`n`$, $`r\frac{1}{2}+`$ and the fixed number $`\beta `$ is either $`1`$ or $`1`$, provide a basis for a copy of $`SK(1,4)`$ inside $`K(1,4)`$. In this section it will be convenient to distinguish these two copies. We therefore denote the copy obtained by setting $`\beta =1`$ simply by $`SK(1,4)`$, while the copy obtained by setting $`\beta =1`$ by $`\overline{SK}(1,4)`$. It is easy to see from our formulas that $`K(1,4)=SK(1,4)+\overline{SK}(1,4)`$. Similarly we distinguish the basis elements of $`SK(1,4)`$ and $`\overline{SK}(1,4)`$ as follows. The generators inside $`SK(1,4)`$ will be denoted by $`L_n,X_n,x_r`$, while generators inside $`\overline{SK}(1,4)`$ will be denoted by $`\overline{L}_n,\overline{X}_n,\overline{x}_r`$, where again $`X=H,E,F`$, $`x=G^{++}`$,$`G^+`$,$`G^+`$,$`G^{}`$. Of course we have $`x_{\frac{1}{2}}=\overline{x}_{\frac{1}{2}}`$, $`L_1=\overline{L}_1`$ and $`L_0=\overline{L}_0`$.
###### Remark 7.1.
The map $`\varphi :SK(1,4)\overline{SK}(1,4)`$ defined by $`\varphi (L_n)=\overline{L}_n`$, $`\varphi (X_n)=\overline{X}_n`$, $`\varphi (G_r^{++})=\overline{G}_r^{++}`$, $`\varphi (G_r^+)=\overline{G}_r^+`$, $`\varphi (G_r^+)=\overline{G}_r^+`$ and $`\varphi (G_r^{})=\overline{G}_r^{}`$, where $`n`$ and $`r\frac{1}{2}+`$ is an isomorphism of Lie superalgebras. Thus all formulas in Section 6 with $`\varphi (L_n)`$, $`\varphi (X_n)`$ and $`\varphi (x_r)`$ replacing $`L_n`$, $`X_n`$ and $`x_r`$, respectively, remain valid.
Let $`𝔤=K(1,4)_+`$ be the annihilation subalgebra of $`K(1,4)`$ so that we have $`𝔤=SK(1,4)_++\overline{SK}(1,4)_+`$, the sum of the corresponding annihilation subalgebras. We have as before $`𝔤=_{j1}𝔤_j`$, where $`j\frac{1}{2}+`$. Furthermore $`𝔤_{}=SK(1,4)_{}=\overline{SK}(1,4)_{}`$ and $`𝔤_0=L_0sl_2\overline{sl}_2cso_4`$, where $`sl_2`$ and $`\overline{sl}_2`$ denote two copies of the Lie algebra $`sl_2`$, generated by $`H_0,E_0,F_0`$ and $`\overline{H}_0,\overline{E}_0,\overline{F}_0`$, respectively.
Let $`U^{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$ be the finite-dimensional irreducible $`sl_2\overline{sl}_2`$-module of highest weight $`(\mathrm{\Lambda },\overline{\mathrm{\Lambda }})_+\times _+`$ on which $`L_0`$ acts as the scalar $`\mathrm{\Delta }`$ and let $`v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$ denote a highest weight vector in $`U^{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$ so that $`H_0v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}=\mathrm{\Lambda }v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$, $`\overline{H}_0v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}=\overline{\mathrm{\Lambda }}v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$ and $`L_0v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}=\mathrm{\Delta }v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$. Regarding $`U^{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$ as a module over $`_0=_{j0}𝔤_j`$ it follows from Theorem 3.1 that every finite irreducible $`𝔤`$-module is a quotient of $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})=\mathrm{Ind}__0^𝔤U^{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$. The unique irreducible quotient will be denoted by $`L_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$.
Now $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$ is a completely reducible $`𝔤_0`$-module, and the subspace of $`E_0\overline{E}_0`$-invariants, denoted by $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})^{E_0,\overline{E}_0}`$, is a free $`[]`$-submodule of $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$. We write down explicit formulas for a $`[]`$-basis for the space $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})^{E_0,\overline{E}_0}`$, which in the case when $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}2`$ is as follows:
$`b_1=v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},b_2=G_{\frac{1}{2}}^{++}v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_3=(\mathrm{\Lambda }G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{++}F_0)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},b_4=(\overline{\mathrm{\Lambda }}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{++}\overline{F}_0)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_5=(\mathrm{\Lambda }\overline{\mathrm{\Lambda }}G_{\frac{1}{2}}^{}\overline{\mathrm{\Lambda }}G_{\frac{1}{2}}^+F_0\mathrm{\Lambda }G_{\frac{1}{2}}^+\overline{F}_0+G_{\frac{1}{2}}^{++}F_0\overline{F}_0)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_6=G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},b_7=\left(\mathrm{\Lambda }(G_{\frac{1}{2}}^+G_{\frac{1}{2}}^++G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^{})2G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+F_0\right)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_8=((\mathrm{\Lambda }1)(\mathrm{\Lambda }G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}+G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+F_0+G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^{}F_0)G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+F_0^2)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_9=G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},b_{10}=\left(\overline{\mathrm{\Lambda }}(G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^{}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+)2G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+\overline{F_0}\right)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_{11}=\left((\overline{\mathrm{\Lambda }}1)(\overline{\mathrm{\Lambda }}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}+G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+\overline{F}_0+G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^{}\overline{F}_0)G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+\overline{F}_0^2\right)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_{12}=G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},b_{13}=(\mathrm{\Lambda }G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_{14}=(\overline{\mathrm{\Lambda }}G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_{15}=\left(\mathrm{\Lambda }G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}(\overline{\mathrm{\Lambda }}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{++}\overline{F}_0)+\overline{\mathrm{\Lambda }}G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+(G_{\frac{1}{2}}^{}F_0G^+F_0\overline{F}_0)\right)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}},`$
$`b_{16}=\left(G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+G_{\frac{1}{2}}^+G_{\frac{1}{2}}^{}(G_{\frac{1}{2}}^+G_{\frac{1}{2}}^++G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^{})\right)v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}.`$
In the case when $`\mathrm{\Lambda }=\overline{\mathrm{\Lambda }}=1`$ (respectively $`\mathrm{\Lambda }=\overline{\mathrm{\Lambda }}=0`$) we have $`b_8=b_{11}=0`$ (respectively $`b_3=b_4=b_5=b_7=b_8=b_{10}=b_{11}=b_{13}=b_{14}=b_{15}=0`$), thus giving us $`14`$ (respectively $`6`$) generators. Other cases are easily described as well, however, we will not need them because of Proposition 7.1 below. Thus we will omit them.
We will, as before, denote the coefficient of $`v_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$ in $`b_i`$ by $`u_i^{\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$ for $`1i16`$. In the case when $`\mathrm{\Lambda }=\overline{\mathrm{\Lambda }}`$, which is the only case we will be concerned with in what follows, we simply write $`u_i^\mathrm{\Lambda }`$ for $`u_i^{\mathrm{\Lambda },\mathrm{\Lambda }}`$ and also $`v_{\mathrm{\Delta },\mathrm{\Lambda }}`$ for $`v_{\mathrm{\Delta },\mathrm{\Lambda },\mathrm{\Lambda }}`$.
###### Proposition 7.1.
If $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$ is a reducible $`𝔤`$-module, then either $`2\mathrm{\Delta }\mathrm{\Lambda }=2\mathrm{\Delta }\overline{\mathrm{\Lambda }}=0`$ or else $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=2\mathrm{\Delta }+\overline{\mathrm{\Lambda }}+2=0`$. In particular if $`\mathrm{\Lambda }\overline{\mathrm{\Lambda }}`$, then $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$ is irreducible.
###### Proof.
As a module over $`SK(1,4)_+`$ we have $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})=U(𝔤_{})U^{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$ is a direct sum of $`\overline{\mathrm{\Lambda }}+1`$ copies of $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$, generated by the highest weight vectors $`\overline{F}_0^jv_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$, where $`0j\overline{\mathrm{\Lambda }}`$. Since the $`\overline{H}_0`$-weights of the $`\overline{F}_0^jv_{\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }}}`$’s are all distinct for distinct $`j`$’s it follows that these modules as $`SK(1,4)_+\overline{H}_0`$-modules are all non-isomorphic. Therefore if $`M_{𝔑_+^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ is irreducible over $`SK(1,4)_+`$, then $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$ is irreducible over $`𝔤`$. From this and Corollary 6.1 we thus conclude that in the case when $`\mathrm{\Delta }2\mathrm{\Lambda }0`$ and $`\mathrm{\Delta }+2\mathrm{\Lambda }+20`$ the $`𝔤`$-module $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$ is irreducible.
By symmetry we conclude that if $`\mathrm{\Delta }2\overline{\mathrm{\Lambda }}0`$ and $`\mathrm{\Delta }+2\overline{\mathrm{\Lambda }}+20`$, then $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$ is irreducible over $`𝔤`$ as well.
Therefore $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$ is possibly reducible only if both $`\mathrm{\Lambda }`$ and $`\overline{\mathrm{\Lambda }}`$ satisfy one of the two linear equations $`\mathrm{\Delta }2x=0`$ and $`\mathrm{\Delta }+2x+2=0`$. But the case $`\mathrm{\Delta }2\mathrm{\Lambda }=0`$ and $`\mathrm{\Delta }+2\overline{\mathrm{\Lambda }}+2=0`$ is not possible, since both $`\mathrm{\Lambda }`$ and $`\overline{\mathrm{\Lambda }}`$ are non-negative integers. By the same token $`\mathrm{\Delta }2\overline{\mathrm{\Lambda }}=0`$ and $`\mathrm{\Delta }+2\mathrm{\Lambda }+2=0`$ is not possible, either. Hence either we have $`\mathrm{\Delta }2\mathrm{\Lambda }=0`$ and $`\mathrm{\Delta }2\overline{\mathrm{\Lambda }}=0`$ or else $`\mathrm{\Delta }+2\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Delta }+2\overline{\mathrm{\Lambda }}+2=0`$. In either case we must have $`\mathrm{\Lambda }=\overline{\mathrm{\Lambda }}`$. ∎
The next step is to analyze proper singular vectors inside $`M_{𝔖^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$. (The definitions of singular vectors and proper singular vectors of $`𝔤`$ are of course analogous.) By Proposition 7.1 proper singular vectors exist only if $`\mathrm{\Lambda }=\overline{\mathrm{\Lambda }}`$ with either $`2\mathrm{\Delta }+\mathrm{\Lambda }=0`$ or $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$.
###### Proposition 7.2.
A complete list of proper singular vectors inside $`M_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\mathrm{\Lambda })`$ is given by:
* $`\alpha b_2`$, $`\alpha 0`$, in the case $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$.
* $`\alpha b_5`$, $`\alpha 0`$, in the case $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }1`$.
###### Proof.
Since as a $`SK(1,4)_+`$-module $`M_{𝔖^4}(\mathrm{\Delta },\mathrm{\Lambda },\mathrm{\Lambda })`$ is a direct sum of $`\mathrm{\Lambda }+1`$ copies of $`M_{𝔑^4}(\mathrm{\Delta },\mathrm{\Lambda })`$ we obtain a description of the vector space spanned by all proper $`SK(1,4)_+`$-singular vectors by virtue of Proposition 6.1. But as a $`\overline{SK}(1,4)_+`$-module $`M_{𝔖^4}(\mathrm{\Delta },\mathrm{\Lambda },\mathrm{\Lambda })`$ is also a direct sum of $`\mathrm{\Lambda }+1`$ copies of $`M_{𝔑^4}(\mathrm{\Delta },\mathrm{\Lambda })`$, from which we obtain similarly a description of the vector space spanned by all proper $`\overline{SK}(1,4)_+`$-singular vectors (see Remark 7.1). The intersection of these two spaces is the space of proper singular vectors.
In the case when $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$ it follows from Proposition 6.1 that the space of proper $`SK(1,4)_+`$-singular vectors is spanned by $`G_{\frac{1}{2}}^{++}\overline{F}_0^jv_{\mathrm{\Delta },\mathrm{\Lambda }}`$, $`G_{\frac{1}{2}}^+\overline{F}_0^jv_{\mathrm{\Delta },\mathrm{\Lambda }}`$ and $`G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+\overline{F}_0^jv_{\mathrm{\Delta },\mathrm{\Lambda }}`$, for $`0j\mathrm{\Lambda }`$. On the other hand the space of proper $`\overline{SK}(1,4)_+`$-singular vectors is spanned by $`G_{\frac{1}{2}}^{++}F_0^jv_{\mathrm{\Delta },\mathrm{\Lambda }}`$, $`G_{\frac{1}{2}}^+F_0^jv_{\mathrm{\Delta },\mathrm{\Lambda }}`$ and $`G_{\frac{1}{2}}^{++}G_{\frac{1}{2}}^+F_0^jv_{\mathrm{\Delta },\mathrm{\Lambda }}`$, for $`0j\mathrm{\Lambda }`$. It is not hard to see that the intersection of these two spaces is the one-dimensional space spanned by $`G_{\frac{1}{2}}^{++}v_{\mathrm{\Delta },\mathrm{\Lambda }}`$, which is $`b_2`$.
Other cases are analogous and so we omit the details. ∎
###### Proposition 7.3.
Suppose that $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$. Then $`L_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\mathrm{\Lambda })`$ is a free $`[]`$-module of rank $`8\mathrm{\Lambda }(\mathrm{\Lambda }+1)`$.
###### Proof.
By Proposition 7.2 $`b_2`$ is a singular vector in $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })`$. Consider $`N`$, the $`𝔤`$-submodule generated by $`b_2`$. Then we have $`N=U(𝔤_{})V`$, where $`V`$ is the irreducible $`sl_2\overline{sl}_2`$-submodule generated by $`b_2`$. Let us compute the space $`N^{E_0,\overline{E}_0}`$, the space of $`(E_0\overline{E}_0)`$-invariants inside $`N`$. Since the $`(H_0,\overline{H}_0)`$-weight of $`b_2`$ is $`(\mathrm{\Lambda }+1,\mathrm{\Lambda }+1)`$, we know that $`N^{E_0,\overline{E}_0}`$ is a free $`[]`$-module generated over $`[]`$ by $`\{u_i^{\mathrm{\Lambda }+1}b_2|1i16\}`$. We have
$`u_1^{\mathrm{\Lambda }+1}b_2=b_2,u_2^{\mathrm{\Lambda }+1}b_2=0,u_3^{\mathrm{\Lambda }+1}b_2=(\mathrm{\Lambda }+2)b_9,u_4^{\mathrm{\Lambda }+1}b_2=(\mathrm{\Lambda }+2)b_6,`$
$`u_5^{\mathrm{\Lambda }+1}b_2={\displaystyle \frac{\mathrm{\Lambda }+2}{2}}(b_7+b_{10}+4(\mathrm{\Lambda }+1)b_1),u_6^{\mathrm{\Lambda }+1}b_2=0,u_7^{\mathrm{\Lambda }+1}b_2=(\mathrm{\Lambda }+3)b_{12},`$
$`u_8^{\mathrm{\Lambda }+1}b_2=(\mathrm{\Lambda }+2)(b_{13}2b_3),u_9^{\mathrm{\Lambda }+1}b_2=0,u_{10}^{\mathrm{\Lambda }+1}b_2=(\mathrm{\Lambda }+3)b_{12}4(\mathrm{\Lambda }+1)b_2,`$
$`u_{11}^{\mathrm{\Lambda }+1}b_2=(\mathrm{\Lambda }+2)(b_{14}2b_4),u_{12}^{\mathrm{\Lambda }+1}b_2=0,u_{13}^{\mathrm{\Lambda }+1}b_2=2(\mathrm{\Lambda }+2)b_9,`$
$`u_{14}^{\mathrm{\Lambda }+1}b_2=2(\mathrm{\Lambda }+2)b_6,u_{15}^{\mathrm{\Lambda }+1}b_2=4(\mathrm{\Lambda }+1)^2b_1(\mathrm{\Lambda }+2)^2b_{16}+(\mathrm{\Lambda }+2)b_7,`$
$`u_{16}^{\mathrm{\Lambda }+1}b_2=4b_{12}.`$
It follows that $`N^{E_0,\overline{E}_0}`$ is generated over $`[]`$ by the set
$`S^\mathrm{\Lambda }=\{`$ $`b_2,b_6,b_7+b_{10}+4(\mathrm{\Lambda }+2)b_1,b_9,b_{12},b_{13}2b_3,b_{14}2b_4,`$
$`b_{16}({\displaystyle \frac{1}{\mathrm{\Lambda }+2}})b_72{\displaystyle \frac{(\mathrm{\Lambda }+1)}{(\mathrm{\Lambda }+2)^2}}^2b_1\}.`$
In the case when $`\mathrm{\Lambda }2`$ it follows from the description of $`S^\mathrm{\Lambda }`$ that $`\{b_1,b_3,b_4,b_5,b_8,b_{10}+2\mathrm{\Lambda }b_1,b_{11},b_{15}\}`$ is a $`[]`$-basis for the $`(E_0\overline{E}_0)`$-invariants of the quotient space $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$. (The choice of $`b_{10}+2\mathrm{\Lambda }b_1`$ instead of just $`b_{10}`$ will be explained later.)
The $`(L_0,H_0,\overline{H}_0)`$-weights of $`b_1`$, $`b_3`$, $`b_4`$, $`b_5`$, $`b_8`$, $`b_{10}+2\mathrm{\Lambda }b_1`$, $`b_{11}`$, $`b_{15}`$ are $`(\mathrm{\Delta },\mathrm{\Lambda },\mathrm{\Lambda })`$, $`(\mathrm{\Delta }+\frac{1}{2},\mathrm{\Lambda }1,\mathrm{\Lambda }+1)`$, $`(\mathrm{\Delta }+\frac{1}{2},\mathrm{\Lambda }+1,\mathrm{\Lambda }1)`$, $`(\mathrm{\Delta }+\frac{1}{2},\mathrm{\Lambda }1,\mathrm{\Lambda }1)`$, $`(\mathrm{\Delta }+1,\mathrm{\Lambda }2,\mathrm{\Lambda })`$, $`(\mathrm{\Delta }+1,\mathrm{\Lambda },\mathrm{\Lambda })`$, $`(\mathrm{\Delta }+1,\mathrm{\Lambda },\mathrm{\Lambda }2)`$, $`(\mathrm{\Delta }+\frac{3}{2},\mathrm{\Lambda }1,\mathrm{\Lambda }1)`$, respectively. Hence $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ is a free $`[]`$-module of rank $`8\mathrm{\Lambda }(\mathrm{\Lambda }+1)`$. So we need to show that $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ is irreducible.
Now $`L_n`$, $`n1`$, together with $`E_0,H_0,F_0`$ and $`\overline{E}_0,\overline{H}_0,\overline{F}_0`$ generate a copy of $`(𝔙_+sl_2\overline{sl}_2)`$, which thus allow us to study the $`(𝔙_+sl_2\overline{sl}_2)`$-module structure of $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$. We can easily check that $`L_n`$, for $`n1`$, annihilates the vectors $`b_1`$, $`b_3`$, $`b_4`$, $`b_5`$, $`b_8`$, $`b_{10}+2\mathrm{\Lambda }b_1`$, $`b_{11}`$, $`b_{15}`$. (We want to point out that $`b_{10}`$ is not annihilated by $`L_n`$, for $`n1`$, hence the choice of $`b_{10}+2\mathrm{\Lambda }b_1`$.) Thus $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ as a $`(𝔙_+sl_2\overline{sl}_2)`$-module is a direct sum of the following eight irreducible modules: $`[]V_1L_{𝔙_+}(\frac{\mathrm{\Lambda }}{2})U^{\mathrm{\Lambda },\mathrm{\Lambda }}`$, $`[]V_3L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }1,\mathrm{\Lambda }+1}`$, $`[]V_4L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }+1,\mathrm{\Lambda }1}`$, $`[]V_5L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }1,\mathrm{\Lambda }1}`$, $`[]V_8L_{𝔙_+}(\frac{\mathrm{\Lambda }+2}{2})U^{\mathrm{\Lambda }2,\mathrm{\Lambda }}`$, $`[]V_{10}L_{𝔙_+}(\frac{\mathrm{\Lambda }+2}{2})U^{\mathrm{\Lambda },\mathrm{\Lambda }}`$, $`[]V_{11}L_{𝔙_+}(\frac{\mathrm{\Lambda }+2}{2})U^{\mathrm{\Lambda },\mathrm{\Lambda }2}`$, $`[]V_{15}L_{𝔙_+}(\frac{\mathrm{\Lambda }+3}{2})U^{\mathrm{\Lambda }1,\mathrm{\Lambda }1}`$, where $`V_i`$ is the irreducible $`sl_2\overline{sl}_2`$-module generated by $`b_i`$, for $`i10`$, and $`V_{10}`$ is generated by $`b_{10}+2\mathrm{\Lambda }b_1`$, and finally $`U^{\mu ,\mu ^{}}`$ denotes the irreducible $`sl_2\overline{sl}_2`$-module of highest weight $`(\mu ,\mu ^{})`$. Note that as $`(𝔙_+sl_2\overline{sl}_2)`$-modules they are all non-isomorphic and thus to show that $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ is irreducible, it suffices to show that one may send a $`(𝔙_+sl_2\overline{sl}_2)`$-highest weight vector in any irreducible $`(𝔙_+sl_2\overline{sl}_2)`$-component to the irreducible component containing the $`𝔤`$-highest weight vectors. This follows from the following computation.
$`G_{\frac{1}{2}}^{}b_3=2(\mathrm{\Lambda }+1)F_0b_1,\overline{G}_{\frac{1}{2}}^{}b_4=2(\mathrm{\Lambda }+1)\overline{F}_0b_1,`$
$`G_{\frac{1}{2}}^{++}b_5=2\mathrm{\Lambda }^2(\mathrm{\Lambda }+1)b_1,E_1b_8=2\mathrm{\Lambda }(\mathrm{\Lambda }1)(\mathrm{\Lambda }+1)b_1,`$
$`\overline{F}_1(b_{10}+2\mathrm{\Lambda }b_1)=2(\mathrm{\Lambda }+2)\overline{F}_0b_1,\overline{E}_1b_{11}=2\mathrm{\Lambda }(\mathrm{\Lambda }1)(\mathrm{\Lambda }+1)b_1,`$
$`\overline{G}_{\frac{3}{2}}^{++}b_{15}=2\mathrm{\Lambda }^2(\mathrm{\Lambda }+1)b_1.`$
Now if $`\mathrm{\Lambda }=1`$ the vectors $`b_8=b_{11}=0`$. Therefore $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ is $`[]V_1[]V_3[]V_4[]V_5[]V_{10}[]V_{15}`$. But then the above calculation also shows that $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ is irreducible. The rank of $`L_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })`$ is then $`4+3+3+1+4+1=16`$, which is equal to $`8\mathrm{\Lambda }(\mathrm{\Lambda }+1)`$ in the case when $`\mathrm{\Lambda }=1`$.
Finally when $`\mathrm{\Lambda }=0`$, the vectors $`b_3=b_4=b_5=b_7=b_8=b_{10}=b_{11}=b_{13}=b_{14}=b_{15}=0`$ and $`S^\mathrm{\Lambda }`$ reduces to $`\{b_2,b_6,b_1,b_9,b_{12},b_{16}\}`$. Hence $`M_{𝔖_+^4}(0,0,0)/N=b_1`$ is the trivial module and so has rank $`0`$. ∎
###### Proposition 7.4.
Suppose that $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }1`$. Then $`L_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\mathrm{\Lambda })`$ is a free $`[]`$-module of rank $`8(\mathrm{\Lambda }+1)(\mathrm{\Lambda }+2)`$.
###### Proof.
By Proposition 7.2 $`b_5`$ is a singular vector in $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda },\mathrm{\Lambda })`$. Let $`N`$ be the $`𝔤`$-submodule generated by $`b_5`$ so that $`N=U(𝔤_{})V`$, where $`V`$ is the irreducible $`sl_2\overline{sl}_2`$-submodule generated by $`b_5`$. Consider $`N^{E_0,\overline{E}_0}`$, the subspace in $`N`$ of $`E_0\overline{E}_0`$-invariants. Now the $`(H_0,\overline{H}_0)`$-weight of $`b_5`$ is $`(\mathrm{\Lambda }1,\mathrm{\Lambda }1)`$ and so $`N^{E_0,\overline{E}_0}`$ is a free $`[]`$-module generated over $`[]`$ by $`\{u_i^{\mathrm{\Lambda }1}b_5|1i16\}`$. We have
$`u_1^{\mathrm{\Lambda }1}b_5=b_5,u_2^{\mathrm{\Lambda }1}b_5={\displaystyle \frac{1}{2}}(b_7+b_{10}),u_3^{\mathrm{\Lambda }1}b_5=\mathrm{\Lambda }b_8,u_4^{\mathrm{\Lambda }1}b_5=\mathrm{\Lambda }b_{11},`$
$`u_5^{\mathrm{\Lambda }1}b_5=0,u_6^{\mathrm{\Lambda }1}b_5=\mathrm{\Lambda }b_{14},u_7^{\mathrm{\Lambda }1}b_5=(\mathrm{\Lambda }1)b_{15},u_8^{\mathrm{\Lambda }1}b_5=0,`$
$`u_9^{\mathrm{\Lambda }1}b_5=\mathrm{\Lambda }b_{13},u_{10}^{\mathrm{\Lambda }1}b_5=(\mathrm{\Lambda }1)b_{15},u_{11}^{\mathrm{\Lambda }1}b_5=0,`$
$`u_{12}^{\mathrm{\Lambda }1}b_5=\mathrm{\Lambda }(\mathrm{\Lambda }b_{16}+b_7),u_{13}^{\mathrm{\Lambda }1}b_5=0,u_{14}^{\mathrm{\Lambda }1}b_5=0,u_{15}^{\mathrm{\Lambda }1}b_5=0,`$
$`u_{16}^{\mathrm{\Lambda }1}b_5=b_{15}.`$
It follows that in the case $`\mathrm{\Lambda }2`$ that $`N^{E_0,\overline{E}_0}`$ is generated over $`[]`$ by the set
$$S^\mathrm{\Lambda }=\{b_5,b_7+b_{10},b_8,b_{11},b_{13},b_{14},b_{15},\mathrm{\Lambda }b_{16}+b_7\}.$$
Hence in this case $`\{b_1,b_2,b_3,b_4,b_6,b_9,b_{10}+2\mathrm{\Lambda }b_1,b_{12}\}`$ is a $`[]`$-basis for the $`(E_0\overline{E}_0)`$-invariants of $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$.
The $`(L_0,H_0,\overline{H}_0)`$-weights of $`b_1`$, $`b_2`$, $`b_3`$, $`b_4`$, $`b_6`$, $`b_9`$, $`b_{10}+2\mathrm{\Lambda }b_1`$, $`b_{12}`$ are $`(\mathrm{\Delta },\mathrm{\Lambda },\mathrm{\Lambda })`$, $`(\mathrm{\Delta }+\frac{1}{2},\mathrm{\Lambda }+1,\mathrm{\Lambda }+1)`$, $`(\mathrm{\Delta }+\frac{1}{2},\mathrm{\Lambda }1,\mathrm{\Lambda }+1)`$, $`(\mathrm{\Delta }+\frac{1}{2},\mathrm{\Lambda }+1,\mathrm{\Lambda }1)`$, $`(\mathrm{\Delta }+1,\mathrm{\Lambda }+2,\mathrm{\Lambda })`$, $`(\mathrm{\Delta }+1,\mathrm{\Lambda },\mathrm{\Lambda }+2)`$, $`(\mathrm{\Delta }+1,\mathrm{\Lambda },\mathrm{\Lambda })`$, $`(\mathrm{\Delta }+\frac{3}{2},\mathrm{\Lambda }+1,\mathrm{\Lambda }+1)`$, respectively. Hence $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda })/N`$ is a free $`[]`$-module of rank $`8(\mathrm{\Lambda }+1)(\mathrm{\Lambda }+2)`$. So we need to show that $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ is irreducible.
Again we will study the $`(𝔙_+sl_2\overline{sl}_2)`$-module structure of $`M_{𝔑_+^4}(\frac{\mathrm{\Lambda }}{2},\mathrm{\Lambda })/N`$. We can check directly that $`L_n`$, for $`n1`$, annihilates the vectors $`b_1`$, $`b_2`$, $`b_3`$, $`b_4`$, $`b_6`$, $`b_9`$, $`b_{10}+2\mathrm{\Lambda }b_1`$, $`b_{12}`$. Thus $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ as a $`(𝔙_+sl_2\overline{sl}_2)`$-module is a direct sum of the following eight irreducible modules: $`[]V_1L_{𝔙_+}(\frac{\mathrm{\Lambda }+2}{2})U^{\mathrm{\Lambda },\mathrm{\Lambda }}`$, $`[]V_2L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }+1,\mathrm{\Lambda }+1}`$, $`[]V_3L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }1,\mathrm{\Lambda }+1}`$, $`[]V_4L_{𝔙_+}(\frac{\mathrm{\Lambda }+1}{2})U^{\mathrm{\Lambda }+1,\mathrm{\Lambda }1}`$, $`[]V_6L_{𝔙_+}(\frac{\mathrm{\Lambda }}{2})U^{\mathrm{\Lambda }+2,\mathrm{\Lambda }}`$, $`[]V_9L_{𝔙_+}(\frac{\mathrm{\Lambda }}{2})U^{\mathrm{\Lambda },\mathrm{\Lambda }+2}`$, $`[]V_{10}L_{𝔙_+}(\frac{\mathrm{\Lambda }}{2})U^{\mathrm{\Lambda },\mathrm{\Lambda }}`$, $`[]V_{12}L_{𝔙_+}(\frac{\mathrm{\Lambda }1}{2})U^{\mathrm{\Lambda }+1,\mathrm{\Lambda }+1}`$, where $`V_i`$ is the irreducible $`sl_2\overline{sl}_2`$-module generated by $`b_i`$, for $`i10`$, and $`V_{10}`$ is generated by $`b_{10}+2\mathrm{\Lambda }b_1`$, and $`U^{\mu ,\mu ^{}}`$ is the irreducible $`sl_2\overline{sl}_2`$-module of highest weight $`(\mu ,\mu ^{})`$. Note these modules are all irreducible. Note further that they are all non-isomorphic. So as before to show that $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ is irreducible, it suffices to show that one may send a $`(𝔙_+sl_2\overline{sl}_2)`$-highest weight vector in any irreducible $`(𝔙_+sl_2\overline{sl}_2)`$-component to the irreducible component containing the $`𝔤`$-highest weight vectors. For this purpose we compute
$`G_{\frac{1}{2}}^{}b_2=2(\mathrm{\Lambda }+1)b_1,\overline{G}_{\frac{1}{2}}^+b_3=2\mathrm{\Lambda }(\mathrm{\Lambda }+1)b_1,`$
$`G_{\frac{1}{2}}^+b_4=2\mathrm{\Lambda }(\mathrm{\Lambda }+1)b_1,F_1b_6=2(\mathrm{\Lambda }+1)b_1,`$
$`\overline{F}_1b_9=2(\mathrm{\Lambda }+1)b_1\overline{F}_1(b_{10}+2\mathrm{\Lambda }b_1)=2\mathrm{\Lambda }\overline{F}_0b_1,`$
$`\overline{G}_{\frac{3}{2}}^{}b_{12}=8(\mathrm{\Lambda }+1)b_1.`$
This settles the case when $`\mathrm{\Lambda }2`$.
In the case when $`\mathrm{\Lambda }=1`$ $`N^{E_0,\overline{E}_0}`$ is generated over $`[]`$ by
$$S^\mathrm{\Lambda }=\{b_5,b_7+b_{10},b_{13},b_{14},b_{16}+b_7,b_{15}\}.$$
Therefore $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda },\mathrm{\Lambda })/N`$ contains a $``$-invariant (and hence $`𝔤`$-invariant) vector $`b_{15}`$. Since in this case the vectors $`b_8=b_{11}=0`$, $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda },\mathrm{\Lambda })/(N+b_{15})`$ as a $`𝔙_+sl_2\overline{sl}_2`$-module is isomorphic to $`[]V_1[]V_2[]V_3[]V_4[]V_6[]V_9[]V_{10}[]V_{12}`$. Every component is irreducible except for $`[]V_{12}`$, which contains a unique (irreducible) $`𝔙_+sl_2\overline{sl}_2`$-submodule isomorphic to $`L_{𝔙_+}(1)U^{2,2}`$ generated by the highest weight vector $`b_{15}`$. But then the above calculation plus the fact that
$$\overline{G}_{\frac{5}{2}}^{}b_{12}=24(\mathrm{\Lambda }+1)b_1$$
also shows that $`M_{𝔖_+^4}(\frac{\mathrm{\Lambda }+2}{2},\mathrm{\Lambda },\mathrm{\Lambda })/(N+b_{15})`$ is irreducible. ∎
We summarize the results of this section in the following theorem.
###### Theorem 7.1.
The modules $`L_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$, for $`\mathrm{\Delta }`$ and $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}_+`$, form a complete list of non-isomorphic finite (over $`[]`$) irreducible $`K(1,4)_+`$-modules. Furthermore $`L_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$ as a $`[]`$-module has rank
* $`8\mathrm{\Lambda }(\mathrm{\Lambda }+1)`$, in the case $`2\mathrm{\Delta }\mathrm{\Lambda }=0`$ and $`\mathrm{\Lambda }=\overline{\mathrm{\Lambda }}`$,
* $`8(\mathrm{\Lambda }+1)(\mathrm{\Lambda }+2)`$, in the case $`2\mathrm{\Delta }+\mathrm{\Lambda }+2=0`$ and $`\mathrm{\Lambda }=\overline{\mathrm{\Lambda }}`$.
* $`16(\mathrm{\Lambda }+1)(\overline{\mathrm{\Lambda }}+1)`$, in all other cases.
Furthermore the $`[]`$-rank of $`L_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})_{\overline{0}}`$ equals the $`[]`$-rank of $`L_{𝔖_+^4}(\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})_{\overline{1}}`$ in all cases.
###### Remark 7.2.
Again the translation into the languages of modules over conformal algebras and of conformal modules is straightforward and hence is omitted. We thus obtain that all finite irreducible modules over the “big” $`N=4`$ conformal superalgebra are of the form $`L_{𝔖^4}(\alpha ,\mathrm{\Delta },\mathrm{\Lambda },\overline{\mathrm{\Lambda }})`$, where $`\alpha ,\mathrm{\Delta }`$ and $`\mathrm{\Lambda },\overline{\mathrm{\Lambda }}_+`$. Again the definition of these modules and the action of the conformal superalgebra on them are easily derived from our explicit description of a $`[]`$-basis in this section. We note that the adjoint module is isomorphic to $`M_{𝔖^4}(0,0,0,0)`$. This module is not simple, since $`K(1,4)`$ is not a simple Lie superalgebra. Its derived algebra $`K(1,4)^{}`$ (which is a simple formal distribution Lie superalgebra) is an ideal in $`K(1,4)`$ of codimension $`1`$ . Thus the annihilation subalgebra of $`K(1,4)^{}`$ and $`K(1,4)`$ are identical, and hence their conformal modules are identical. Therefore the results in this section also give explicit description of irreducible conformal modules over $`K(1,4)^{}`$. We finally remark that the $`K(1,4)^{}`$ as a conformal module over $`K(1,4)`$ corresponds to $`L_{𝔖^4}(0,\frac{1}{2},1,1)`$.
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# References
The stopping of swift protons in matter and its implication for astrophysical fusion reactions
C.A. Bertulani<sup>(a,b)(∗)</sup> and D.T. de Paula<sup>(a)</sup><sup>1</sup><sup>1</sup>1E-mails: bertu@if.ufrj.br, dani@if.ufrj.br
<sup>(a)</sup>Instituto de Física, Universidade Federal do Rio de Janeiro, 21945-970 Rio de Janeiro, RJ, Brazil
<sup>(b)</sup>Brookhaven National Laboratory, Physics Department, Upton, NY 11973-5000, USA
> The velocity dependence of the stopping power of swift protons and deuterons in low energy collisions is investigated. At low projectile energies the stopping is mainly due to nuclear stopping and charge exchange of the electron. The second mechanism dominates after E$`{}_{p}{}^{}200`$ eV. A dynamical treatment of the charge exchange mechanism based on two-center electronic wavefunctions yields very transparent results for the exchange probability. We predict that the stopping cross sections vary approximately as v$`{}_{}{}^{1.35}{}_{p}{}^{}`$ for projectile protons on hydrogen targets in the 1 keV energy region.
Nuclear fusion reactions proceed in stars at extremely low energies, e.g., of the order of $`10`$ keV in our sun . At such low energies it is extremely difficult to measure the cross sections for charged particles at laboratory conditions due to the large Coulomb barrier. One often uses a theoretical model to extrapolate the experimental data to the low-energy region. Such extrapolations are sometimes far from reliable, due to unknown features of the low-energy region. E.g., there might exist unknown resonances along the extrapolation, or even some simple effect which one was not aware of before. One of these effects is the laboratory atomic screening of fusion reactions . It is well known that the laboratory measurements of low energy fusion reactions are strongly influenced by the presence of the atomic electrons. This effect has to be corrected for in order to relate the fusion cross sections measured in the laboratory with those at the stellar environment. Another screening effect, arising from free electrons in the stellar plasma, will not be treated here. For about one decade, until 1996, one observed a large discrepancy between the experimental data and the best models available to treat the screening effect. The simplest (and perhaps the best of these models), the so-called adiabatic model, predicts that as the projectile nucleus penetrates the electronic cloud of the target the electrons become more bound and the projectile energy increases by energy conservation. Since the fusion cross sections increase strongly with the projectile’s energy, this tiny amount of energy gain (of order of 10-100 eV) leads to a large effect on the measured cross sections. However, in order to explain the experimental data, it is necessary an extra-amount of energy - about twice the value obtained by the adiabatic model. This is puzzling, since more refined dynamical models, e.g., time-dependent Hartree-Fock , include electronic excitation and thus yield a screening energy which is smaller than that obtained with the adiabatic model.
This problem was apparently solved in 1996 by Langanke and collaborators and by Bang and collaborators , who observed that the experimental data for $`{}_{}{}^{3}He(d,`$ p$`)^4He`$ \- the reaction for which the screening effect was best studied - was probably obtained with a erroneous extrapolation of the stopping power for deuterons in helium targets to the low energy regime. The fusion reaction occurs at a point inside the target after the projectile has slowed down by interactions with the atomic targets. In the experimental analysis one needs to correct for this energy loss in order to assign the right projectile energy value for that reaction. These corrections were usually based on the Andersen-Ziegler table of stopping power of low energy particles . Due to the lack of experimental information on the stopping power at the extreme low projectile energies needed for astrophysical purposes, the Anderson-Ziegler tabulation was extrapolated to the required energy; another example of a dangerous extrapolation procedure. In fact, Golser and Semrad observed a strong departure of their experimental data from the extrapolations based on the Andersen-Ziegler tables for the stopping of low energy protons on helium targets. Grande and Schwietz performed a dynamical calculation of the energy dependence of the stopping power for this system and confirmed that the extrapolation procedure cannot be extended to the very low energies. Whereas at higher energies the stopping is mainly due to the ionization of the target electrons, at the astrophysical energies it is mainly due to charge-exchange between the target and the projectile. Refs. and use these arguments to explain the long standing discrepancy between theory and experiment for the low energy dependence of the reaction $`{}_{}{}^{3}He(d,`$ p$`)^4He`$. Other reactions of astrophysical interest (e.g., those listed in by Rolfs and collaborators ) should also be corrected for this effect.
In this work we address the problem of the stopping of very low energy ions in matter. To simplify matters, we study the system $`p+H`$, which is the simplest one can think of. It displays important features of the stopping power and has the advantage of allowing a very simple solution.
Our approach is based on the solution of the time-dependent Schrödinger equation for the electron in a dynamical two-center field. The static two-center $`p+H`$ system has been solved by Edward Teller in 1930 . He showed that as the distance between the protons decreases the hydrogen orbitals split into two or more orbitals, depending on its degeneracy in the two-center system. Analogous problems are well known in quantum systems . For example, take two identical potential wells at a certain distance. For large distances the states in one well are degenerated with the states in the other potential well. As they approach this degeneracy is lifted due to the influence of barrier tunneling. Thus, the lowest energy state of hydrogen, $`1`$s, splits into the 1s$`\sigma `$ and the 2p$`\sigma `$ states as the protons approach each other. The 1s$`\sigma `$ state is space symmetrical, while the 2p$`\sigma `$ state is antisymmetric. As the proton separation distance decreases their respective energies decrease. At $`R1`$ Å the energy of the 2p$`\sigma `$ state starts to increase again, while the energy of the 1s$`\sigma `$ state continues to decrease. For proton distances much smaller than $`1`$ Å the 1s$`\sigma `$ and the 2p$`\sigma `$ energies correspond to those of the first and second states of the He atom, respectively .
Let us now consider the dynamical case. The full time-dependent wavefunction for the system can be expanded in terms of two-center states, $`\varphi _n(t)`$, governed by the Schrödinger’s equation
$$\left[H_0+V_p\left(t\right)\right]\varphi _n(t)=E_n\left(t\right)\varphi _n(t),\mathrm{with}H_0=\widehat{𝐩}_e^2/2m_e+V_T,$$
(1)
where $`V_p\left(t\right)=e^2/\left|𝐫+𝐑/2\right|`$ is the electron-projectile proton interaction potential and $`V_T=e^2/\left|𝐫𝐑/2\right|`$ is the electron-target proton interaction for a proton-proton separation distance $`R(t)`$. Note that in our formulation the two-center wave functions also depend on time, as well as their energies $`E_n\left(t\right)`$. The full electronic wavefunction is obtained by a sum over all orthonormal two-center states
$$|\mathrm{\Psi }\left(t\right)=\underset{n}{}a_n\left(t\right)|\varphi _n\left(t\right),\mathrm{with}d^3r\varphi _n\left(t\right)\varphi _m\left(t\right)=\delta _{nm}.$$
(2)
Inserting this expansion in eq. (1) we obtain
$$i\mathrm{}\frac{d}{dt}a_m\left(t\right)=E_m\left(t\right)a_m\left(t\right)i\mathrm{}\underset{n}{}a_n\left(t\right)m\left|\frac{d}{dt}\right|n.$$
(3)
Using (1) one can easily show that, for $`mn`$,
$$m\left|\frac{d}{dt}\right|n=\frac{m\left|dV_p/dt\right|n}{E_n\left(t\right)E_m\left(t\right)},(mn).$$
(4)
Moreover, using the second relation of eq. (2), one can show that $`m\left|\frac{d}{dt}\right|m=0`$, if $`|m`$ is real. This indeed will be our case. Our basis, $`|n\left(t\right)`$, is formed by two-center states at a given time $`t`$, i.e., a given proton separation distance, $`R`$. These wavefunctions are real. Thus, the final coupled-channels equation for the two-center problem is given by
$$i\mathrm{}\frac{d}{dt}a_m\left(t\right)=E_m\left(t\right)a_m\left(t\right)i\mathrm{}\underset{mn}{}a_n\left(t\right)\frac{m\left|dV_p/dt\right|n}{E_n\left(t\right)E_m\left(t\right)}.$$
(5)
At very low proton energies $`\left(E_p1\mathrm{k}\mathrm{e}\mathrm{V}\right)`$ it is fair to assume that only the low-lying states are involved in the electronic dynamics. Only at proton energies of order of 25 keV the proton velocity is comparable to the electron velocity, v$`{}_{e}{}^{}\alpha c`$. Thus, the evolution of the system is almost adiabatic at $`E_p10`$ keV. The higher states require too much excitation energy and belong to different degeneracy multiplets. The initial electronic wavefunction is a clear superposition of 1s$`\sigma `$ and 2p$`\sigma `$ two-center states. One thus expects that only these states are relevant for the calculation. In fact, at these energies the population of the 2p atomic state in charge exchange is much less than the population of the 1s atomic state. These assumptions are well supported by the calculations of Grande and Schwietz , who have used a dynamical approach based on target-centered wavefunctions. In their approach one has to include a great amount of target-centered states in order to represent well the strong distortion of the wavefunction as the projectile closes in the target. We also have assumed that the proton follows a classical trajectory determined by an impact parameter $`b`$.
Eq. (5) does not look like the usual form of coupled-channels equations in the theory of the time-dependent Schrödinger equation. But we can put it in such form by rewriting the equation as
$$i\mathrm{}\frac{d}{dt}\left(\begin{array}{c}a_+\\ a_{}\end{array}\right)=\left(\begin{array}{cc}V_++E_0& iW\\ iW& V_{}+E_0\end{array}\right)\left(\begin{array}{c}a_+\\ a_{}\end{array}\right),$$
(6)
where the indices $`+`$ and $``$ refer to the $`1`$s$`\sigma `$ and $`2`$p$`\sigma `$ states, respectively, $`E_0=13.6`$ eV, $`V_\pm \left(t\right)=E_\pm (t)E_0`$, and
$$W\left(t\right)=\mathrm{}\frac{\mathrm{\Psi }_+\left|dV_p/dt\right|\mathrm{\Psi }_{}}{E_+\left(t\right)E_{}\left(t\right)}\mathrm{}\frac{\mathrm{\Psi }_{1\mathrm{s}\sigma }\left(t\right)\left|dV_p/dt\right|\mathrm{\Psi }_{2\mathrm{p}\sigma }\left(t\right)}{E_{1\mathrm{s}\sigma }\left(t\right)E_{2\mathrm{p}\sigma }\left(t\right)}.$$
(7)
In this form, the potentials $`V_\pm \left(t\right)`$ and $`W\left(t\right)`$ act like potentials in usual coupled-channels equations. We use the formalism of Teller to calculate the wavefunctions $`\mathrm{\Psi }_\pm \left(R\right)`$ at different inter-proton distances, $`R(t)`$, corresponding to a particular time $`t`$. The static Schrödinger equation is solved in elliptical coordinates. This yields two coupled differential equations which can be solved by expanding the solutions in Taylor series. A set of recurrence relations is obtained for the expansion coefficients when the boundary conditions are used. The energies $`E_{1\mathrm{s}\sigma }\left(R\right)`$ and $`E_{2\mathrm{p}\sigma }\left(R\right)`$ are obtained by adjusting the constant which separates the two coupled equations to its correct matching value.
When $`t\pm \mathrm{}`$, $`V_\pm 0`$ and $`W0`$. The initial state, an electron localized in the target can be written in terms of the degenerate symmetric, $`\mathrm{\Psi }_+=`$ $`\mathrm{\Psi }_{1\mathrm{s}\sigma }`$, and anti-symmetric, $`\mathrm{\Psi }_{}=`$ $`\mathrm{\Psi }_{2\mathrm{p}\sigma }`$, states:
$$\mathrm{\Phi }_T=\frac{1}{\sqrt{2}}\left(\mathrm{\Psi }_++\mathrm{\Psi }_{}\right),\mathrm{at}t\mathrm{},$$
(8)
where both $`\mathrm{\Phi }_T`$, and $`\mathrm{\Psi }_\pm `$ are normalized wavefunctions. If the electron is localized in the projectile, the wavefunction $`\mathrm{\Phi }_p=\left(\mathrm{\Psi }_+\mathrm{\Psi }_{}\right)/\sqrt{2},`$when $`t\mathrm{}`$ is used. We will consider only the condition of eq.$`\left(\text{8}\right),`$ namely, an electron localized at the target at $`t\mathrm{}`$. These relations are well known quantum mechanical results; the asymptotic two-center wavefunctions can be written as combinations of target- and projectile-centered 1s-wavefunctions: $`\mathrm{\Psi }_\pm =\left(\mathrm{\Phi }_p\pm \mathrm{\Phi }_T\right)/\sqrt{2}`$.
Starting with a target localized electron we assign the initial conditions $`a_\pm =1/\sqrt{2}`$ at $`t\mathrm{}`$ and solve the equation $`\left(\text{6}\right)`$ numerically. Although at $`t\mathrm{}`$ the probabilities $`\left|a_\pm \right|^2`$ remain very close to 1/2, the amplitudes $`a_\pm `$ acquire phases which change the relative population of the projectile and the target 1s state. We correct for energy conservation which feeds the increasing binding energy of the electron back to an increasing relative motion energy of the two protons as they come closer. This is specially important as $`E_p`$ becomes of order of hundreds of eV, and smaller. In figure 1 we show the time dependence of $`V_\pm \left(t\right)`$ and $`W\left(t\right)`$ for $`E_p=10`$ keV and a nearly central collision, $`b=0.1`$ Å. One observes that the potentials $`V_\pm \left(t\right)`$ extend much farther out than $`W\left(t\right).`$ Moreover, we find that as $`E_p`$ decreases the potential $`W`$ decreases faster than the projectile’s velocity, $`v_p`$. This is mainly due to the derivative of in eq. (7). At $`E_p100`$ eV the potential $`W`$ loses its relevance as compared to $`V_\pm `$, which have no dependence on $`v_p`$. This becomes clear in figure 2. In this figure we show the exchange probability as a function of the impact parameter for two projectile energies. The solid line is the full solution of eq. $`\left(\text{6}\right).`$ The dashed line is the approximation obtained when we set $`W=0`$ in eq. $`\left(\text{6}\right).`$ In the later case, the equations decouple and it is straightforward to show that the exchange probability is given by
$$P_{exch}=\left|\underset{\pm }{}a_\pm \left(\mathrm{}\right)\mathrm{\Phi }_T|\mathrm{\Psi }_\pm \left(\mathrm{}\right)\right|^2=\frac{1}{2}+\frac{1}{2}\mathrm{cos}\left\{\frac{1}{\mathrm{}}_{\mathrm{}}^{\mathrm{}}\left[E_{}\left(t\right)E_+\left(t\right)\right]𝑑t\right\}.$$
(9)
At $`E_p=10`$ keV there is an appreciable difference between the full calculation and the approximation (9). But, for $`E_p=100`$ eV the results are practically equal, except for very small impact parameters at which the potential $`W`$ still has an effect.
One observes that the exchange probability is not constant at small impact parameters, but oscillates wildly around 0.5, specially for low projectile energies. One might naively assume that because the collision is almost adiabatic, the system loses memory of to which nucleus the electron is bound after the collision. Thus, for small impact parameters one would expect a 50% probability of finding the electron in one of the nuclei at $`t=\mathrm{}`$. However, this is not what happens. From eq. (9) we see that minima of the probability occur for impact parameters satisfying the relation
$$_{\mathrm{}}^{\mathrm{}}\left[E_{}\left(t\right)E_+\left(t\right)\right]𝑑t=2\pi \mathrm{}\left(n+1/2\right),n=0,1,2,\mathrm{},N.$$
(10)
This relation looks familiar, of course. It simply states that the interference between the 1s$`\sigma `$ and the 2p$`\sigma `$ states induces oscillations in the exchange probability. The electron tunnels back and forth between the projectile and the target during the ingoing and the outgoing part of the trajectory. When the interaction time is an exact multiple of the oscillation time, a minimum in the exchange probability occurs. The average probability over the smaller impact parameters is indeed 0.5. As the impact parameter decreases from infinity, the first maximum in the exchange probability indicates the beginning of the region of strong exchange probability. One sees that at low proton energies this starts at $`b3`$ Å. The size of the hydrogen atom is about 0.5 Å and thus the electron travels in a forbidden region (tunnels) of about 2 Å from the target to the projectile. This is possible because of the strong interference between the 1s$`\sigma `$ and the 2p$`\sigma `$ states, which for some trajectories satisfy the quantum relation (10).
To obtain the stopping power we need the total cross section for charge exchange, $`\sigma =2\pi P_{exch}b𝑑b`$. This is shown in figure 3. The solid line is the full coupled-channels calculation, while the dashed line uses approximation (9) for the exchange probability. We observe that the approximation (9) reproduces well the full calculation even at the highest energies. The reason is that the potential $`W`$ is always smaller than $`V_\pm `$ for large impact parameters which have more weight on the integral cross section. We also compare our calculations with the lowest energy data of McClure . The formalism developed here is inappropriate for energies in the tens of keV range and higher, as the projectile velocity becomes comparable to or higher than the electron velocity. This implies that two-center states with higher energy and even continuum states (ionization) should be included in the calculation. For $`E_p0`$, the charge exchange cross section becomes the constant value $`\sigma \left(E_p=0\right)=37.88`$ $`\times 10^{16}`$ cm<sup>2</sup>. This happens because, when $`E_p0`$ and as the projectile nears the targets, the increasing electron binding in the two-center system acts as a push in the relative motion energy to compensate for energy conservation. The average result is that the cross section for charge exchange becomes approximately constant for projectile energies of tens of eV and below.
In figure 4 we show the stopping cross section of the proton. The stopping cross section is defined as $`S=_i\mathrm{\Delta }E_i\sigma _i`$ , where $`\mathrm{\Delta }E_i`$ is the energy loss of the projectile in a process denoted by $`i`$. The stopping power, $`S_P=dE/dx`$, the energy loss per unit length of the target material, is related to the stopping cross section by $`S=S_P/N`$, where $`N`$ is the atomic density of the material. In our charge exchange mechanism the electron is transferred to the ground state of the projectile and the energy transfer is given by $`\mathrm{\Delta }E=m_ev_p^2/2`$, where $`v_p`$ is the projectile velocity. Assuming that there is a few free electrons in the material (e.g., in a hydrogen gas) only one more stopping mechanism at very low energies should be considered: the nuclear stopping power. This is simply the elastic scattering of the projectile off the target nuclei. The projectile energy is partially transferred to the recoil energy of the target atom. The stopping cross section for this mechanism has been extensively studied by Lindhard and collaborators (see, e.g., ref. ). The nuclear stopping includes the effect of the electron screening of the nuclear charges.
The dotted line in figure 4 gives the energy transfer by means of nuclear stopping, while the solid line are our results for the charge-exchange stopping mechanism. The data points are from the tabulation of Andersen and Ziegler . We see that the nuclear stopping dominates at the lowest energies, while the charge-exchange stopping is larger for proton energies greater than 200 eV. Since we neglect the difference between molecular and atomic hydrogen targets, there is a limitation to compare our results with the experimental data. But, the order of magnitude agreement is very good in view of our simplifying assumptions. We do not consider the change of the charge state of the protons as they penetrate the target material. The exchange mechanism transforms the protons into H atoms. These again interact with the target atoms. The can loose their electron again by transfer to the 1s state of the target .
The best fit to our calculation for the stopping power for proton energies in the range 100 eV - 1 keV yields $`Sv_p^{1.35}`$. This contrasts with the extrapolation $`Sv_p`$, based on the Andersen-Ziegler table. But, this discrepancy is much less than the one obtained by Golser and Semrad for helium targets, who found a stopping power for protons $`Sv_p^{3.34}`$ for protons in the energy range of 4 keV. No data at lower energies are available in this case. But, the Golser and Semrad data, for proton energies above 3 keV, firmly indicate that a high power dependence on the projectile velocity will be also valid at lower energies, in contrast the predictions from the Andersen and Ziegler tables . One cannot extend our calculations to helium targets as the initial wavefunction cannot be described in terms of a simple sum of two-center states. A much larger two-center basis is necessary. Since the electrons in the helium target are more bound than in the proton, the charge-exchange probability must be much smaller than in the case of hydrogen targets. One thus should indeed expect a much stronger dependence of the stopping on the projectile velocity. At very low energies, of the order of some hundreds of eV, the stopping cross section should be entirely dominated by nuclear stopping, even more than for hydrogen targets.
The $`p+pd+e^++\nu _e`$ reaction is a very important one occurring in, e.g., our sun. But, it proceeds via the weak interaction and its cross section is extremely small for studies under the laboratory conditions . Fortunately, a good theoretical model exists for this reaction . Other reactions could be strongly influenced by the stopping power of protons and deuterons due to the charge-exchange mechanism. They can be relevant for the study of $`d+D`$ reactions in stellar interiors and fusion reactors. Another application is the D(p, $`\gamma `$)$`{}_{}{}^{3}He`$ reaction which is important for the hydrogen burning in stars. In our sun the most effective energy of this reaction is $`E_{c.m.}=6.5\pm 3.3`$ keV at $`T=15\times 10^6`$ K. At this energy one expects that the charge-exchange stopping cross section should be as important as the ionization cross section. Experimental data exist at the lowest energy value of 16 keV . Although the extrapolation based on theory appears to be under control in this case, it is worthwhile to consider a better study of the stopping power for this reaction. The steep rise of the fusion cross sections at astrophysical energies amplifies all effects leading to a slight modification of the projectile energy . Our results show that the stopping mechanism does not follow a universal pattern for all systems. This calls for improved theoretical studies of charge-exchange effects and for their independent experimental verification.
Acknowledgments We would like to express our gratitude to profs. A.B. Balantekin, S.R. Souza and L.F. Canto for useful comments and suggestions during the development of this work. This work was partially supported by the Brazilian agencies: CNPq, FAPERJ, FUJB, and by the MCT/ FINEP/CNPq(PRONEX) (contract 41.96.0886.00). (\*) John Simon Guggenheim fellow.
Figure Captions
Fig. 1 \- Time dependence of the interaction potentials $`V_\pm \left(t\right)`$ and $`W\left(t\right)`$ for $`E_p=10`$ keV and a nearly central collision, $`b=0.1`$ Å.
Fig. 2 \- The exchange probability as a function of the impact parameter for two projectile energies. The solid line is the full solution of eq. $`\left(\text{6}\right).`$ The dashed line is the approximation obtained when we set $`W=0`$ in eq. $`\left(\text{6}\right).`$
Fig. 3 \- The solid line is the full coupled-channels calculation for the charge-exchange cross section, while the dashed line uses approximation (9) for the exchange probability. The experimental data are from McClure .
Fig. 4 \- The stopping cross section of protons on H-targets. The dotted line in gives the energy transfer by means of nuclear stopping, while the solid line are our results for the charge-exchange stopping mechanism. The data points are from the tabulation of Andersen and Ziegler .
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# Stochastic Monotonicity and Realizable Monotonicity
## 1 Introduction
We will discuss two notions of monotonicity for probability measures on a finite partially ordered set (poset). Let $`𝒮`$ be a finite poset and let $`(P_1,P_2)`$ be a pair of probability measures on $`S`$. (We use a calligraphic letter $`𝒮`$ in order to distinguish the set $`S`$ from the same set equipped with a partial ordering $``$.) A subset $`U`$ of $`S`$ is said to be an up-set in $`𝒮`$ (or increasing set) if $`yU`$ whenever $`xU`$ and $`xy`$. We say that $`P_1`$ is stochastically smaller than $`P_2`$, denoted
(1.1)
$$P_1P_2,$$
if
(1.2)
$$P_1(U)P_2(U)\text{ for every up-set }U\text{ in }𝒮\text{.}$$
The relation introduced in (1.1)–(1.2) is clearly reflexive and transitive. Antisymmetry follows easily using our assumption that $`S`$ is finite, so the relation defines a partial ordering on the class of probability measures on $`S`$. (For a careful discussion on the matter of antisymmetry in a rather general setting for infinite $`S`$, see .)
The following characterization of stochastic ordering was established by Strassen and fully investigated by Kamae, Krengel, and O’Brien . Suppose that there exists a pair $`(𝐗_1,𝐗_2)`$ of $`S`$-valued random variables \[defined on some probability space $`(\mathrm{\Omega },,)`$\] satisfying the properties
(1.3)
$$𝐗_1𝐗_2$$
and
(1.4)
$$(𝐗_i)=P_i()\text{ for }i=1,2\text{}$$
Then we have
$$P_1(U)=(𝐗_1U)=(𝐗_1U,𝐗_1𝐗_2)(𝐗_2U)=P_2(U),$$
for every up-set $`U`$ in $`𝒮`$. Thus, the conditions (1.3)–(1.4) necessitate (1.1). Moreover, Strassen’s work shows that (1.1) is in fact sufficient for the existence of a probability space $`(\mathrm{\Omega },,)`$ and a pair $`(𝐗_1,𝐗_2)`$ of $`S`$-valued random variables on $`(\mathrm{\Omega },,)`$ satisfying (1.3)–(1.4). \[Equivalently, we need only require that (1.3) hold almost surely.\]
Now let $`𝒜`$ be a finite poset. Let $`(P_\alpha :\alpha A)`$ be a system of probability measures on $`S`$. We call $`(P_\alpha :\alpha A)`$ a realizably monotone system if there exists a system $`(𝐗_\alpha :\alpha A)`$ of $`S`$-valued random variables defined on some probability space $`(\mathrm{\Omega },,)`$ such that
(1.5)
$$𝐗_\alpha 𝐗_\beta \text{ whenever }\alpha \beta \text{ }$$
and
(1.6)
$$(𝐗_\alpha )=P_\alpha ()\text{ for every }\alpha A\text{}$$
In such a case we shall say that $`(𝐗_\alpha :\alpha A)`$ realizes the monotonicity of $`(P_\alpha :\alpha A)`$. Since the conditions (1.3)–(1.4) imply (1.1), the conditions (1.5)–(1.6) applied pairwise imply
(1.7)
$$P_\alpha P_\beta \text{ whenever }\alpha \beta \text{}$$
The system $`(P_\alpha :\alpha A)`$ is said to be stochastically monotone if it satisfies (1.7). We have shown that stochastic monotonicity is necessary for realizable monotonicity.
In light of Strassen’s characterization of stochastic ordering, one might guess that stochastic monotonicity is also sufficient for realizable monotonicity. It is perhaps surprising that the conjecture is false in general, as the following example shows.
###### Example 1.1.
Let
(1.8)
$$𝒮=𝒜:=\text{}$$
be the usual 2-dimensional Boolean algebra with $`x<y,z`$ and $`y,z<w`$ (and, of course, $`x<w`$ by transitivity). Define a system $`(P_x,P_y,P_z,P_w)`$ of probability measures on $`S`$ by
$$P_\xi :=\{\begin{array}{cc}\text{unif}\{x,y\}\hfill & \text{ if }\xi =x\text{}\hfill \\ \text{unif}\{x,w\}\hfill & \text{ if }\xi =y\text{}\hfill \\ \text{unif}\{y,z\}\hfill & \text{ if }\xi =z\text{}\hfill \\ \text{unif}\{y,w\}\hfill & \text{ if }\xi =w\text{}\hfill \end{array}$$
where $`\text{unif}(B)`$ denotes the uniform probability measure on a set $`B`$. Clearly, $`(P_x,P_y,P_z,P_w)`$ is stochastically monotone. Now suppose that there exists a system $`(𝐗_x,𝐗_y,𝐗_z,𝐗_w)`$ which realizes the monotonicity of $`(P_x,P_y,P_z,P_w)`$. Considering the event $`𝐗_x=y`$, realizable monotonicity forces
$$(𝐗_x=y)=(𝐗_x=y,𝐗_y=w,𝐗_z=y,𝐗_w=w)=\frac{1}{2}.$$
Similarly, we find
$$(𝐗_z=z)=(𝐗_x=x,𝐗_z=z,𝐗_w=w)=\frac{1}{2}.$$
Noting that the above two events are disjoint, we conclude $`(𝐗_w=w)=1`$, which is a contradiction. Thus, $`(P_x,P_y,P_z,P_w)`$ cannot be realizably monotone.
Given a pair $`(𝒜,𝒮)`$ of posets, if the two notions of monotonicity—stochastic and realizable—are equivalent, then we say that monotonicity equivalence holds for $`(𝒜,𝒮)`$. The counterexample in Example 1.1 was discovered independently by Ross ; we are grateful to Robin Pemantle for pointing this out to us. Ross reduced the question of monotonicity equivalence for general infinite posets $`𝒜`$ (and given $`𝒮`$) to consideration of the same question for every finite induced subposet of $`𝒜`$. Thus we regard our work as a useful complement to his. As an historical aside, we note that what we call a realizably monotone system, Ross called a coherent family.
To give an example where monotonicity equivalence holds, we next consider the case where $`𝒮`$ is a linearly ordered set.
###### Example 1.2.
Let $`𝒜`$ be any poset and let $`𝒮`$ be a linearly ordered set. Suppose that $`(P_\alpha :\alpha A)`$ is a stochastically monotone system of probability measures on $`S`$. For each $`\alpha A`$, define the inverse probability transform $`P_\alpha ^1`$ by
(1.9)
$$P_\alpha ^1(t):=\mathrm{min}\left\{xS:t<F_\alpha (x)\right\}\text{ for }t[0,1)\text{,}$$
where $`F_\alpha `$ is the distribution function of $`P_\alpha `$ \[i.e., $`F_\alpha (x):=P_\alpha (\{\xi S:\xi x\})`$ for each $`xS`$\]. Given a single uniform random variable $`𝐔`$ on $`[0,1)`$, we can construct a system $`(𝐗_\alpha :\alpha A)`$ of $`S`$-valued random variables via
$$𝐗_\alpha :=P_\alpha ^1(𝐔)\text{ for each }\alpha A\text{.}$$
Then $`(𝐗_\alpha :\alpha A)`$ realizes the monotonicity and therefore $`(P_\alpha :\alpha A)`$ is realizably monotone. Thus monotonicity equivalence holds for $`(𝒜,𝒮)`$.
The goal of our investigation is to determine for precisely which pairs $`(𝒜,𝒮)`$ of posets monotonicity equivalence holds. Let us discuss here the usefulness of such a determination. It is (structurally) simple to say which systems $`(P_\alpha :\alpha A)`$ are stochastically monotone. Indeed, one need only determine all up-sets $`U`$ of $`𝒮`$, and then $`(P_\alpha :\alpha A)`$ is stochastically monotone if and only if $`P_\alpha (U)P_\beta (U)`$ for all such $`U`$ whenever $`\alpha \beta `$. For realizable monotonicity, an analogous result is Theorem 2.9, but the necessary and sufficient condition there involves an infinite collection of inequalities. We know how to reduce, for each $`(𝒜,𝒮)`$, the infinite collection to a finite one, but (1) there seems to be in general no nice structural characterization of the resulting finite collection, and (2) the computations needed to do the reduction can be massive even for fairly small $`𝒜`$ and $`𝒮`$ (Chapter 7 of ).
Thus, when monotonicity equivalence holds, we learn that realizable monotonicity has the same simple structure as stochastic monotonicity. And when monotonicity equivalence fails, we learn that testing a system $`(P_\alpha :\alpha A)`$ for stochastic monotonicity does not suffice as a test for realizable monotonicity. In the latter case, for fixed $`(𝒜,𝒮)`$ and a single numerically specified system $`(P_\alpha :\alpha A)`$, we can determine whether or not the system is realizably monotone by constructing a system $`(𝐗_\alpha :\alpha A)`$ subject to the marginal condition (1.6) so as to maximize the probability $`(𝐗_\alpha 𝐗_\beta \text{ whenever }\alpha \beta )`$. Indeed, the system $`(P_\alpha :\alpha A)`$ is realizably monotone if and only if the maximum value equals $`1`$. The construction can be carried out using linear programming with variables corresponding to the values of the joint probability mass function for $`(𝐗_\alpha :\alpha A)`$.
For further discussion along these lines, see .
Of particular interest in our study of realizable monotonicity is the case $`𝒜=𝒮`$. Here the system $`(P(x,):xS)`$ of probability measures can be considered as a Markov transition matrix $`𝐏`$ on the state space $`S`$. Recently, Propp and Wilson and Fill have introduced algorithms to produce observations distributed perfectly according to the long-run distribution of a Markov chain. Both algorithms apply most readily and operate most efficiently when the state space $`𝒮`$ is a poset and a suitable monotonicity condition holds. Of the many differences between the two algorithms, one is that the appropriate notion of monotonicity for the Propp–Wilson algorithm is realizable monotonicity, while for Fill’s algorithm it is stochastic monotonicity; see Remark 4.5 in . Here the properties (1.5)–(1.6) are essential for the Propp–Wilson algorithm to be able to generate transitions simultaneously from every state in such a way as to preserve ordering relations. For further discussion of these perfect sampling algorithms, see and . In Theorem 4.3 we show that the two notions of monotonicity are equivalent if and only if the poset $`𝒮`$ is “acyclic”, which is characterized by possession of a Hasse diagram (the standard graphical representation of partial ordering) that is cycle-free. For example, the Hasse diagram of a linearly ordered set is a vertical path such as the one in Figure 6.1(b), and therefore has no cycle. On the other hand, the 2-dimensional Boolean algebra whose Hasse diagram is displayed in (1.8) is not acyclic. See Section 2.1 for precise terminology.
In the present paper we study the notion of realizable monotonicity when $`𝒜`$ and $`𝒮`$ are both finite posets. In Section 2.3 we review a general result for the existence of a probability measure with specified marginals and present the extensibility problem. Sections 2.12.2 are prepared to introduce definitions and several key notions in studying posets. In Section 2.4 we formulate the monotonicity equivalence problem from the viewpoint of the extensibility problem. Section 3 is rather short, introducing four subclasses—Classes B, Y, W and Z—that partition the class of connected posets $`𝒮`$.
At this juncture we intend to provide the reader with an overview of the main results of this paper. In Section 4 we present the first case of our investigation, where $`𝒮`$ is in Class B. Kamae, Krengel, and O’Brien showed that if $`𝒜`$ is a linearly ordered set then monotonicity equivalence holds for $`(𝒜,𝒮)`$. We generalize this result (in our finite setting) to the case of an acyclic poset $`𝒜`$ (Theorem 4.1); see Section 2.1 for the definition of an acyclic poset. Theorem 4.2 gives an exact answer to our central question of monotonicity equivalence when $`𝒮`$ is a poset of Class B. In Section 5 we proceed to the second case of our investigation, where $`𝒮`$ is in Class Y. There we first present Proposition 5.2 and discuss its proof in Section 5.1. This turns out to be a useful result concerning probability measures on a poset of Class Y, leading to Theorem 5.1, which in turn answers our monotonicity equivalence question when $`𝒮`$ is a poset of Class Y. If $`𝒮`$ is a poset of Class Z, then we can show that monotonicity equivalence holds for any poset $`𝒜`$ (Theorem 6.1). In Section 6 we give the proof by using a generalization of inverse probability transform. When $`𝒮`$ is a poset of Class W, we have devised a further generalization of inverse probability transform, which results in constructing a rather large class of posets $`𝒜`$ for which monotonicity equivalence holds. We refer the reader to for the results of our investigation of Class W.
## 2 Posets and the monotonicity equivalence problem
In Section 2.1 we briefly summarize the material on posets that we need for our study. An important assertion is that if a poset is non-acyclic then the poset has an induced cyclic subposet. In Section 2.2 we prove this (Lemma 2.4) among other results concerning induced cyclic posets. In Section 2.3 we review the well-known results of Strassen on the existence of a probability measure with specified marginals; our review is tailored somewhat to fit our application to realizable monotonicity. In Section 2.4 we discuss realizable monotonicity in terms of the existence of a probability measure with specified marginals. Propositions 2.142.15 are presented in Section 2.4; these allow $`𝒜`$ and $`𝒮`$ both to be connected posets in our later investigations.
### 2.1 Posets
We devote this subsection to introducing definitions and notation related to partial ordering. By a poset $`𝒮`$ we shall mean a finite set $`S`$ (the qualifier will not again be explicitly noted) together with a partial ordering $``$. The (unordered) set $`S`$ is called the ground set of $`𝒮`$. Most of the basic poset terminology adopted here can be found in Stanley or Trotter . Throughout this subsection, $`𝒮`$ and $`𝒮^{}`$ denote posets.
* Dual poset, up-set, down-set. The dual of $`𝒮`$, denoted $`𝒮^{}`$, is the poset on the same ground set $`S`$ as $`𝒮`$ such that $`x_1x_2`$ in $`𝒮^{}`$ if and only if $`x_1x_2`$ in $`𝒮`$. A subset $`U`$ of $`S`$ is said to be an up-set (or increasing set) in $`𝒮`$ if $`yU`$ whenever $`xU`$ and $`xy`$. A down-set $`V`$ in $`𝒮`$ is defined to be an up-set in $`𝒮^{}`$. Note that $`U`$ is an up-set in $`𝒮`$ if and only if $`SU`$ is a down-set in $`𝒮`$. For any subset $`B`$ of $`S`$, we can define the down-set $`B`$ generated by $`B`$ in the usual fashion:
$$B:=\{\xi S:\xi \eta \text{ for some }\eta B\}.$$
We simply write $`x_1,\mathrm{},x_n`$ for $`\{x_1,\mathrm{},x_n\}`$.
* Cover graph. For $`x,yS`$, we say that $`y`$ covers $`x`$ if $`x<y`$ in $`𝒮`$ and no element $`z`$ of $`S`$ satisfies $`x<z<y`$. The Hasse diagram of a poset is the directed graph whose vertices are the elements of the poset and whose arcs are those ordered pairs $`(x,y)`$ such that $`y`$ covers $`x`$. \[By convention, if $`y`$ covers $`x`$, then $`y`$ is drawn above $`x`$ in the Hasse diagram (as represented in the plane); this indicates the direction of each arc.\] We define the cover graph $`(S,_𝒮)`$ of $`𝒮`$ by considering the Hasse diagram of $`𝒮`$ as an undirected graph. That is, the edge set $`_𝒮`$ consists of those unordered pairs $`\{x,y\}`$ such that either $`x`$ covers $`y`$ or $`y`$ covers $`x`$ in $`𝒮`$.
* Subposet. We shall need to distinguish among several, somewhat subtly different, notions of subposet. We say that a poset $`𝒮^{}`$ is a subposet of $`𝒮`$ if $`S^{}`$ is (or, by extension and when there is no possibility of confusion, is isomorphic to) a subset of $`S`$ and $`xy`$ in $`𝒮^{}`$ implies $`xy`$ in $`𝒮`$ for $`x,yS^{}`$. When we speak of an induced subposet $`𝒮^{}`$ of $`𝒮`$, we mean that for $`x,yS^{}`$ we have $`xy`$ in $`𝒮^{}`$ if and only if $`xy`$ in $`𝒮`$. On the other hand, we call a (not necessarily induced) subposet $`𝒮^{}`$ a subposet via induced cover subgraph of $`𝒮`$ if $`y`$ covers $`x`$ in $`𝒮^{}`$ for $`x,yS^{}`$ precisely when $`y`$ covers $`x`$ in $`𝒮`$, that is, when the cover graph $`(S^{},_𝒮^{})`$ of $`𝒮^{}`$ is an induced subgraph of the cover graph $`(S,_𝒮)`$ of $`𝒮`$. Clearly, a subposet via induced cover subgraph of $`𝒮`$ with ground set $`S^{}`$ is a subposet of the subposet of $`𝒮`$ induced by ground set $`S^{}`$. In Example 2.1 we illustrate differences between these two notions of subposet.
Let $`(S^{},^{})`$ be a (not necessarily induced) subgraph of the cover graph $`(S,_𝒮)`$ of $`𝒮`$. Then $`(S^{},^{})`$ is the cover graph $`(S^{},_𝒮^{})`$ of a (not necessarily induced) subposet $`𝒮^{}`$ of $`𝒮`$. Here, $`y`$ covers $`x`$ in $`𝒮^{}`$ if and only if $`y`$ covers $`x`$ in $`𝒮`$ and $`\{x,y\}^{}`$. In this sense, a subgraph $`(S^{},^{})`$ of $`(S,_𝒮)`$ can be considered as a subposet of $`𝒮`$.
* Chain, height. We call a poset $`𝒮`$ a chain if any two elements of $`S`$ are comparable in $`𝒮`$. When we say that a subposet $`𝒮^{}`$ is a chain in $`𝒮`$, we mean that $`𝒮^{}`$ is a chain and an induced subposet of $`𝒮`$. The height $`n`$ of a poset $`𝒮`$ is the number of elements in a maximum-sized chain in $`𝒮`$. That is, $`𝒮`$ has height $`n`$ if and only if $`𝒮`$ has an $`n`$-element chain, but no $`(n+1)`$-element chain, as an induced subposet.
* Path, upward path, downward path. We call a (not necessarily induced) subposet $`𝒮^{}`$ of $`𝒮`$ a path if the cover graph $`(S^{},_𝒮^{})`$ of $`𝒮^{}`$ is (i) a path (in the usual graph-theoretic sense) and (ii) a (not necessarily induced) subgraph of the cover graph $`(S,_𝒮)`$ of $`𝒮`$. A sequence $`(x_0,x_1,\mathrm{},x_{n1})`$ denotes a path from $`x_0`$ to $`x_{n1}`$ with vertex set $`\{x_0,x_1,\mathrm{},x_{n1}\}`$ and edge set $`\{\{x_{i1},x_i\}:1in1\}`$. We say that a path $`(x_0,\mathrm{},x_{n1})`$ is upward (respectively, downward) in $`𝒮`$ if $`x_i`$ covers $`x_{i1}`$ in $`𝒮`$ ($`x_{i1}`$ covers $`x_i`$, respectively) for each $`i=1,\mathrm{},n1`$. Note that any upward or downward path in $`𝒮`$ is a chain in $`𝒮`$, but that the converse is not true. We illustrate chains and paths in Example 2.1(iv)–(vi).
* Cycle. We call a (not necessarily induced) subposet $`𝒮^{}`$ of $`𝒮`$ a cycle (or a cyclic subposet) if the cover graph $`(S^{},_𝒮^{})`$ of $`𝒮^{}`$ is (i) a cycle (in the usual graph-theoretic sense) and (ii) a (not necessarily induced) subgraph of the cover graph $`(S,_𝒮)`$ of $`𝒮`$. In Example 2.1 we demonstrate that a cyclic subposet may be neither an induced subposet nor a subposet via induced cover subgraph. If a poset $`𝒮`$ has a cyclic subposet, then we call $`𝒮`$ a non-acyclic poset. In keeping with the foregoing definitions, we call the reference poset $`𝒮`$ itself a cycle if the cover graph $`(S,_𝒮)`$ of $`𝒮`$ is a cycle. A sequence $`(x_0,x_1,\mathrm{},x_{n1},x_0)`$ with $`n4`$ denotes a cycle with vertex set $`\{x_0,x_1,\mathrm{},x_{n1}\}`$ and edge set $`\{\{x_{i1},x_i\}:1in\}`$. Here, indices are interpreted modulo $`n`$. (Note that a cyclic subposet must consist of at least four elements.)
* Connected poset, disjoint union. We say that $`𝒮`$ is connected if its cover graph $`(S,_𝒮)`$ is connected. The components of $`𝒮`$ are its maximal connected induced subposets. If $`S`$ and $`S^{}`$ are disjoint, then we can construct the disjoint union of $`𝒮`$ and $`𝒮^{}`$, denoted $`𝒮+𝒮^{}`$, as a poset on the ground set $`SS^{}`$ by declaring $`xy`$ in $`𝒮+𝒮^{}`$ precisely when either (i) $`x,yS`$ and $`xy`$ in $`𝒮`$, or (ii) $`x,yS^{}`$ and $`xy`$ in $`𝒮^{}`$. Thus any poset $`𝒮`$ is the disjoint union of its components.
* Acyclic poset, leaf. We say that a poset $`𝒮`$ is acyclic if $`𝒮`$ has no cyclic subposet. We call an element $`x`$ of $`S`$ a leaf in $`𝒮`$ if the edge set $`_𝒮`$ of the cover graph of $`𝒮`$ has a unique element $`\{x,y\}`$ for some $`yS`$. Note that if $`x`$ is a leaf in $`𝒮`$ then $`x`$ must be either maximal or minimal in $`𝒮`$. If $`𝒮`$ is a connected acyclic poset with $`|S|2`$, then there are at least two leaves in $`𝒮`$ (see, e.g., ).
* Poset-isomorphism, some named posets, subdivision. $`𝒮`$ is said to be poset-isomorphic to $`𝒮^{}`$ if there exists an bijection $`\varphi `$ from $`S`$ to $`S^{}`$ such that $`xy`$ in $`𝒮`$ if and only if $`\varphi (x)\varphi (y)`$ in $`𝒮^{}`$. In this paper, we call the $`2`$-dimensional Boolean algebra a diamond. \[See the figure in (1.8).\] Furthermore, we call the posets of Figure 2.1 and their duals (a) the bowtie, (b) Y-posets, (c) W-posets, and (d) the $`k`$-crown, respectively. The bowtie is the same as the $`2`$-crown. We may simply call $`𝒮`$ a crown if $`𝒮`$ is the $`k`$-crown for some $`k2`$. We say that a poset $`𝒮^{}`$ is a subdivision of $`𝒮`$ if the induced subposet of $`𝒮^{}`$ on ground set $`\{zS^{}:xzy\text{ in }𝒮^{}\}`$ is a chain whenever $`y`$ covers $`x`$ in $`𝒮`$.
###### Example 2.1.
Let
$`𝒮:=`$
be a poset. Here we give a number of examples to illustrate subtle distinctions in the definitions of subposets, paths, and cycles. Let
$`𝒮_1^{}:=`$ $`𝒮_2^{}:=`$ and $`𝒮_3^{}:=`$
be subposets of $`𝒮`$. Then, (i) $`𝒮_1^{}`$ is an induced subposet but not a subposet via induced cover subgraph, (ii) $`𝒮_2^{}`$ is a subposet via induced cover subgraph but not an induced subposet, and (iii) $`𝒮_3^{}`$ is both an induced subposet and a subposet via induced cover subgraph. Let
$`𝒮_4^{}:=`$ $`𝒮_5^{}:=`$ and $`𝒮_6^{}:=`$
be subposets of $`𝒮`$. Then, (iv) $`𝒮_4^{}`$ is a chain but not a path in $`𝒮`$, (v) $`𝒮_5^{}`$ is a chain and an upward path from $`w`$ to $`r`$, and (vi) $`𝒮_6^{}`$ is a path between $`w`$ and $`u`$ but neither an upward path nor a downward path. Let
$`𝒮_7^{}:=`$ $`𝒮_8^{}:=`$ and $`𝒮_9^{}:=`$
be cyclic subposets of $`𝒮`$. Then, (vii) $`𝒮_7^{}`$ is neither an induced subposet nor a subposet via induced cover subgraph, (viii) $`𝒮_8^{}`$ is a subposet via induced cover subgraph but not an induced subposet, and (ix) $`𝒮_9^{}`$ is both an induced subposet and a subposet via induced cover subgraph. We note that if a cyclic subposet is an induced subposet, then it must be a subposet via induced cover subgraph. We will show this (Lemma 2.2) in Section 2.2.
### 2.2 Induced cyclic subposets
For developments later in this paper, a study of cyclic subposets turns out to be crucial, and for this we must also study path subposets. Since the material here is irrelevant until Section 4, the reader may wish to return to the present subsection after reading Section 3.
Let $`𝒮`$ be a poset. A path or a cycle $`𝒱`$ (with ground set $`V`$) in $`𝒮`$ is by definition a subposet of the subposet $`𝒱^{}`$ via induced cover subgraph of $`𝒮`$ on $`V`$, and $`𝒱^{}`$ is in turn a subposet of the induced subposet $`𝒱^{\prime \prime }`$ of $`𝒮`$ on $`V`$. So if this $`𝒱`$ is equal to $`𝒱^{\prime \prime }`$, then $`𝒱=𝒱^{}`$. Thus, we have established
###### Lemma 2.2.
Let $`𝒱`$ be a path or a cycle in $`𝒮`$ and have ground set $`V`$. If $`𝒱`$ is the induced subposet of $`𝒮`$ on $`V`$, then $`𝒱`$ is the subposet via induced cover subgraph of $`𝒮`$ on $`V`$.
An upward (or downward) path is a chain and therefore it is both (i) an induced subposet and (ii) a subposet via induced cover subgraph. As shown in Example 2.1(vi), a path in general may be neither of these; however, we can always devise a path with such properties which substitutes.
###### Lemma 2.3.
Suppose that there exists a path from $`x`$ to $`y`$ in $`𝒮`$. Then there is a path $`𝒱`$ from $`x`$ to $`y`$ in $`𝒮`$ which is an induced subposet of $`𝒮`$.
###### Proof..
Partially ordering the up-sets and (separately) the down-sets in $`𝒮`$ by set inclusion, let $`U_0`$ be a minimal up-set in $`𝒮`$ containing the vertices of some path from $`x`$ to $`y`$. (By assumption, $`S`$ is an up-set containing such a path, so $`U_0`$ exists.) Let $`V_0`$ be a minimal down-set in $`𝒮`$ such that $`U_0V_0`$ contains the vertices of some path from $`x`$ to $`y`$. (Again, $`S`$ is a down-set satisfying this condition, so $`V_0`$ exists.) Let $`𝒲`$ be a path from $`x`$ to $`y`$. We can label minimal and maximal elements of $`𝒲`$ and count the segments of $`𝒲`$ alternating upward and downward as follows: As the path $`𝒲`$ is traversed from $`x`$ to $`y`$, the path traces out either an upward or a downward path from $`z_0=x`$ to $`z_1`$, either a downward or an upward path from $`z_1`$ to $`z_2`$, etc., alternatingly, as illustrated in
$`𝒲=`$
Then we can find a path $`𝒲`$ in $`U_0V_0`$, with ground set $`W`$, from $`x`$ to $`y`$ so that $`m1`$ is as small as possible, that is, a path $`𝒲`$ from $`x`$ to $`y`$ satisfying $`WU_0V_0`$ such that the number $`m`$ of segments alternating upward and downward is smallest among such paths.
We first claim that any two minimal elements in $`𝒲`$ are incomparable in $`𝒮`$ and that any two maximal elements in $`𝒲`$ are incomparable in $`𝒮`$. To see this, suppose that $`z_i>z_j`$ for two minimal elements $`z_i`$ and $`z_j`$ in $`𝒲`$. Then there is some downward path $`(u_0,\mathrm{},u_l)`$ in $`𝒮`$ from $`z_i=u_0`$ to $`u_l=z_j`$. Let $`u_k<u_0`$ be the first element among the $`u_i`$’s with $`i0`$ which belongs to $`W`$. So $`u_1,\mathrm{},u_{k1}W`$. Since $`u_0`$, $`u_k`$ are in $`U_0V_0`$, all the vertices of the downward path $`(u_0,\mathrm{},u_k)`$ are also. Replacing the part of $`𝒲`$ between $`u_0`$ and $`u_k`$ by the path $`(u_0,\mathrm{},u_k)`$ (possibly in reverse order), we get a path from $`x`$ to $`y`$ with a lower value of $`m`$, which is impossible. Thus, any two minimal elements in $`𝒲`$ are incomparable in $`𝒮`$. The same holds for any two maximal elements in $`𝒲`$.
To finish the proof, we claim that the path $`𝒲`$ is the induced subposet of $`𝒮`$ on $`W`$, that is, that no pairs of elements of $`W`$ are comparable in $`𝒮`$, beyond those specified by the poset $`𝒲`$. To see this, suppose that there are elements $`w`$ and $`w^{}`$ of $`W`$ which are comparable in $`𝒮`$ but not in $`𝒲`$. By a replacement scheme similar to that employed in the preceding paragraph, we can get a path $`𝒲^{}`$ (with ground set $`W^{}`$) from $`x`$ to $`y`$ in $`U_0V_0`$ which bypasses some $`z_i`$. If we define $`V_0^{}:=W^{}`$ to be the down-set in $`𝒮`$ generated by $`W^{}`$ and (dually) $`U_0^{}`$ to be the up-set
$$U_0^{}:=\{\xi S:\eta \xi \text{ for some }\eta W^{}\},$$
then $`U_0^{}U_0`$ and $`V_0^{}V_0`$, and $`W^{}U_0^{}V_0^{}`$. If $`z_i`$ is minimal in $`𝒲`$, then (by the preceding paragraph) $`z_i`$ is incomparable in $`𝒮`$ with any minimal element of $`𝒲^{}`$, and therefore $`z_iU_0^{}`$. Similarly, if $`z_i`$ is maximal in $`𝒲`$, then $`z_iV_0^{}`$. Thus, we have contradicted the minimality of $`U_0`$ or $`V_0`$ according as $`z_i`$ is minimal or maximal in $`𝒲`$, which completes the proof.
The next three lemmas make it possible to construct an induced cyclic subposet with specific properties when a poset $`𝒮`$ has a cycle. The first of these lemmas is a simple corollary to Lemma 2.3 which ensures that any non-acyclic poset has an induced cyclic subposet.
###### Lemma 2.4.
Suppose that a poset $`𝒮`$ has a cycle $`(x_0,x_1,\mathrm{},x_{n1},x_0)`$. Then $`𝒮`$ has an induced cyclic subposet $`(y_0,y_1,\mathrm{},y_{m1},y_0)`$ such that $`x_0=y_0`$ and $`x_{n1}=y_{m1}`$.
###### Proof..
Let $`𝒮^{}`$ be the poset on ground set $`S`$ obtained by deleting the edge joining $`x_{n1}`$ and $`x_0`$ from the Hasse diagram of $`𝒮`$. Since $`(x_0,x_1,\mathrm{},x_{n1})`$ is a path from $`x_0`$ to $`x_{n1}`$ in $`𝒮^{}`$, by Lemma 2.3 there is a path $`(y_0,y_1,\mathrm{},y_{m1})`$ from $`y_0=x_0`$ to $`y_{m1}=x_{n1}`$ which is an induced subposet of $`𝒮^{}`$. Then, the cycle $`(y_0,y_1,\mathrm{},y_{m1},y_0)`$ is as desired.
###### Lemma 2.5.
If a poset $`𝒮`$ has a pair $`(x,y)`$ of elements such that there exist at least two unequal upward paths from $`x`$ to $`y`$ in $`𝒮`$, then $`𝒮`$ has an induced cyclic subposet which is a subdivided diamond.
###### Proof..
Let $`(u_0,\mathrm{},u_k)`$ and $`(v_0,\mathrm{},v_l)`$ be two unequal upward paths from $`x`$ to $`y`$ in $`𝒮`$. Without loss of generality, the two paths differ in their first step, i.e., $`u_1v_1`$. Clearly $`u_1`$ and $`v_1`$ are incomparable. Let $`U:=\{\xi S:u_1<\xi \text{ and }v_1<\xi \}`$. Then $`U`$ is nonempty because $`yU`$. Let $`y^{}`$ be a minimal (in $`𝒮`$) element of $`U`$. Let $`W_1`$ and $`W_2`$ be upward paths from $`u_1`$ to $`y^{}`$ and from $`v_1`$ to $`y^{}`$, respectively. Then, for any $`\xi W_1\{y^{}\}`$ and $`\eta W_2\{y^{}\}`$, $`\xi `$ and $`\eta `$ are incomparable; otherwise, the minimality of $`y^{}`$ is contradicted. Thus, the ground set $`\{x\}W_1W_2`$ gives the desired induced subposet.
###### Lemma 2.6.
If a poset $`𝒮`$ has a cycle with height at least $`3`$, then $`𝒮`$ has an induced cyclic subposet $`𝒱`$ with height at least $`3`$.
###### Proof..
If $`𝒮`$ has a diamond as an induced subposet, then by Lemma 2.5 we can find a subdivided diamond $`𝒱`$ which is an induced cyclic subposet of $`𝒮`$. Clearly $`𝒱`$ has height at least $`3`$.
Suppose now that $`𝒮`$ has no diamond as an induced subposet. Let $`(x_0,x_1,\mathrm{},x_{n1},x_0)`$ be a cycle with height at least $`3`$ such that $`x_1<x_0<x_{n1}`$. Let $`S^{}:=S\{x_0\}`$ and let $`𝒮^{}`$ with ground set $`S^{}`$ be the subposet of $`𝒮`$ via induced cover graph. So $`(x_1,\mathrm{},x_{n1})`$ is a path in $`𝒮^{}`$. By Lemma 2.3, there is a path $`𝒰=(u_0,u_1,\mathrm{},u_{k1})`$ in $`𝒮^{}`$ from $`u_0=x_1`$ to $`u_{k1}=x_{n1}`$ which is an induced subposet of $`𝒮^{}`$. If $`𝒰`$ is an upward path from $`x_1`$ to $`x_{n1}`$, then $`𝒮`$ has two distinct upward paths from $`x_1`$ to $`x_{n1}`$ which, by Lemma 2.5, contradicts our assumption. Thus, $`𝒰`$ is not an upward path and in particular $`x_1`$ and $`x_{n1}`$ are incomparable in $`𝒰`$ but comparable in $`𝒮`$. This implies that $`𝒰`$ is not an induced subposet of $`𝒮`$. Let $`𝒰^{}=(u_i,u_{i+1},\mathrm{},u_i^{})`$ be a minimal segment of the path $`𝒰`$ which is not an induced subposet of $`𝒮`$, that is, a segment $`𝒰^{}`$ such that (i) $`𝒰^{}`$ is not an induced subposet of $`𝒮`$, and (ii) any proper segment of $`𝒰^{}`$ is an induced subposet of $`𝒮`$. Then $`u_i`$ and $`u_i^{}`$ must be incomparable in $`𝒰`$ but comparable in $`𝒮`$. Without loss of generality, we may assume that $`u_i<u_i^{}`$ in $`𝒮`$. Then there is a downward path $`𝒲=(w_0,w_1,\mathrm{},w_{l1})`$ in $`𝒮`$ from $`w_0=u_i^{}`$ to $`w_{l1}=u_i`$.
Let $`𝒱:=(u_i,u_{i+1},\mathrm{},u_i^{},w_1,\mathrm{},w_{l1})`$ be a cycle in $`𝒮`$. Since $`𝒰`$ is an induced subposet of $`𝒮^{}`$, we must have $`x_0W`$ and in particular $`𝒱`$ has height at least $`3`$. We claim that $`𝒱`$ is an induced subposet of $`𝒮`$; establishing the claim will complete the proof of the lemma. To prove the claim, suppose that $`u_j`$ and $`w_j^{}`$ are comparable in $`𝒮`$ but incomparable in $`𝒱`$ for some $`u_jU^{}\{u_i,u_i^{}\}`$ and $`w_j^{}W`$. Then $`u_j`$ is comparable with $`u_i`$ or $`u_i^{}`$ in $`𝒮`$ according as $`w_j^{}<u_j`$ or $`w_j^{}>u_j`$. If the pair (either $`u_j,u_i`$ or $`u_j,u_i^{}`$) are comparable in $`𝒱`$, then there are two unequal upward paths (the one through $`w_j^{}`$ and the other consisting of a segment of $`𝒱`$) with common ends in $`𝒮`$. By Lemma 2.5, this contradicts our diamond-free assumption. If the pair is incomparable in $`𝒱`$, then some proper segment of $`𝒰^{}`$ is not an induced subposet of $`𝒮`$, and this contradicts the minimality of $`𝒰^{}`$.
Lemma 2.4 implies that if a poset $`𝒮`$ is non-acyclic then it has an induced cyclic subposet. The next result gives various sufficient conditions for $`𝒮`$ to be non-acyclic.
###### Proposition 2.7.
If a poset $`𝒮`$ has an induced subposet $`𝒮^{}`$ poset-isomorphic to any of the following posets, then $`𝒮`$ is non-acyclic:
* the diamond;
* a subdivided crown with height at least $`3`$;
* the $`k`$-crown for some $`k3`$;
* the following “double-bowtie” poset:
(2.1)
###### Proof..
(i) If $`𝒮^{}`$ is the diamond as in (1.8), then there are at least two unequal upward paths from $`x`$ to $`w`$. By Lemma 2.5, $`𝒮`$ has a cycle (namely, a subdivided diamond).
(ii) Suppose that $`𝒮^{}`$ is a subdivision of the $`k`$-crown displayed and labeled in Figure 2.1(d). Since by assumption $`𝒮^{}`$ has height at least $`3`$, without loss of generality we may assume that there exists $`z^{}S^{}`$ such that $`x_0<z^{}<y_0`$. Then we can find an upward path $`(x_0,\mathrm{},z,z^{},\mathrm{},y_0)`$ in $`𝒮`$ from $`x_0`$ to $`y_0`$ with height at least $`3`$. Since $`z^{}`$ is incomparable in $`𝒮`$ with each of $`x_1,\mathrm{},x_{k1}`$ and each of $`y_1,\mathrm{},y_{k1}`$, no upward path in $`𝒮`$ from any $`x_i`$ to any $`y_j`$ contains the directed edge $`(z,z^{})`$ unless $`(i,j)=(0,0)`$. Let $`(S,_𝒮)`$ be the cover graph of $`𝒮`$. Then $`x_0`$ and $`z^{}`$ are connected in the graph $`(S,_𝒮\{\{z,z^{}\}\})`$, which implies that $`𝒮`$ is non-acyclic.
(iii) Suppose that $`𝒮^{}`$ is the $`k`$-crown for some $`k3`$, as displayed and labeled in Figure 2.1(d). Since the set $`B:=\{\xi S:x_0\xi y_0\text{ and }x_0\xi y_{k1}\}`$ is nonempty, we can find a maximal element $`x_0^{}`$ in the set $`B`$. Then $`x_0^{}`$ is incomparable with $`x_1,\mathrm{},x_{k1},y_1,\mathrm{},y_{k2}`$ and therefore the subposet of $`𝒮`$ induced by the ground set $`\{x_0^{},x_1,\mathrm{},x_{k1},y_0,\mathrm{},y_{k1}\}`$ is again a $`k`$-crown. Thus, we may without loss of generality assume that the $`k`$-crown has no element $`\xi `$ satisfying $`x_0<\xi y_0`$ and $`x_0<\xi y_{k1}`$. Let $`(x_0,z,\mathrm{},y_0)`$ (with $`z=y_0`$ possible) be an upward path from $`x_0`$ to $`y_0`$. By our assumption, an upward path from $`x_0`$ to $`y_{k1}`$ does not contain the directed edge $`(x_0,z)`$. Furthermore, no upward path from any $`x_i`$ to any $`y_j`$ \[except when $`(i,j)=(0,0)`$\] contains the directed edge $`(x_0,z)`$ either; otherwise, $`𝒮^{}`$ is not an induced subposet of $`𝒮`$. Thus we see that $`x_0`$ and $`z`$ are connected in the graph $`(S,_𝒮\{\{x_0,z\}\})`$, which implies that $`𝒮`$ is non-acyclic.
(iv) Suppose that $`𝒮^{}`$ is the poset as in (2.1). Since the set $`B:=\{\xi S:x_1\xi y_1\text{ and }x_1\xi y_2\}`$ is nonempty, we can find a maximal element $`x_1^{}`$ in the set $`B`$. If $`x_1^{}`$ and $`x_2`$ are comparable, then we have $`x_2<x_1^{}`$. Noticing that $`x_1^{}`$ is incomparable with $`x_3`$ and $`y_3`$, the cycle $`(x_2,x_1^{},y_2,x_3,y_3,x_2)`$ is an induced subposet of $`𝒮`$ which is a subdivided $`2`$-crown with height $`3`$. Then (ii) implies that $`𝒮`$ is non-acyclic. So we may assume that $`x_1^{}`$ and $`x_2`$ are incomparable. But then an upward path $`(x_1^{},z,\mathrm{},y_1)`$ in $`𝒮`$ from $`x_1^{}`$ to $`y_1`$ (with $`z=y_1`$ possible) does not share the directed edge $`(x_1^{},z)`$ with any upward path from $`x_2`$ to either $`y_1`$ or $`y_2`$. Moreover, an upward path from $`x_1^{}`$ to $`y_2`$ does not contain the directed edge $`(x_1^{},z)`$; otherwise, maximality of $`x_1^{}`$ is contradicted. Thus, we see that $`x_1^{}`$ and $`z`$ are connected in the graph $`(S,_𝒮\{\{x_1^{},z\}\})`$, which implies that $`𝒮`$ is non-acyclic.
### 2.3 Extensibility and Strassen’s theorem
Strassen’s pioneering work on the existence of probability measures with specified marginals has been influential for the development of the theory and applications of stochastic ordering (e.g., ). We will treat briefly the general subject of probability measures with specified marginals and review some results essential for our later development. Since we restrict attention to finite sets in the present paper, some of the results presented here are greatly simplified by our not needing to deal with topological and other technical matters. (For an interesting review of the subject matter in a general topological setting, see .)
Let $`A`$ and $`S`$ be finite sets and let $`S^A`$ be the collection of all functions $`𝐱=(x_\alpha :\alpha A)`$ from $`A`$ into $`S`$. For $`𝐱S^A`$ and $`\alpha A`$, $`\pi _\alpha (𝐱)`$ will denote the $`\alpha `$-coordinate of $`𝐱`$. Let $`\alpha A`$ be fixed. Then $`\pi _\alpha `$, the $`\alpha `$-projection from $`S^A`$ to the $`\alpha `$-coordinate space $`S`$, is a surjective map from $`S^A`$ to $`S`$. Given a probability measure $`Q`$ on $`S^A`$, we define the probability measure $`Q\pi _\alpha ^1`$ on $`S`$ in the usual way via $`(Q\pi _\alpha ^1)(B):=Q(\pi _\alpha ^1(B))`$ for any subset $`B`$ of $`S`$.
Consider the set of all signed measures on $`S^A`$ as a normed vector space equipped with a suitable topology. Strassen established the following theorem.
###### Theorem 2.8.
(Strassen ) Let $`\mathrm{\Lambda }`$ be a nonempty convex closed subset of probability measures on $`S^A`$ and let $`(P_\alpha :\alpha A)`$ be a system of probability measures on $`S`$. Then there exists a probability measure $`Q\mathrm{\Lambda }`$ such that
(2.2)
$$Q\pi _\alpha ^1=P_\alpha \text{for every }\alpha A$$
if and only if
(2.3)
$$\underset{\alpha A}{}\left(\underset{\xi S}{}P_\alpha (\{\xi \})f_\alpha (\xi )\right)sup\{\underset{𝐱S^A}{}Q(\{𝐱\})\left(\underset{\alpha A}{}f_\alpha \pi _\alpha \right)(𝐱):Q\mathrm{\Lambda }\}$$
for any system $`(f_\alpha :\alpha A)`$ of real-valued functionals on $`S`$.
Let $`\mathrm{\Delta }`$ be a nonempty subset of $`S^A`$. Then we say that a system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ is extensible on $`\mathrm{\Delta }`$ if there exists a probability measure $`Q`$ on $`S^A`$ satisfying (2.2) and
(2.4)
$$Q(\mathrm{\Delta })=1.$$
Let $`\mathrm{\Lambda }_\mathrm{\Delta }`$ be the set of all probability measures on $`S^A`$ satisfying (2.4). Clearly, $`\mathrm{\Lambda }_\mathrm{\Delta }`$ is nonempty, closed, and convex, so Theorem 2.8 applies to it. Observe that $`\mathrm{\Lambda }_\mathrm{\Delta }`$ is the convex hull of the set $`\{\delta _𝐱:𝐱\mathrm{\Delta }\}`$ where $`\delta _𝐱`$ denotes the point-mass probability at $`𝐱`$. Thus, the following theorem is a special case of Theorem 2.8.
###### Theorem 2.9.
Let $`\mathrm{\Delta }`$ be a nonempty subset of $`S^A`$ and let $`(P_\alpha :\alpha A)`$ be a system of probability measures on $`S`$. Then $`(P_\alpha :\alpha A)`$ is extensible on $`\mathrm{\Delta }`$ if and only if
(2.5)
$$\underset{\alpha A}{}\left(\underset{\xi S}{}P_\alpha (\{\xi \})f_\alpha (\xi )\right)sup\{\underset{\alpha A}{}f_\alpha \pi _\alpha (𝐱):𝐱\mathrm{\Delta }\}$$
for any system $`(f_\alpha :\alpha A)`$ of real-valued functionals on $`S`$.
###### Remark 2.10.
Throughout this paper, we use the term “system” in place of “family” to refer to a collection of probability measures, random variables, or real-valued functionals. When a partial ordering on the index set $`A`$ is introduced in the later discussion, the usage becomes more appropriate. We have co-opted the term “extensibility” for a system of probability measures from a use by Vorob’ev in a slightly different setting. Vorob’ev’s “extensibility” problem is now generally called the marginal problem in the literature (e.g. ).
### 2.4 The monotonicity equivalence problem
Since realizable monotonicity always implies stochastic monotonicity, the monotonicity equivalence problem for a given pair $`(𝒜,𝒮)`$ of posets is to either verify or disprove Statement 2.11.
###### Statement 2.11.
For the given pair $`(𝒜,𝒮)`$ of posets, every stochastically monotone system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ is realizably monotone.
We first formulate the monotonicity equivalence problem as a special case of the extensibility problem of Section 2.3. Let $`𝒜`$ and $`𝒮`$ be finite posets. We say that an element $`𝐱=(x_\alpha :\alpha A)`$ of $`S^A`$ is monotone if $`x_\alpha x_\beta `$ in $`𝒮`$ whenever $`\alpha \beta `$ in $`𝒜`$. Define $`\mathrm{\Delta }`$ to be the collection of all monotone elements of $`S^A`$. Given a stochastically monotone system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$, we say that a probability measure $`Q`$ on $`S^A`$ realizes the monotonicity if it satisfies (2.2) and (2.4). Observe that a system $`(𝐗_\alpha :\alpha A)`$ of $`S`$-valued random variables is merely an $`S^A`$-valued random variable distributed as a probability measure $`Q`$ on $`S^A`$. Clearly, the existence of $`(𝐗_\alpha :\alpha A)`$ satisfying (1.5)–(1.6) is equivalent to the existence of a probability measure $`Q`$ on $`S^A`$ satisfying (2.2) and (2.4). Thus, $`(P_\alpha :\alpha A)`$ is realizably monotone if and only if it is extensible on $`\mathrm{\Delta }`$. This formulation establishes Theorem 2.9 as a necessary and sufficient condition for realizable monotonicity.
We now present, without proof, some first simple results on the monotonicity equivalence problem. The upshot of these results is that we need only consider connected posets in our investigation of monotonicity equivalence.
###### Lemma 2.12.
Suppose that $`𝒜^{}`$ is a (not necessarily induced) subposet of $`𝒜`$. If $`(P_\alpha :\alpha A)`$ is realizably monotone, then so is $`(P_\alpha :\alpha A^{})`$.
###### Lemma 2.13.
Suppose that $`𝒮^{}`$ is an induced subposet of $`𝒮`$. If monotonicity equivalence holds for $`(𝒜,𝒮)`$, then it holds for $`(𝒜,𝒮^{})`$.
###### Proposition 2.14.
Suppose that $`𝒜`$ is the disjoint union of nonempty posets $`𝒜_1`$ and $`𝒜_2`$. Then monotonicity equivalence holds for $`(𝒜,𝒮)`$ if and only if it holds for both $`(𝒜_1,𝒮)`$ and $`(𝒜_2,𝒮)`$.
###### Proposition 2.15.
Suppose that $`𝒮`$ is the disjoint union of nonempty posets $`𝒮_1`$ and $`𝒮_2`$. Then monotonicity equivalence holds for $`(𝒜,𝒮)`$ if and only if it holds for both $`(𝒜,𝒮_1)`$ and $`(𝒜,𝒮_2)`$.
The next proposition is an immediate consequence of the observation that the collection $`\mathrm{\Delta }`$ of all monotone elements of $`S^A`$ for $`(𝒜,𝒮)`$ is equal to the corresponding collection for $`(𝒜^{},𝒮^{})`$.
###### Proposition 2.16.
Monotonicity equivalence holds for $`(𝒜,𝒮)`$ if and only if it holds for $`(𝒜^{},𝒮^{})`$.
## 3 Subclasses of connected posets
As explained by Propositions 2.14 and 2.15, we assume without further notice that $`𝒜`$ and $`𝒮`$ are connected posets throughout the remainder of the paper. We partition the collection of connected posets $`𝒮`$ into the following four subclasses. We say that
* $`𝒮`$ is in Class $`\mathrm{B}`$, denoted $`𝒮\mathrm{B}`$, if $`𝒮`$ has either a cycle or an induced bowtie;
* $`𝒮`$ is in Class $`\mathrm{Y}`$, denoted $`𝒮\mathrm{Y}`$, if (i) $`𝒮\mathrm{B}`$, and (ii) $`𝒮`$ has an induced Y-poset;
* $`𝒮`$ is in Class $`\mathrm{W}`$, denoted $`𝒮\mathrm{W}`$, if (i) $`𝒮\mathrm{B}\mathrm{Y}`$, and (ii) $`𝒮`$ has an induced W-poset;
* $`𝒮`$ is in Class $`\mathrm{Z}`$, denoted $`𝒮\mathrm{Z}`$, if $`𝒮\mathrm{B}\mathrm{Y}\mathrm{W}`$.
Note that a poset $`𝒮`$ in Class B may be acyclic. For example, let
$`𝒮:=`$
be an acyclic poset. Then the subposet of $`𝒮`$ induced by $`\{x_1,x_2,z_1,z_2\}`$ is the bowtie; thus, $`𝒮\mathrm{B}`$. If the cover graph $`(S,_𝒮)`$ of a given poset $`𝒮`$ has an element $`x`$ whose degree is at least $`3`$, then $`𝒮`$ must have either a Y-poset or a W-poset as an induced subposet and therefore $`𝒮\mathrm{B}\mathrm{Y}\mathrm{W}`$. This implies that Class Z consists precisely of those posets $`𝒮`$ whose cover graph $`(S,_𝒮)`$ is a path (and the nature of whose Hasse diagram is therefore “zig-zag,” which explains our choice of “Z”).
Given a poset $`𝒮`$, we call $`𝒜`$ a poset of monotonicity equivalence or of monotonicity inequivalence (for $`𝒮`$) according as Statement 2.11 is true or false for the pair $`(𝒜,𝒮)`$. The question of monotonicity equivalence raised in Section 1 can be recast as that of determining, for each $`𝒮`$, the class $`(𝒮)`$ of all posets $`𝒜`$ of monotonicity equivalence for $`𝒮`$. For a poset $`𝒮`$ in Class B, Y, or Z, we can characterize the class $`(𝒮)`$ precisely. Furthermore, the class $`(𝒮)`$ is the same for every $`𝒮`$ of the same class among Classes B, Y, and Z. In the rest of this paper, we will show that
* for every $`𝒮B`$, $`(𝒮)`$ is the collection of all acyclic posets $`𝒜`$ (Theorem 4.2);
* for every $`𝒮Y`$, $`(𝒮)`$ is the collection of posets $`𝒜`$ such that $`𝒜`$ is enlargeable to an acyclic poset (Theorem 5.1);
* for every $`𝒮Z`$, $`(𝒮)`$ is the class of all posets $`𝒜`$ (Theorem 6.1).
For a poset $`𝒮`$ of Class W, we can exhibit a large subclass of $`(𝒮)`$. But the assertion that the class $`(𝒮)`$ is the same for every $`𝒮`$ of Class W is false. Our investigation of Class W is presented in the companion paper .
## 4 The monotonicity equivalence problem on Class B
In this section, we solve the monotonicity equivalence problem when $`𝒮`$ is a poset of Class B. The main results of this section are summarized in the following two theorems.
###### Theorem 4.1.
If $`𝒜`$ is an acyclic poset, then monotonicity equivalence holds for $`(𝒜,𝒮)`$ for any $`𝒮`$.
###### Theorem 4.2.
Let $`𝒮`$ be a poset of Class B. Then monotonicity equivalence holds for $`(𝒜,𝒮)`$ if and only if $`𝒜`$ is an acyclic poset.
In Section 4.1 we briefly review a well-known characterization of stochastic ordering and then prove Theorem 4.1. Theorem 4.1 establishes a sufficient condition for a poset $`𝒜`$ of monotonicity equivalence which is applicable to any poset $`𝒮`$. But further generalization is not possible when $`𝒮`$ is a poset of Class B. In Section 4.2 we present various counterexamples where monotonicity equivalence fails for a pair $`(𝒜,𝒮)`$ of non-acyclic posets. In Section 4.3 we build on these counterexamples to complete the proof of Theorem 4.2. Theorems 4.1 and 4.2 can be immediately combined to settle the monotonicity equivalence question for Markov transition matrices, where $`𝒜=𝒮`$ (cf. the end of Section 1):
###### Theorem 4.3.
If $`𝒜=𝒮`$, then monotonicity equivalence holds for $`(𝒜,𝒮)`$ if and only if $`𝒮`$ is an acyclic poset.
### 4.1 Stochastic ordering and acyclic index posets $`𝒜`$
To supplement the characterization of stochastic ordering described in Section 1, Kamae, Krengel, and O’Brien introduced an equivalent condition in terms of upward kernels. Using their condition, they showed that, for a sequence $`(P_1,P_2,\mathrm{})`$ of probability measures on a common poset $`𝒮`$, we have $`P_iP_{i+1}`$ for $`i=1,2,\mathrm{}`$ if and only if there exists a sequence $`(𝐗_1,𝐗_2,\mathrm{})`$ such that $`𝐗_i𝐗_{i+1}`$ and $`(𝐗_i)=P_i()`$ for $`i=1,2,\mathrm{}`$. We will show (Theorem 4.1) that this result can be generalized to a system $`(P_\alpha :\alpha A)`$ of probability measures whenever $`𝒜`$ is an acyclic poset.
Let $`𝒮`$ be a poset. A function $`k`$ from $`S\times 2^S`$ to $`[0,1]`$ is called a stochastic kernel on $`S`$ if $`k(x,)`$ is a probability measure on $`S`$ for every $`xS`$. A stochastic kernel $`k`$ on $`𝒮`$ is said to be upward if
(4.1)
$$k(x,\{\xi S:x\xi \text{ in }𝒮\})=1\text{ for each }xS\text{}$$
We collect several characterizations of stochastic ordering in the following proposition.
###### Proposition 4.4.
(Kamae, Krengel, and O’Brien ) Let $`(P_1,P_2)`$ be a pair of probability measures on $`S`$. Then the following conditions are equivalent:
* $`P_1P_2`$;
* $`P_1(V)P_2(V)`$ for every down-set $`V`$ in $`𝒮`$;
* there exists a pair $`(𝐗_1,𝐗_2)`$ of $`S`$-valued random variables satisfying (1.3)(1.4);
* there exists an upward kernel $`k`$ such that
(4.2)
$$P_2()=\underset{xS}{}P_1(\{x\})k(x,).$$
Now consider the monotonicity equivalence problem. The equivalence of (a) and (c) in Proposition 4.4 can be extended to equivalence between stochastic monotonicity and realizable monotonicity in the case that $`𝒜`$ is acyclic. The precise result has already been stated as Theorem 4.1.
###### Proof of Theorem 4.1..
We prove the claim of Theorem 4.1 by induction over the cardinality of $`A`$. The claim is vacuous when $`|A|=1`$. We now suppose that the claim is true for an acyclic poset $`𝒜^{}`$ when $`|A^{}|=n1`$ for fixed $`n2`$, and consider an acyclic poset $`𝒜`$ with cardinality $`|A|=n`$. Let $`a`$ be a leaf in $`𝒜`$. Without loss of generality, we assume that $`a`$ is maximal in $`𝒜`$; thus, there is a unique element $`b`$ which is covered by $`a`$. We consider the subposet $`𝒜^{}`$ of $`𝒜`$ induced by the ground set $`A^{}=A\{a\}`$.
Let $`(P_\alpha :\alpha A)`$ be a stochastically monotone system of probability measures on $`S`$. Then the subsystem $`(P_\alpha :\alpha A^{})`$ is stochastically monotone. Since $`𝒜^{}`$ is an acyclic poset and $`|A^{}|=n1`$, by the induction hypothesis there exists a probability measure $`Q^{}`$ on $`S^A^{}`$ which realizes the monotonicity. Since $`P_bP_a`$, by Proposition 4.4 there exists an upward kernel $`k`$ satisfying (4.2) for the pair $`(P_b,P_a)`$ of probability measures. We can define a probability measure $`Q`$ on $`S^A`$ by
$$Q(\{𝐱\}):=Q^{}(\{\pi _A^{}(𝐱)\})k(\pi _b(𝐱),\{\pi _a(𝐱)\})\text{ for }𝐱S^A\text{,}$$
where $`\pi _A^{}`$ denotes the projection from $`S^A`$ to $`S^A^{}`$ and $`\pi _\alpha `$ denotes the $`\alpha `$-projection from $`S^A`$ to $`S`$ for each $`\alpha A`$. In words, this says simply that we couple together the probability measures $`P_\alpha `$, $`\alpha A^{}`$ using $`Q^{}`$ and then extend the multivariate coupling to $`P_a`$ using the upward kernel $`k`$ from $`P_b`$ to $`P_a`$. Observe that $`Q^{}`$ couples the probability measures $`(P_\alpha `$: $`\alpha A^{})`$ correctly, and that $`a>\alpha A^{}`$ in $`𝒜`$ implies $`b\alpha `$ in $`𝒜^{}`$. So coupling $`P_a`$ to $`P_b`$ correctly automatically couples $`P_a`$ to each $`P_\alpha `$ ($`\alpha A^{}`$) correctly. Thus, $`Q`$ realizes the monotonicity of $`(P_\alpha :\alpha A)`$, and therefore the claim holds for $`𝒜`$.
### 4.2 Monotonicity inequivalence on Class B
The objective of this subsection is to present several examples of monotonicity inequivalence. To establish such an example we must exhibit a pair $`(𝒜,𝒮)`$ of posets and a specific system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ which is stochastically but not realizably monotone. We have already done this when both $`𝒜`$ and $`𝒮`$ are diamonds: see Example 1.1. Our simple examples, including Example 1.1, will serve as building blocks for more complex counterexamples that establish quite general negative results.
###### Example 4.5.
Let
$`𝒜:=`$
be the bowtie and let $`𝒮`$ be the diamond as in (1.8). Define a system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ by
(4.3)
$$P_\alpha :=\{\begin{array}{cc}\text{unif}\{x,w\}\hfill & \text{ if }\alpha =a_0\text{}\hfill \\ \text{unif}\{y,z\}\hfill & \text{ if }\alpha =a_1\text{}\hfill \\ \text{unif}\{y,w\}\hfill & \text{ if }\alpha =b_0\text{}\hfill \\ \text{unif}\{z,w\}\hfill & \text{ if }\alpha =b_1\text{}\hfill \end{array}$$
The system is clearly stochastically monotone.
To see that it is not realizably monotone, suppose that there exists a system $`(𝐗_\alpha :\alpha A)`$ of $`S`$-valued random variables which realizes the monotonicity. Considering the event $`\{𝐗_{b_0}=y\}`$, in order to maintain monotonicity we must have
$$(𝐗_{b_0}=y)=(𝐗_{a_0}=x,𝐗_{a_1}=y,𝐗_{b_0}=y,𝐗_{b_1}=w)=\frac{1}{2}.$$
Similarly, we must have
$$(𝐗_{a_0}=w)=(𝐗_{a_0}=w,𝐗_{b_0}=w,𝐗_{b_1}=w)=\frac{1}{2}.$$
Since the above two events are disjoint, we must have $`(𝐗_{b_1}=w)=1`$, which contradicts $`(𝐗_{b_1}=w)=\frac{1}{2}`$.
###### Example 4.6.
Let $`𝒜`$ be the bowtie and let $`𝒮`$ be a $`k`$-crown with $`k2`$. The posets $`𝒜`$ and $`𝒮`$ are displayed and labeled in Example 4.5 and Figure 2.1(d), respectively. Define a system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ by
(4.4)
$$P_\alpha :=\{\begin{array}{cc}\frac{k1}{k}\text{unif}(S_1\{x_1\})+\frac{1}{k}\text{unif}(\{y_0,x_1\})\hfill & \text{ if }\alpha =a_0\text{}\hfill \\ \frac{1}{k}\text{unif}(\{x_0\})+\frac{k1}{k}\text{unif}(S\{x_0,y_0\})\hfill & \text{ if }\alpha =a_1\text{}\hfill \\ \frac{1}{k}\text{unif}(\{y_{k1}\})+\frac{k1}{k}\text{unif}(S\{x_0,y_{k1}\})\hfill & \text{ if }\alpha =b_0\text{}\hfill \\ \frac{k1}{k}\text{unif}(\{x_0\}(S_2\{y_0,y_{k1}\}))+\frac{1}{k}\text{unif}(\{y_0,y_{k1}\})\hfill & \text{ if }\alpha =b_1\text{}\hfill \end{array}$$
where
$$S_1:=\{x_0,x_1,\mathrm{},x_{k1}\}\text{ and }S_2:=\{y_0,y_1,\mathrm{},y_{k1}\}.$$
Then $`(P_\alpha :\alpha A)`$ is stochastically monotone.
Now let $`\mathrm{\Delta }`$ be the collection of all monotone elements of $`S^A`$ and let
$$U_\alpha :=\{\begin{array}{cc}\{y_0\}\hfill & \text{ if }\alpha =a_0\text{}\hfill \\ S_2\{y_0\}\hfill & \text{ if }\alpha =a_1\text{}\hfill \\ S_1\{x_0\}\hfill & \text{ if }\alpha =b_0\text{}\hfill \\ \{x_0\}\hfill & \text{ if }\alpha =b_1\text{}\hfill \end{array}$$
This builds a system $`(I_{U_\alpha }:\alpha A)`$ of real-valued functions on $`S`$, where $`I_{U_\alpha }`$ denotes the indicator function of a subset $`U_\alpha `$ of $`S`$. It is not hard to verify that
$$\underset{\alpha A}{}\left(I_{U_\alpha }\pi _\alpha \right)(𝐱)1\text{ for any }𝐱\mathrm{\Delta }\text{}$$
Since
$$\underset{\alpha A}{}P_\alpha (U_\alpha )=1+\frac{1}{2k},$$
by Theorem 2.9 we have shown that $`(P_\alpha :\alpha A)`$ is not realizably monotone.
###### Remark 4.7.
The specific systems $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ presented in Examples 4.5 and 4.6 will be used in our later discussions. We further define probability measures $`P_{\widehat{0}}`$ and $`P_{\widehat{1}}`$ on $`S`$ for each example.
* In Example 4.5, let $`P_{\widehat{0}}:=\delta _x`$ and $`P_{\widehat{1}}:=\delta _w`$. Clearly, $`P_{\widehat{0}}PP_{\widehat{1}}`$ for any probability measure $`P`$ on $`S`$.
* In Example 4.6, let $`P_{\widehat{0}}:=\text{unif}(S_1)`$ and $`P_{\widehat{1}}:=\text{unif}(S_2)`$. Then we have $`P_{\widehat{0}}P_\alpha P_{\widehat{1}}`$ for every $`P_\alpha `$ defined at (4.4).
###### Example 4.8.
Let
(4.5)
$$𝒜:=\text{}$$
be the diamond and let $`𝒮`$ be a $`k`$-crown for $`k2`$ as in Figure 2.1(d). Define a system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ by
$$P_\alpha :=\{\begin{array}{cc}\text{unif}(S_1)\hfill & \text{ if }\alpha =a\text{}\hfill \\ \text{unif}(\{y_0\}(S_1\{x_0\}))\hfill & \text{ if }\alpha =b\text{}\hfill \\ \text{unif}(\{y_{k1}\}(S_1\{x_0\}))\hfill & \text{ if }\alpha =c\text{}\hfill \\ \text{unif}(S_2)\hfill & \text{ if }\alpha =d\text{}\hfill \end{array}$$
where $`S_1`$ and $`S_2`$ are defined as in Example 4.6. Then $`(P_\alpha :\alpha A)`$ is stochastically monotone.
Now let $`\mathrm{\Delta }`$ be the collection of all monotone elements of $`S^A`$ and let
$$U_\alpha :=\{\begin{array}{cc}\{x_0\}\hfill & \text{ if }\alpha =a\text{}\hfill \\ \{y_0\}(S_1\{x_0\})\hfill & \text{ if }\alpha =b\text{}\hfill \\ \{y_{k1}\}(S_1\{x_0\})\hfill & \text{ if }\alpha =c\text{}\hfill \\ \mathrm{}\hfill & \text{ if }\alpha =d\text{}\hfill \end{array}$$
Then we have
$$\underset{\alpha A}{}\left(I_{U_\alpha }\pi _\alpha \right)(𝐱)2\text{ for any }𝐱\mathrm{\Delta }\text{}$$
To see this, suppose that the sum is $`3`$ for some monotone element $`𝐱`$. Then we must have $`(\pi _a(𝐱),\pi _b(𝐱),\pi _c(𝐱))=(x_0,y_0,y_{k1})`$, which is impossible. Since
$$\underset{\alpha A}{}P_\alpha (U_\alpha )=2+\frac{1}{k},$$
we deduce from Theorem 2.9 that $`(P_\alpha :\alpha A)`$ is not realizably monotone.
Examples 1.1 and 4.5 both employ a certain probabilistic argument which assumes the existence of certain random variables and leads to a contradiction. Here we introduce a lemma which is useful in conjunction with such probabilistic arguments when we extend monotonicity equivalence beyond our previously considered counterexamples.
###### Lemma 4.9.
Let $`𝒮`$ be a poset and let $`(𝐗_1,𝐗_2)`$ be a pair of $`S`$-valued random variables. If $`(𝐗_1)=(𝐗_2)`$ and $`(𝐗_1𝐗_2)=1`$, then $`(𝐗_1=𝐗_2)=1`$.
###### Proof..
Notice that for any $`\xi S`$,
$$(𝐗_1\xi )=(𝐗_1\xi ,𝐗_1𝐗_2)(𝐗_2\xi )+(𝐗_1=\xi <𝐗_2).$$
Since $`(𝐗_1\xi )=(𝐗_2\xi )`$, we deduce $`(𝐗_1=\xi <𝐗_2)=0`$. Thus we obtain
$$(𝐗_1<𝐗_2)=\underset{\xi S}{}(𝐗_1=\xi <𝐗_2)=0,$$
which completes the proof.
Now let
(4.6)
$$𝒜_k:=\text{}$$
be a $`k`$-crown. If we have a known case of monotonicity inequivalence for a pair $`(𝒜_k,𝒮)`$ of posets, then we can apply Lemma 4.9 to extend monotonicity inequivalence to $`(𝒜_k^{},𝒮)`$ whenever $`k^{}k`$.
###### Proposition 4.10.
Let $`𝒜_k`$ be a $`k`$-crown as in (4.6). Given a pair $`(𝒜_k,𝒮)`$ of posets, suppose that there exists a stochastically monotone system $`(P_\alpha :\alpha A_k)`$ of probability measures on $`S`$ which is not realizably monotone. Then if $`k^{}k`$, we can define
$$P_\alpha :=\{\begin{array}{cc}P_\alpha \hfill & \text{if }\alpha A_k\text{}\hfill \\ P_{b_{k1}}\hfill & \text{if }\alpha A_k^{}A_k\text{}\hfill \end{array}$$
to enlarge $`(P_\alpha :\alpha A_k)`$ to a stochastically monotone system $`(P_\alpha :\alpha A_k^{})`$ which is not realizably monotone for the pair $`(𝒜_k^{},𝒮)`$.
###### Proof..
Since $`P_{b_{k1}}=P_{a_k}=\mathrm{}=P_{a_{k^{}1}}=P_{b_{k^{}1}}`$, we see that the system $`(P_\alpha :\alpha A_k^{})`$ is stochastically monotone. To see that it is not realizably monotone, suppose that there exists a system $`(𝐗_\alpha :\alpha A_k^{})`$ of $`S`$-valued random variables which realizes the monotonicity of $`(P_\alpha :\alpha A_k^{})`$. By applying Lemma 4.9 repeatedly, we (almost surely) have $`𝐗_{b_{k1}}=𝐗_{a_k}=\mathrm{}=𝐗_{a_{k^{}1}}=𝐗_{b_{k^{}1}}`$. But then (after perhaps taking care of null sets) $`(𝐗_\alpha :\alpha A_k)`$ realizes the monotonicity of $`(P_\alpha :\alpha A_k)`$, which is a contradiction.
As an immediate corollary to Proposition 4.10, we can extend Examples 4.5 and 4.6 to allow $`𝒜`$ to be the $`k`$-crown for arbitrary $`k2`$. In summary, from the counterexamples in Example 1.1 and Examples 4.54.8 we have derived
###### Proposition 4.11.
Let $`𝒜`$ and $`𝒮`$ each be either a diamond or a crown. Then monotonicity equivalence fails for $`(𝒜,𝒮)`$.
### 4.3 The proof of Theorem 4.2
Let $`𝒮`$ be a poset of Class B. Then we can find either (i) a $`2`$-crown as an induced subposet of $`𝒮`$ or (ii) a cycle as a (not necessarily induced) subposet of $`𝒮`$. If $`𝒮`$ has a cycle, then, by Lemma 2.4, $`𝒮`$ has an induced cyclic subposet $`𝒱`$. It is possible to label the cycle $`𝒱`$ and to fix a starting point and orientation of the cycle so that, as the cycle is traversed, it traces out an upward path from $`z_0`$ to $`z_1`$, then a downward path from $`z_1`$ to $`z_2`$, then an upward path from $`z_2`$ to $`z_3`$, etc., finishing with a downward path from $`z_{2k1}`$ to $`z_0`$, as illustrated in
(4.7)
$$𝒱=\text{}$$
If $`k=1`$, then $`𝒱`$ is a subdivided diamond; otherwise, $`𝒱`$ is a subdivided $`k`$-crown ($`k2`$). This observation gives a different characterization of Class B.
###### Lemma 4.12.
A poset $`𝒮`$ is in Class B if and only if $`𝒮`$ has either the diamond or a crown as an induced subposet.
###### Proof..
We have already seen that a poset $`𝒮`$ of Class B has either the diamond or a crown as an induced subposet. If $`𝒮`$ has an induced $`2`$-crown, then $`𝒮`$ is in Class B by definition. If $`𝒮`$ has an induced subposet which is either the diamond or a $`k`$-crown for some $`k3`$, then, by Proposition 2.7, $`𝒮`$ is non-acyclic and therefore $`𝒮`$ is in Class B, again by definition.
Now we turn to the proof of Theorem 4.2. If $`𝒜`$ is an acyclic poset, then, by Theorem 4.1, monotonicity equivalence holds for $`(𝒜,𝒮)`$. Thus, the remaining task is to show that if $`𝒜`$ is a non-acyclic poset, then monotonicity equivalence fails for $`(𝒜,𝒮)`$. By Lemmas 2.13 and 4.12, it suffices to show that monotonicity equivalence fails for $`(𝒜,𝒮^{})`$ whenever $`𝒮^{}`$ is either the diamond or a crown. We complete the
###### Proof of Theorem 4.2..
Let $`𝒜`$ be a non-acyclic poset and $`𝒮^{}`$ be either the diamond or an $`m`$-crown for some $`m2`$. We will construct a stochastically monotone system $`(P_\alpha :\alpha A)`$ of probability measures on $`S^{}`$ which is not realizably monotone, by dividing the construction into two cases.
Case I. Suppose that $`𝒜`$ has a diamond $`𝒜^{}`$ as an induced subposet. Let $`𝒜^{}`$ be labeled as in Example 4.8. By Examples 1.1 and 4.8, there exists a stochastically monotone system $`(P_\alpha :\alpha A^{})`$ of probability measures on $`S^{}`$ which is not realizably monotone. It then suffices by Lemma 2.12 to show that the system can be enlarged to a stochastically monotone system $`(P_\alpha :\alpha A)`$ of probability measures on $`S^{}`$.
For this, define a partition $`A_a,A_b,A_c`$, and $`A_d`$ of $`A`$ by
$$A_\alpha :=\{\begin{array}{cc}A\{\alpha A:\text{ }b\alpha \text{ or }c\alpha \}\hfill & \text{ if }\alpha =a\text{}\hfill \\ \{b\}\hfill & \text{ if }\alpha =b\text{}\hfill \\ \{c\}\hfill & \text{ if }\alpha =c\text{}\hfill \\ \{\alpha A:\text{ }b<\alpha \text{ or }c<\alpha \}\hfill & \text{ if }\alpha =d\text{}\hfill \end{array}$$
Then we can extend $`(P_\alpha :\alpha A^{})`$ to $`(P_\alpha :\alpha A)`$ by putting
$$P_\alpha :=P_\beta ,\alpha A_\beta $$
for each $`\beta A^{}`$. It is routine to check that this extension maintains stochastic monotonicity, that is, that if $`\alpha _1<\alpha _2`$, then $`P_{\alpha _1}P_{\alpha _2}`$. This is true if $`\alpha _1,\alpha _2A_\beta `$ for some $`\beta `$ and also if $`\alpha _1A_a`$ or $`\alpha _2A_d`$. If $`\alpha _1\{b,c\}`$ and $`\alpha _1<\alpha _2`$, then $`\alpha _2A_d`$. If $`\alpha _1A_d`$ and $`\alpha _1<\alpha _2`$, then $`\alpha _2A_d`$. So stochastic monotonicity is clear.
Case II. Suppose that $`𝒜`$ has no diamond as an induced subposet. By Lemma 2.4, $`𝒜`$ has an induced cyclic subposet $`𝒜^{}`$. In the same way as what we did in (4.7), we can label the cycle $`𝒜^{}`$ as illustrated in
$`𝒜^{}=`$
By our Case II assumption, $`𝒜^{}`$ must be a subdivided $`k`$-crown for some $`k2`$. Let $`𝒜^{\prime \prime }`$ be the $`k`$-crown $`(a_0,b_0,a_1,b_1,\mathrm{},a_{k1},b_{k1},a_0)`$. By Example 4.5 and 4.6 (and then further using Proposition 4.10, if necessary), there exists a stochastically monotone system $`(P_\alpha ^{\prime \prime }:\alpha A^{\prime \prime })`$ of probability measures on $`S`$ which is not realizably monotone.
Let $`P_{\widehat{0}}^{\prime \prime }`$ and $`P_{\widehat{1}}^{\prime \prime }`$ be defined as in Remark 4.7 so that $`P_{\widehat{0}}^{\prime \prime }P_\alpha ^{\prime \prime }P_{\widehat{1}}^{\prime \prime }`$ for all $`\alpha A^{\prime \prime }`$. Consider the partition $`\{A_\beta ^{}`$: $`\beta A^{\prime \prime }\}`$ of $`A^{}`$, where
$$A_{a_i}^{}:=\{\alpha A^{}:\text{ }a_i\alpha <b_{i1}\text{ or }a_i\alpha <b_i\}$$
and $`A_{b_i}^{}:=\{b_i\}`$ for $`i=0,\mathrm{},k1`$. By letting
$$A_{\widehat{1}}:=\{\alpha AA^{}:\alpha >\beta \text{ for some }\beta A^{}\}$$
and $`A_{\widehat{0}}:=A(A^{}A_{\widehat{1}})`$, we can define a system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ by
$$P_\alpha :=\{\begin{array}{cc}P_\beta ^{\prime \prime }\hfill & \text{ if }\alpha A_\beta ^{}\text{ for some }\beta A^{\prime \prime }\text{}\hfill \\ P_{\widehat{1}}^{\prime \prime }\hfill & \text{ if }\alpha A_{\widehat{1}}\text{}\hfill \\ P_{\widehat{0}}^{\prime \prime }\hfill & \text{ if }\alpha A_{\widehat{0}}\text{}\hfill \end{array}$$
this system extends $`(P_\alpha ^{\prime \prime }:\alpha A^{\prime \prime })`$.
We claim that $`(P_\alpha :\alpha A)`$ is stochastically monotone. Let $`\alpha _1<\alpha _2`$. If $`\alpha _1A_{\widehat{0}}`$ or $`\alpha _2A_{\widehat{1}}`$, then $`P_{\alpha _1}P_{\alpha _2}`$. This is also trivial if $`\alpha _1,\alpha _2A^{}`$. If $`\alpha _1A^{}`$, then $`\alpha _2A^{}A_{\widehat{1}}`$, so $`P_{\alpha _1}P_{\alpha _2}`$. If $`\alpha _1A_{\widehat{1}}`$, then $`\alpha _2A^{}A_{\widehat{1}}`$. We need only show that it is impossible to have both $`\alpha _1A_{\widehat{1}}`$ and $`\alpha _2A^{}`$. Indeed, then $`\alpha _1A^{}`$, but for some $`\beta A^{}`$ we have $`\beta <\alpha _1<\alpha _2`$. But then there are two distinct upward paths from $`\beta `$ to $`\alpha _2`$ in $`𝒜`$, namely the one using edges in the cover graph $`(A^{},_𝒜^{})`$ and one containing $`\alpha _1A^{}`$. This violates Lemma 2.5, since we are assuming that $`𝒜`$ has no induced diamond. Thus, we have established the claim and, by Lemma 2.12, $`(P_\alpha :\alpha A)`$ cannot be realizably monotone.
## 5 The monotonicity equivalence problem on Class Y
In Section 5 we investigate the monotonicity equivalence problem when $`𝒮\mathrm{Y}`$. The goal of this section is to prove the following theorem. Two reformulations of the necessary and sufficient condition here are given in Proposition 5.11.
###### Theorem 5.1.
Let $`𝒮`$ be a poset of Class $`\mathrm{Y}`$. Then monotonicity equivalence holds for $`(𝒜,𝒮)`$ if and only if there exists an acyclic poset $`\stackrel{~}{𝒜}`$ which has $`𝒜`$ as an induced subposet.
Thus, some posets $`𝒜`$ of monotonicity equivalence may be non-acyclic. As an instructive example, let
(5.1)
$$𝒜:=\text{}$$
be a poset where $`a_i<b_j`$ for all $`i=1,\mathrm{},m`$ and all $`j=1,\mathrm{},n`$. Then $`𝒜`$ is a poset of monotonicity equivalence for any $`𝒮\mathrm{B}`$. To see this without resorting to Theorem 5.1, grant the following proposition for now.
###### Proposition 5.2.
Suppose that $`𝒮\mathrm{B}`$. Let $`P_{a,1},\mathrm{},P_{a,m},P_{b,1},\mathrm{},P_{b,n}`$ be probability measures on $`S`$ satisfying
$$P_{a,i}P_{b,j}\text{for all }i=1,\mathrm{},m\text{ and all }j=1,\mathrm{},n\text{}$$
Then there exists a probability measure $`P_0`$ on $`S`$ such that
$$P_{a,i}P_0P_{b,j}\text{for all }i=1,\mathrm{},m\text{ and all }j=1,\mathrm{},n\text{}$$
Now suppose that $`(P_\alpha :\alpha A)`$ is a stochastically monotone system of probability measures on $`S`$. Proposition 5.2 implies that there is a probability measure $`P_0`$ on $`S`$ such that $`P_{a_i}P_0P_{b_j}`$ for all $`i=1,\mathrm{},m`$ and all $`j=1,\mathrm{},n`$. Define an acyclic poset $`\stackrel{~}{𝒜}`$ on the set $`\stackrel{~}{A}=A\{c\}`$ by means of the Hasse diagram
(5.2)
$$\stackrel{~}{𝒜}:=\text{}$$
Then the poset $`𝒜`$ is an induced subposet of $`\stackrel{~}{𝒜}`$. By letting $`P_c:=P_0`$, we can enlarge the system $`(P_\alpha :\alpha A)`$ to the system $`(P_\alpha :\alpha \stackrel{~}{A})`$, which remains stochastically monotone. By applying Theorem 4.1 and then Lemma 2.12, we see that both $`(P_\alpha :\alpha \stackrel{~}{A})`$ and $`(P_\alpha :\alpha A)`$ are realizably monotone. Thus, we have shown that monotonicity equivalence holds for $`(𝒜,𝒮)`$.
In Section 5.1 we attend to the proof of Proposition 5.2. A large class of posets $`𝒜`$ of monotonicity inequivalence is presented in Section 5.2, leading to the proof of Theorem 5.1 in Section 5.3.
### 5.1 Probability measures on an acyclic poset
The goal of this subsection is to prove Proposition 5.2. We begin this subsection by introducing a natural partial ordering on a connected acyclic graph (i.e., a tree) when one vertex is specified to become a top element, that is, to be made larger than every other vertex. Let $`𝒮`$ be a connected acyclic poset and let $`\tau `$ be a fixed leaf of $`𝒮`$. Declare $`x_\tau y`$ for $`x,yS`$ if and only if the (necessarily existent and unique) path $`(\tau ,\mathrm{},x)`$ from $`\tau `$ to $`x`$ contains the path $`(\tau ,\mathrm{},y)`$ from $`\tau `$ to $`y`$ as a segment. This introduces another partial ordering $`_\tau `$ on the same ground set $`S`$ (see ). We call this new poset $`(S,_\tau )`$ a rooted tree (rooted at $`\tau `$). (Comparison of the poset $`𝒮`$ and a rooted tree is illustrated in Example 5.4.) The element $`\tau `$ is clearly the maximum of the rooted tree $`(S,_\tau )`$ and is called the root. If $`x`$ covers $`y`$ in $`(S,_\tau )`$, then $`y`$ is called a successor of $`x`$, and $`x`$ is called the predecessor of $`y`$.
For each $`xS`$, we define a section of rooted tree by
$$(,x]:=\{\xi S:\xi _\tau x\},$$
that is, the down set in $`(S,_\tau )`$ generated by $`x`$ \[cf. Section 2.1(1)\]. Every section $`(,x]`$ is either a down-set or an up-set in $`𝒮`$, and which of these holds can be determined from the cover relation of $`𝒮`$. We state this as the following lemma.
###### Lemma 5.3.
Let $`𝒮`$ be a connected acyclic poset. For every $`xS`$, $`(,x]`$ is either a down-set or an up-set in $`𝒮`$. If $`x\tau `$, then there is a unique predecessor $`w`$ of $`x`$, and the edge $`\{x,w\}`$ belongs to the cover graph $`(S,_𝒮)`$ of $`𝒮`$. Moreover, $`(,x]`$ is (i) a down-set or (ii) an up-set in $`𝒮`$ according as (i) $`w`$ covers $`x`$ or (ii) $`x`$ covers $`w`$ in $`𝒮`$.
###### Proof..
If $`x=\tau `$, then $`(,\tau ]=S`$ is both a down-set and an up-set in $`𝒮`$. If $`x\tau `$, then $`x<_\tau \tau `$ and there is a unique predecessor of $`x`$; otherwise, the uniqueness of the path is contradicted. Let $`w`$ be the predecessor of $`x`$. Clearly, $`\{x,w\}`$ belongs to the cover graph of $`𝒮`$. Suppose that $`w`$ covers $`x`$ in $`𝒮`$. We claim that $`(,x]`$ is a down-set in $`𝒮`$, that is, that $`\eta (,x]`$ whenever $`\eta \xi `$ in $`𝒮`$ for some $`\xi (,x]`$. \[Since we have the same rooted tree $`(S,_\tau )`$ for the dual $`𝒮^{}`$, in proving the claim we will also settle the case that $`x`$ covers $`w`$ in $`𝒮`$.\] To see this, look at the paths from the root $`\tau `$ to $`\xi `$ and $`\eta `$, say, $`(u_0,\mathrm{},u_{n1})`$ from $`\tau =u_0`$ to $`u_{n1}=\xi `$ and $`(v_0,\mathrm{},v_{m1})`$ from $`\tau =v_0`$ to $`v_{m1}=\eta `$. For some $`k`$, the two paths descend the same vertices until the $`k`$th vertex, then split at the $`(k+1)`$st vertex. The path from $`\xi `$ to $`\eta `$ is then $`(u_{n1},\mathrm{},u_{k+1},u_k,v_{k+1},\mathrm{},v_{m1})`$, which is downward in $`𝒮`$ by assumption. That $`\xi _\tau x`$ implies that $`(u_{i1},u_i)=(w,x)`$ for some $`i`$. Furthermore, we have $`ik`$; otherwise, the downward path from $`\xi `$ to $`\eta `$ contains the directed edge $`(x,w)`$, which is impossible. Thus, the path from $`\tau `$ to $`\eta `$ contains the vertex $`x`$, which implies that $`\eta _\tau x`$.
###### Example 5.4.
Let $`𝒮`$ be the poset of Class Y displayed in Figure 5.1(a). By choosing the leaf $`\tau `$ of $`𝒮`$ as the root, we obtain the rooted tree $`(S,_\tau )`$ illustrated in Figure 5.1(b). For example, $`r`$ covers its predecessor $`q`$ in $`𝒮`$. By Lemma 5.3, the section $`(,r]=\{r,p,t,x,y,z\}`$ is an up-set in $`𝒮`$, which we can confirm immediately from Figure 5.1(a).
Now let $`P`$ be a probability measure on $`S`$. We define the distribution function of $`P`$ by
(5.3)
$$F(x):=P((,x])\text{ for each }xS\text{.}$$
It satisfies
(5.4)
$$F(\tau )=1$$
and
(5.5)
$$\underset{\xi 𝒞(x)}{}F(\xi )F(x)\text{ for every }xS\text{,}$$
where $`𝒞(x)`$ denotes the set of all successors of $`x`$ (and the summation is defined to be zero if $`𝒞(x)=\mathrm{}`$). Conversely, if a nonnegative function $`F`$ on $`S`$ satisfies the properties (5.4)–(5.5), then it is the distribution function of the probability measure $`P`$ determined uniquely via
(5.6)
$$P(\{x\}):=F(x)\underset{\xi 𝒞(x)}{}F(\xi )\text{ for each }xS\text{.}$$
Furthermore, stochastic ordering can be characterized in terms of distribution functions, as stated in the following lemma.
###### Lemma 5.5.
Let $`P_i`$ be a probability measure on $`S`$ and let $`F_i`$ be the distribution function of $`P_i`$, for each $`i=1,2`$. Then $`P_1P_2`$ if and only if for every $`xS`$ we have
* $`F_1(x)F_2(x)`$ if $`(,x]`$ is an up-set in $`𝒮`$, and
* $`F_1(x)F_2(x)`$ if $`(,x]`$ is a down-set in $`𝒮`$.
###### Proof..
By (1.2) and its trivial consequence Proposition 4.4(b), $`P_1P_2`$ clearly implies the conditions (i)–(ii). We proceed to the converse. Since any up-set $`U`$ in $`𝒮`$ is the disjoint union of the components $`V_1,\mathrm{},V_m`$ of the subgraph of $`(S,_𝒮)`$ induced by $`U`$ and $`V_1,\mathrm{},V_m`$ are all up-sets in $`𝒮`$, to prove $`P_1P_2`$ it suffices to show (1.2) for every up-set $`U`$ which induces a connected subgraph of $`(S,_𝒮)`$. If a set $`U`$ induces a connected subgraph of $`(S,_𝒮)`$, then we can find $`xU`$ and incomparable elements $`y_1,\mathrm{},y_k`$ of $`(,x]`$ in $`(S,_\tau )`$ so that
$$U=(,x]\left(\underset{i=1}{\overset{k}{}}(,y_i]\right),$$
where (as usual) the union is empty if $`k=0`$. Furthermore, suppose that $`U`$ is an up-set in $`𝒮`$. If $`x=\tau `$, then $`(,x]=S`$ is trivially an up-set in $`𝒮`$; otherwise, $`x`$ covers its predecessor $`w`$ in $`𝒮`$ and, by Lemma 5.3, $`(,x]`$ is an up-set in $`𝒮`$. Similarly, $`(,y_i]`$ is a down-set in $`𝒮`$ for each $`i=1,\mathrm{},k`$. Therefore, we have
$$P_1(U)=F_1(x)\underset{i=1}{\overset{k}{}}F_1(y_i)F_2(x)\underset{i=1}{\overset{k}{}}F_2(y_i)=P_2(U),$$
which establishes the sufficiency of the conditions (i)–(ii).
Because of Lemma 5.5, we write $`F_1F_2`$ if a pair $`(F_1,F_2)`$ of distribution functions on $`S`$ satisfies Lemma 5.5(i)–(ii) for every $`xS`$.
We now turn to the proof of Proposition 5.2. Let $`𝒮\mathrm{B}`$ and let $`F_{a,1},\mathrm{},F_{a,m}`$, $`F_{b,1},\mathrm{},F_{b,m}`$ be the distribution functions satisfying
$$F_{a,i}F_{b,j}\text{ for all }i=1,\mathrm{},m\text{ and all }j=1,\mathrm{},n\text{}$$
Then define the function $`\theta `$ on $`S`$ by
(5.7)
$$\theta (x):=\{\begin{array}{cc}\mathrm{max}\{F_{a,i}(x):i=1,\mathrm{},m\}\hfill & \text{ if }(,x]\text{ is an up-set in }𝒮\text{}\hfill \\ \mathrm{max}\{F_{b,j}(x):j=1,\mathrm{},n\}\hfill & \text{ if }(,x]\text{ is a down-set in }𝒮\text{}\hfill \end{array}$$
for $`xS`$. We first present the following lemma.
###### Lemma 5.6.
Let $`𝒮\mathrm{B}`$. Suppose that $`xS`$ and that $`v_1,\mathrm{},v_l`$ are mutually incomparable elements of $`(,x]`$ in $`(S,_\tau )`$.
* If $`(,x]`$ is a down-set in $`𝒮`$ and $`v_1,\mathrm{},v_lx`$ in $`𝒮`$, then
$$\underset{i=1}{\overset{l}{}}\theta (v_i)F_{a,j}(x)\text{for all }j=1,\mathrm{},m\text{.}$$
* If $`(,x]`$ is an up-set in $`𝒮`$ and $`v_1,\mathrm{},v_lx`$ in $`𝒮`$, then
$$\underset{i=1}{\overset{l}{}}\theta (v_i)F_{b,j^{}}(x)\text{for all }j^{}=1,\mathrm{},n\text{.}$$
###### Proof..
Suppose that the hypotheses in (a) hold. If $`v_1=x`$, then $`l=1`$ and the inequality clearly holds. Otherwise, $`v_ix`$ for every $`i=1,\mathrm{},l`$. Since the path $`(v_i,u_i,\mathrm{},x)`$ is upward and $`u_i`$ covers $`v_i`$ in $`𝒮`$, by Lemma 5.3(b) $`(,v_i]`$ is a down-set in $`𝒮`$. Therefore we have
$$\underset{i=1}{\overset{l}{}}\theta (v_i)\underset{i=1}{\overset{l}{}}F_{a,j}(v_i)F_{a,j}(x)$$
for all $`j=1,\mathrm{},m`$, as desired. The case (b) is reduced to (a) by considering the dual $`𝒮^{}`$.
We now define a nonnegative function $`F_0`$ on $`S`$ inductively. If $`x`$ is a minimal element in $`(S,_\tau )`$, then assign $`F_0(x):=\theta (x)`$. If $`x`$ is a nonminimal element in $`(S,_\tau )`$ and $`F_0(\xi )`$, $`\xi 𝒞(x)`$, have all been assigned, then set
$$F_0(x):=\mathrm{max}\{\theta (x),\underset{\xi 𝒞(x)}{}F_0(\xi )\}.$$
Clearly $`F_0`$ satisfies (5.5). We complete the proof of Proposition 5.2 by showing that $`F_0`$ satisfies (5.4) and
(5.8)
$$F_{a,i}F_0F_{b,j}\text{ for all }i=1,\mathrm{},m\text{ and all }j=1,\mathrm{},n\text{}$$
Thus, $`F_0`$ is a distribution function with the desired property.
###### Proof of Proposition 5.2..
We first claim that for every $`xS`$, there are incomparable elements $`v_1,\mathrm{},v_l`$ of $`(,x]`$ in $`(S,_\tau )`$ which satisfy both the hypotheses of one of Lemma 5.6(a),(b) and also
(5.9)
$$F_0(x)=\underset{i=1}{\overset{l}{}}\theta (v_i).$$
We will show this by induction over the cardinality of $`(,x]`$. If $`|(,x]|=1`$, then $`x`$ is a minimal element in $`(S,_\tau )`$ and indeed $`F_0(x)=\theta (x)`$.
Suppose that the claim holds for any $`xS`$ such that $`|(,x]|n1`$. Let $`xS`$ satisfy $`|(,x]|=n2`$. If $`x=\tau `$, then recall that $`\tau `$ is a leaf in $`𝒮`$ so that $`𝒞(\tau )`$ is a singleton, say $`\{y\}`$. By the induction hypothesis, we can find incomparable elements $`v_1,\mathrm{},v_l`$ of $`(,y]`$ in $`(S,_\tau )`$ which satisfy both the hypotheses of one of Lemma 5.6(a),(b) and (5.9) for $`y`$. If Lemma 5.6(a) obtains, then
$$F_0(y)=\underset{j=1}{\overset{l}{}}\theta (v_j)F_{a,1}(y)1.$$
A similar derivation concludes that $`F_0(y)1`$ when Lemma 5.6(b) obtains. Therefore, we have $`F_0(\tau )=\theta (\tau )=1`$, which proves (5.4). Furthermore, the claim holds for $`x=\tau `$.
If $`x\tau `$, then $`x`$ has a predecessor $`y_0`$ and successors $`y_1,\mathrm{},y_k`$ for some $`k1`$ (recalling our assumption $`|(,x]|=n2`$). Without loss of generality, we may assume that $`y_0`$ covers $`x`$ in $`𝒮`$. (We can treat the case that $`x`$ covers $`y_0`$ in $`𝒮`$ in exactly the same way by considering the dual $`𝒮^{}`$.) Then, by Lemma 5.3, $`(,x]`$ is a down-set in $`𝒮`$. Since $`𝒮`$ cannot have a bowtie as an induced subposet, only the following three cases can occur: (I) $`x`$ covers $`y_1,\mathrm{},y_k`$ in $`𝒮`$, or (II) $`y_1,\mathrm{},y_k`$ cover $`x`$ in $`𝒮`$, or, with $`k2`$, (III) $`y_1,\mathrm{},y_{k1}`$ cover $`x`$, and $`x`$ covers $`y_k`$, in $`𝒮`$. The induced subposet of $`𝒮`$ on $`\{x,y_0,y_1,\mathrm{},y_k\}`$ for each of these three cases is illustrated in the following figure.
| | | | | |
| --- | --- | --- | --- | --- |
| Case I | | Case II | | Case III |
Case I. If $`F_0(x)=\theta (x)`$, the claim is obvious. Otherwise, we have $`F_0(x)=_{i=1}^kF_0(y_i)`$. For each $`i=1,\mathrm{},k`$, the section $`(,y_i]`$ is a down-set in $`𝒮`$ by Lemma 5.3, and therefore by the induction hypothesis we have incomparable elements $`v_1^{(i)},\mathrm{},v_{l_i}^{(i)}`$ of $`(,y_i]`$ in $`(S,_\tau )`$ satisfying Lemma 5.6(a) and (5.9) for $`y_i`$. Thus we have
$$F_0(x)=\underset{i=1}{\overset{k}{}}F_0(y_i)=\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{l_i}{}}\theta (v_j^{(i)}),$$
and $`xv_j^{(i)}`$ for all $`i,j`$. Since the $`v_j^{(i)}`$’s are incomparable in $`(S,_\tau )`$, the claim holds for $`x`$.
Case II. For each $`i=1,\mathrm{},k`$, $`(,y_i]`$ is an up-set in $`𝒮`$ by Lemma 5.3, and therefore by the induction hypothesis we have incomparable elements $`v_1^{(i)},\mathrm{},v_{l_i}^{(i)}`$ of $`(,y_i]`$ in $`(S,_\tau )`$ satisfying Lemma 5.6(b) and (5.9) for $`y_i`$. By applying Lemma 5.6(b) to (5.9), we have
$$\underset{i=1}{\overset{k}{}}F_0(y_i)=\underset{i=1}{\overset{k}{}}\underset{j=1}{\overset{l_i}{}}\theta (v_j^{(i)})\underset{i=1}{\overset{k}{}}F_{b,1}(y_i)F_{b,1}(x)\theta (x),$$
which implies that $`F_0(x)=\theta (x)`$. Thus the claim holds for $`x`$.
Case III. By Lemma 5.3, $`(,y_k]`$ is a down-set in $`𝒮`$ and therefore by the induction hypothesis we can find incomparable elements $`v_1^{(k)},\mathrm{},v_{l_k}^{(k)}`$ of $`(,y_k]`$ in $`(S,_\tau )`$ satisfying Lemma 5.6(a) and (5.9) for $`y_k`$. Since $`𝒮`$ has no bowtie as an induced subposet and $`y_0,y_1y_kv_1^{(k)},\mathrm{},v_{l_k}^{(k)}`$ in $`𝒮`$, we have $`l_k=1`$. Thus, we can find some $`j_0`$ so that
$$F_0(y_k)=\theta (v_1^{(k)})=F_{b,j_0}(v_1^{(k)})F_{b,j_0}(y_k).$$
For each $`i=1,\mathrm{},k1`$, the section $`(,y_i]`$ is an up-set in $`𝒮`$ by Lemma 5.3, and therefore by the induction hypothesis we have incomparable elements $`v_1^{(i)},\mathrm{},v_{l_i}^{(i)}`$ of $`(,y_i]`$ in $`(S,_\tau )`$ satisfying Lemma 5.6(b) and (5.9) for $`y_i`$. By applying Lemma 5.6(b) to (5.9), we obtain
$`{\displaystyle \underset{i=1}{\overset{k}{}}}F_0(y_i)`$ $`F_{b,j_0}(y_k)+{\displaystyle \underset{i=1}{\overset{k1}{}}}{\displaystyle \underset{j=1}{\overset{l_i}{}}}\theta (v_j^{(i)})`$
$`{\displaystyle \underset{i=1}{\overset{k}{}}}F_{b,j_0}(y_i)F_{b,j_0}(x)\theta (x),`$
which implies that $`F_0(x)=\theta (x)`$. Thus, we have established the claim.
In order to show (5.8), it suffices to show that if $`(,x]`$ is an up-set in $`𝒮`$ then we have
$$F_{a,i}(x)F_0(x)F_{b,j}(x)\text{ for all }i=1,\mathrm{},m\text{ and all }j=1,\mathrm{},n\text{}$$
(Again, we can verify the case that $`(,x]`$ is a down-set in $`𝒮`$ by considering the dual $`𝒮^{}`$.) Suppose that $`(,x]`$ is an up-set in $`𝒮`$. Then we can find incomparable elements $`v_1,\mathrm{},v_l`$ of $`(,x]`$ in $`(S,_\tau )`$ satisfying Lemma 5.6(b) and (5.9). By applying Lemma 5.6(b) to (5.9), we have
$$F_{a,i}(x)\theta (x)F_0(x)=\underset{i=1}{\overset{l}{}}\theta (v_i)F_{b,j}(x).$$
This completes the proof.
### 5.2 Monotonicity inequivalence on Class Y
In this subsection, we present various examples, each with a poset $`𝒮`$ from Class Y, of posets $`𝒜`$ of monotonicity inequivalence. The next example turns out to be a building block for all the other examples.
###### Example 5.7.
Let $`𝒜_0`$ be the diamond given in (4.5) and let $`𝒮_0`$ be the Y-poset as in Figure 2.1(b). Define a system $`(P_a,P_b,P_c,P_d)`$ of probability measures on $`S_0`$ by
(5.10)
$$P_\alpha :=\{\begin{array}{cc}\text{unif}\{x,y\}\hfill & \text{ if }\alpha =a\text{}\hfill \\ \text{unif}\{x,w\}\hfill & \text{ if }\alpha =b\text{}\hfill \\ \text{unif}\{y,w\}\hfill & \text{ if }\alpha =c\text{}\hfill \\ \text{unif}\{z,w\}\hfill & \text{ if }\alpha =d\text{}\hfill \end{array}$$
It is clearly stochastically monotone with respect to $`(𝒜_0,𝒮_0)`$. We can prove that it is not realizably monotone by contradiction. Assume that there exists a system $`(𝐗_\alpha :\alpha A_0)`$ of $`S_0`$-valued random variables which realizes the monotonicity. Then we have
$`(𝐗_b=x)`$ $`=(𝐗_a=x,𝐗_b=x,𝐗_c=w,𝐗_d=w)={\displaystyle \frac{1}{2}},`$
$`(𝐗_c=y)`$ $`=(𝐗_a=y,𝐗_b=w,𝐗_c=y,𝐗_d=w)={\displaystyle \frac{1}{2}}.`$
Therefore, we have $`(𝐗_d=w)1`$, which contradicts the requirement $`(𝐗_d=w)=P_d(\{w\})=1/2`$. Thus monotonicity equivalence fails for $`(𝒜_0,𝒮_0)`$.
In Example 5.7, the dual $`𝒜_0^{}`$ is the diamond again. By Proposition 2.16, monotonicity equivalence fails for both $`(𝒜_0^{},𝒮_0)`$ and $`(𝒜_0,𝒮_0^{})`$. Now let $`𝒮`$ be any poset of Class Y. Since $`𝒮`$ has either the Y-poset $`𝒮_0`$ or its dual $`𝒮_0^{}`$ as an induced subposet, by Lemma 2.13 monotonicity equivalence fails for $`(𝒜_0,𝒮)`$. Thus, there exists a system $`(\stackrel{~}{P}_a,\stackrel{~}{P}_b,\stackrel{~}{P}_c,\stackrel{~}{P}_d)`$ of probability measures on $`S`$ which is stochastically but not realizably monotone with respect to $`(𝒜_0,𝒮)`$. In the next three examples (Examples 5.85.10), we take $`𝒮`$ to be any poset of Class Y.
###### Example 5.8.
Suppose that $`𝒜`$ has a cycle with height at least $`3`$. Then monotonicity equivalence fails for $`(𝒜,𝒮)`$.
To see this, observe by Lemma 2.6 that $`𝒜`$ has an induced cyclic subposet $`𝒜^{}=(a_0,a_1,\mathrm{},a_{n1},a_0)`$ with height at least $`3`$. Without loss of generality we may assume that $`𝒜^{}`$ has a maximal upward path $`(a_k^{},a_{k^{}+1},\mathrm{},a_{n1},a_0,a_1,\mathrm{},a_k)`$ with height at least $`3`$ for some $`1kk^{}n1`$, as illustrated in
(5.11)
$$𝒜^{}=\text{}$$
Note that $`k+2k^{}`$, since $`a_k`$ does not cover $`a_k^{}`$. Then we can define a system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ by
$$P_\alpha :=\{\begin{array}{cc}\stackrel{~}{P}_b\hfill & \text{ if }\alpha =a_0\hfill \\ \stackrel{~}{P}_d\hfill & \text{ if }a_0<\alpha \hfill \\ \stackrel{~}{P}_a\hfill & \text{ if }\alpha <a_0\hfill \\ \stackrel{~}{P}_c\hfill & \text{ otherwise.}\hfill \end{array}$$
Then the system is stochastically monotone.
We now show by contradiction that it is not realizably monotone. Suppose that we have a system $`(𝐗_\alpha :\alpha A)`$ of $`S`$-valued random variables which realizes the monotonicity. Then we have $`𝐗_{a_k^{}}𝐗_{a_0}𝐗_{a_k}`$. Since $`P_{a_{k+1}}=\mathrm{}=P_{a_{k^{}1}}=\stackrel{~}{P}_c`$, by applying Lemma 4.9 repeatedly we obtain (almost surely) $`𝐗_{a_{k+1}}=\mathrm{}=𝐗_{a_{k^{}1}}`$. Since $`a_{k+1}<a_k`$ and $`a_k^{}<a_{k^{}1}`$, we have $`𝐗_{a_k^{}}𝐗_{a_{k^{}1}}=𝐗_{a_{k+1}}𝐗_{a_k}`$. Therefore, the system $`(𝐗_{a_k^{}},𝐗_{a_0},𝐗_{a_{k+1}},𝐗_{a_k})`$ of $`S`$-valued random variables realizes the monotonicity of the system $`(\stackrel{~}{P}_a,\stackrel{~}{P}_b,\stackrel{~}{P}_c,\stackrel{~}{P}_d)`$ in terms of $`(𝒜_0,𝒮)`$. But this contradicts Example 5.7.
###### Example 5.9.
Suppose that $`𝒜`$ has a $`k`$-crown $`𝒜_k`$ as an induced subposet for some $`k3`$. Then monotonicity equivalence fails for $`(𝒜,𝒮)`$.
To see this, let $`𝒜_k`$ be as labeled in (4.6), let $`U:=\{\alpha A:a_0\alpha \}`$ be the up-set in $`𝒜`$ generated by $`a_0`$, and let $`V:=b_{k1}`$ be the down-set in $`𝒜`$ generated by $`b_{k1}`$. Then we define a system $`(P_\alpha :\alpha A)`$ of probability measures on $`S`$ by
$$P_\alpha :=\{\begin{array}{cc}\stackrel{~}{P}_b\hfill & \text{ if }\alpha UV\hfill \\ \stackrel{~}{P}_d\hfill & \text{ if }\alpha UV^c\hfill \\ \stackrel{~}{P}_a\hfill & \text{ if }\alpha U^cV\hfill \\ \stackrel{~}{P}_c\hfill & \text{ otherwise (i.e., }\alpha UV\text{).}\hfill \end{array}$$
Suppose that $`\alpha <\beta `$ in $`𝒜`$. If $`\alpha U^cV`$, then $`P_\alpha =\stackrel{~}{P}_aP_\beta `$. If $`\alpha UV`$, then $`\beta U`$ and $`P_\beta `$ is either $`\stackrel{~}{P}_b`$ or $`\stackrel{~}{P}_d`$; thus, $`P_\alpha =\stackrel{~}{P}_bP_\beta `$. If $`\alpha UV`$, then $`\beta V`$ and $`P_\beta `$ is either $`\stackrel{~}{P}_c`$ or $`\stackrel{~}{P}_d`$; thus, $`P_\alpha =\stackrel{~}{P}_cP_\beta `$. If $`\alpha UV^c`$, then $`\beta UV^c`$, and $`P_\alpha =\stackrel{~}{P}_d\stackrel{~}{P}_d=P_\beta `$. In each case that $`\alpha <\beta `$, we have shown $`P_\alpha P_\beta `$. Therefore, the system is stochastically monotone.
Suppose now that we have a system $`(𝐗_\alpha :\alpha A)`$ of $`S`$-valued random variables which realizes the monotonicity. Since $`P_{a_0}=P_{b_{k1}}=\stackrel{~}{P}_b`$, by Lemma 4.9 we (almost surely) have $`𝐗_{a_0}=𝐗_{b_{k1}}`$ and therefore $`𝐗_{a_{k1}}𝐗_{a_0}𝐗_{b_0}`$. Since $`P_{a_1}=P_{b_1}=\mathrm{}=P_{a_{k2}}=P_{b_{k2}}=\stackrel{~}{P}_c`$, by applying Lemma 4.9 repeatedly we obtain (almost surely) $`𝐗_{a_1}=𝐗_{b_1}=\mathrm{}=𝐗_{a_{k2}}=𝐗_{b_{k2}}`$ and therefore $`𝐗_{a_{k1}}𝐗_{a_1}𝐗_{b_0}`$, which implies that the system $`(𝐗_{a_{k1}},𝐗_{a_0},𝐗_{a_1},𝐗_{b_0})`$ of $`S`$-valued random variables realizes the monotonicity of the system $`(\stackrel{~}{P}_a,\stackrel{~}{P}_b,\stackrel{~}{P}_c,\stackrel{~}{P}_d)`$ indexed by the diamond (4.5). But this contradicts the discussion following Example 5.7. Hence $`(P_\alpha :\alpha A)`$ is not realizably monotone.
###### Example 5.10.
Suppose that $`𝒜`$ has
(5.12)
$$𝒜^{}=\text{}$$
as an induced subposet. Then monotonicity equivalence fails for $`(𝒜,𝒮)`$.
To see this, define a system $`𝒫=(P_\alpha :\alpha A)`$ of probability measures on $`S`$ by
$$P_\alpha :=\{\begin{array}{cc}\stackrel{~}{P}_b\hfill & \text{ if }\alpha b_1b_3\text{ }\hfill \\ \stackrel{~}{P}_c\hfill & \text{ if }\alpha b_3b_1\text{ }\hfill \\ \stackrel{~}{P}_a\hfill & \text{ if }\alpha b_1b_3\text{ }\hfill \\ \stackrel{~}{P}_d\hfill & \text{ otherwise (i.e., }\alpha b_1,b_3\text{). }\hfill \end{array}$$
Suppose that $`\alpha <\beta `$ in $`𝒜`$. If $`\alpha b_1b_3`$, then $`P_\alpha =\stackrel{~}{P}_aP_\beta `$. If $`\alpha b_1b_3`$, then $`\beta b_3`$ and $`P_\beta `$ is either $`\stackrel{~}{P}_b`$ or $`\stackrel{~}{P}_d`$; thus, $`P_\alpha =\stackrel{~}{P}_bP_\beta `$. If $`\alpha b_3b_1`$, then $`\beta b_1`$ and $`P_\beta `$ is either $`\stackrel{~}{P}_c`$ or $`\stackrel{~}{P}_d`$; thus, $`P_\alpha =\stackrel{~}{P}_cP_\beta `$. If $`\alpha b_1,b_3`$, then $`\beta b_1,b_3`$; thus, $`P_\alpha =\stackrel{~}{P}_d\stackrel{~}{P}_d=P_\beta `$. Therefore, the system is stochastically monotone.
Suppose now that there exists a system $`(𝐗_\alpha :\alpha A)`$ of $`S`$-valued random variables which realizes the monotonicity. By Lemma 4.9, we (almost surely) have $`𝐗_{a_1}=𝐗_{b_1}`$ and $`𝐗_{a_3}=𝐗_{b_3}`$. Thus we have found a system $`(𝐗_{a_2},𝐗_{a_1},𝐗_{a_3},𝐗_{b_2})`$ of $`S`$-valued random variables which realizes the monotonicity of the system $`(\stackrel{~}{P}_a,\stackrel{~}{P}_b,\stackrel{~}{P}_c,\stackrel{~}{P}_d)`$ indexed by the diamond (4.5). But this again contradicts the discussion following Example 5.7. Therefore, $`(P_\alpha :\alpha A)`$ is not realizably monotone.
### 5.3 The proof of Theorem 5.1
At the beginning of this Section 5, we saw that a non-acyclic poset $`𝒜`$ can sometimes be enlarged to an acyclic poset $`\stackrel{~}{𝒜}`$. But such an enlargement is not always possible. In fact, by Proposition 2.7, if $`𝒜`$ has an induced subposet which is one of the posets (i)–(iv) in Proposition 2.7 \[i.e., (i) the diamond, (ii) a subdivided crown with height at least $`3`$, (iii) the $`k`$-crown for some $`k3`$, or (iv) the double-bowtie poset\], then such an enlargement is not possible. It turns out that a non-acyclic poset can be enlarged to an acyclic poset if and only if none of the posets (i)–(iv) in Proposition 2.7 is an induced subposet; this relates the examples in Section 5.2 to Theorem 5.1.
###### Proposition 5.11.
Let $`𝒜`$ be a connected poset. The following conditions (a)(c) for $`𝒜`$ are equivalent:
* there exists an acyclic poset $`\stackrel{~}{𝒜}`$ which has $`𝒜`$ as an induced subposet;
* any induced cyclic subposet of $`𝒜`$ is a $`2`$-crown, and no induced subposet of $`𝒜`$ is the double-bowtie (5.12);
* no induced subposet of $`𝒜`$ is one of the posets (i)(iv) in Proposition 2.7.
###### Proof of (a) $``$ (c) and (c) $``$ (b)..
Suppose that there exists an induced subposet $``$ of $`𝒜`$ which is poset-isomorphic to one of the posets (i)–(iv) in Proposition 2.7. If there is an acyclic poset $`\stackrel{~}{𝒜}`$ which has $`𝒜`$ as an induced subposet, then $``$ is also an induced subposet of $`\stackrel{~}{𝒜}`$; by Proposition 2.7, this is impossible. We have thus shown that (a) $``$ (c).
To prove (c) $``$ (b), observe that a cycle is simply a subdivision of either the diamond or a crown. So if $`𝒜`$ has an induced cyclic subposet $``$ which is not a $`2`$-crown, then we can find an induced subposet of $``$ (automatically, of course, an induced subposet of $`𝒜`$) that is one of the posets (i)–(iii) in Proposition 2.7. Thus, the failure to satisfy the condition (b) implies the existence of an induced subposet which is one of the posets (i)–(iv) in Proposition 2.7.
In preparation for proving (b) $``$ (a), we introduce a new operation that welds two posets into one, as follows. Suppose that two posets $`𝒜^{}`$ and $`𝒜^{\prime \prime }`$ share a single element $`c`$ (i.e., that $`A^{}A^{\prime \prime }=\{c\}`$). Then the two Hasse diagrams of $`𝒜^{}`$ and $`𝒜^{\prime \prime }`$ can be drawn in the same plane with their own vertices and arcs independently except for the vertex $`c`$ to be shared by the two diagrams; this introduces a merged diagram on the vertex set $`A^{}A^{\prime \prime }`$. We call the poset represented by this Hasse diagram the union of $`𝒜^{}`$ and $`𝒜^{\prime \prime }`$ joined at $`c`$ and denote it by $`𝒜^{}\stackrel{c}{}𝒜^{\prime \prime }`$. We list some easily verified properties of the welding operation here:
* If $`𝒜^{}`$ and $`𝒜^{\prime \prime }`$ are connected, then $`𝒜^{}\stackrel{c}{}𝒜^{\prime \prime }`$ is connected.
* If $`𝒜^{}`$ and $`𝒜^{\prime \prime }`$ are acyclic, then $`𝒜^{}\stackrel{c}{}𝒜^{\prime \prime }`$ is acyclic.
* Both $`𝒜^{}`$ and $`𝒜^{\prime \prime }`$ are induced subposets of $`𝒜^{}\stackrel{c}{}𝒜^{\prime \prime }`$.
* Suppose that two posets $`𝒜_0^{}`$ and $`𝒜_0^{\prime \prime }`$ share a single element $`c`$ and that $`𝒜^{}`$ and $`𝒜^{\prime \prime }`$ (also sharing $`c`$) are induced subposets of $`𝒜_0^{}`$ and $`𝒜_0^{\prime \prime }`$, respectively. Then $`𝒜^{}\stackrel{c}{}𝒜^{\prime \prime }`$ is an induced subposet of $`𝒜_0^{}\stackrel{c}{}𝒜_0^{\prime \prime }`$.
We now continue our preparation for the proof of (b) $``$ (a) in Proposition 5.11. Lemma 5.12 provides machinery to split a poset into two smaller ones; this enables us to devise induction arguments in proving both Theorem 5.1 and (b) $``$ (a) in Proposition 5.11.
###### Lemma 5.12.
Let $`𝒜`$ be a connected non-acyclic poset. Suppose that $`𝒜`$ satisfies the condition (b) of Proposition 5.11. Then we can construct a pair $`𝒜_0^{}`$ and $`𝒜_0^{\prime \prime }`$ of connected posets (with ground sets $`A_0^{}`$ and $`A_0^{\prime \prime }`$, respectively) such that, for some $`c`$,
* both $`𝒜_0^{}`$ and $`𝒜_0^{\prime \prime }`$ satisfy the condition (b) of Proposition 5.11;
* $`A_0^{}A_0^{\prime \prime }=\{c\}`$, $`(A_0^{}A_0^{\prime \prime })\{c\}=A`$, and $`|A_0^{}|,|A_0^{\prime \prime }|<|A|`$; and
* $`𝒜`$ is the subposet of $`𝒜_0=𝒜_0^{}\stackrel{c}{}𝒜_0^{\prime \prime }`$ induced by $`A`$.
###### Proof..
Let $`𝒢`$ be the collection of all subsets $`B`$ of $`A`$ such that the subposet via induced cover subgraph of $`𝒜`$ on the ground set $`B`$ is a poset of the form (5.1) for some $`m,n2`$. Note that if a subposet $``$ via induced cover subgraph is of the form (5.1) then $``$ is an induced subposet. By Lemma 2.4 and the condition (b) of Proposition 5.11, $`𝒜`$ has an induced cyclic subposet which is necessarily a $`2`$-crown; thus, $`𝒢`$ is nonempty. Choose a maximal subset $`B_0`$ from $`𝒢`$. Then let $`_0`$ be the subposet of $`𝒜`$ induced by $`B_0`$ and label it as in (5.1).
Consider the Hasse diagram of $`𝒜`$ as represented in the plane. First remove the arcs from each of $`a_1,\mathrm{},a_m`$ to each of $`b_1,\mathrm{},b_n`$. Since the elements $`b_1,\mathrm{},b_n`$ are all drawn above the elements $`a_1,\mathrm{},a_m`$, we can insert a new vertex $`c`$ above the elements $`a_1,\mathrm{},a_m`$ but below the elements $`b_1,\mathrm{},b_n`$, and then add new arcs from $`a_1,\mathrm{},a_m`$ to $`c`$ and from $`c`$ to $`b_1,\mathrm{},b_n`$. This creates a new poset $`𝒜_0`$ with ground set $`A_0:=A\{c\}`$, as illustrated in
$`𝒜_0=`$
The subposet of $`𝒜_0`$ induced by $`A`$ introduces the arc from each of $`a_1,\mathrm{},a_m`$ to each of $`b_1,\mathrm{},b_n`$, thus restoring the Hasse diagram of $`𝒜`$.
We claim that $`𝒜_0`$ does not have any cycle which contains an upward path $`(a_i,c,b_j)`$ with $`a_i,b_jB_0`$. Granting the claim for the remainder of this paragraph, we define
$`A_0^{}`$ $`:=\{\alpha A_0:\text{ a path }(c,a_i,\mathrm{},\alpha )\text{ exists in }𝒜_0\text{ for some }a_iB_0\};`$
$`A_0^{\prime \prime }`$ $`:=\{\alpha A_0:\text{ a path }(c,b_j,\mathrm{},\alpha )\text{ exists in }𝒜_0\text{ for some }b_jB_0\}.`$
By convention, we include $`c`$ both in $`A_0^{}`$ and in $`A_0^{\prime \prime }`$. Let $`𝒜_0^{}`$ and $`𝒜_0^{\prime \prime }`$ be the subposets of $`𝒜_0`$ induced by $`A_0^{}`$ and $`A_0^{\prime \prime }`$, respectively. Clearly, $`𝒜_0^{}`$ and $`𝒜_0^{\prime \prime }`$ are both connected posets. By observing that $`𝒜_0`$ is connected, we find $`A_0^{}A_0^{\prime \prime }=A_0`$. The claim implies that $`A_0^{}A_0^{\prime \prime }=\{c\}`$ and that there are no edges between $`A^{}:=A_0^{}\{c\}`$ and $`A^{\prime \prime }:=A_0^{\prime \prime }\{c\}`$ in the cover graph of $`𝒜_0`$. Therefore, $`𝒜_0=𝒜_0^{}\stackrel{c}{}𝒜_0^{\prime \prime }`$, which implies (iii). Since $`\{a_1,\mathrm{},a_m\}A_0^{}`$, $`\{b_1,\mathrm{},b_n\}A_0^{\prime \prime }`$, and $`m,n2`$, we have $`|A_0^{}|,|A_0^{\prime \prime }|<|A|`$, as desired in (ii). To see (i) for $`𝒜_0^{}`$ (the same argument works for $`𝒜_0^{\prime \prime }`$), suppose that $`𝒜_0^{}`$ has an induced subposet $`𝒱`$ which violates the condition (b) of Proposition 5.11. If $`VA^{}`$, then let $`V^{}:=V`$; otherwise, let $`V^{}:=(V\{c\})\{b_1\}`$. Then the subposet $`𝒱^{}`$ of $`𝒜`$ induced by $`V^{}`$ is poset-isomorphic to $`𝒱`$. Furthermore, if $`𝒱`$ is a cycle in $`𝒜_0^{}`$, then $`𝒱^{}`$ is so in $`𝒜`$. But the existence of such an induced subposet of $`𝒜`$ contradicts the assumed condition (b) of Proposition 5.11. So (i) holds, and Lemma 5.12 is established modulo a proof of the claim.
Now we show the claim. To prove this by contradiction, we further enlarge the poset $`𝒜_0`$ as follows. We first remove the arcs from $`c`$ to each of $`b_1,\mathrm{},b_n`$ from the Hasse diagram of $`𝒜_0`$. Then a new element $`c^{}`$ is drawn above $`c`$ but below each of $`b_1,\mathrm{},b_n`$, and the arc from $`c`$ to $`c^{}`$ and the arcs from $`c^{}`$ to each of $`b_1,\mathrm{},b_n`$ are introduced in the diagram. This creates a new poset $`𝒜_1`$, as illustrated in
$`𝒜_1=`$
Clearly, $`𝒜_0`$ is an induced subposet of $`𝒜_1`$. If we can show that $`𝒜_1`$ has no cycle which contains the edge $`\{c,c^{}\}`$, then $`𝒜_0`$ has no cycle which contains an upward path $`(a_i,c,b_j)`$ for some $`a_i,b_jB_0`$, establishing the claim and Lemma 5.12. Thus, to obtain a contradiction, suppose that $`𝒜_1`$ has a cycle which contains the edge $`\{c,c^{}\}`$. By Lemma 2.4, we may assume that such a cycle, say $`(a_1,c,c^{},b_1,u_1,\mathrm{},u_k,a_1)`$, is an induced cyclic subposet of $`𝒜_1`$. Then the cycle $`(a_1,b_1,u_1,\mathrm{},u_k,a_1)`$ in $`𝒜`$ is an induced subposet of $`𝒜`$, and therefore by condition (b) of Proposition 5.11 the induced cyclic subposet $`(a_1,b_1,u_1,\mathrm{},u_k,a_1)`$ of $`𝒜`$ must be a $`2`$-crown, and therefore $`k=2`$. Note that $`u_iB_0`$ for $`i=1,2`$; otherwise, the cycle $`(a_1,c,c^{},b_1,u_1,u_2,a_1)`$ cannot be an induced subposet of $`𝒜_1`$.
Write $`a_0:=u_1`$ and $`b_0:=u_2`$. Now consider the comparability in $`𝒜`$ between $`\{a_0,b_0\}`$ and $`B_0`$. If $`a_0`$ is comparable with some $`a_i`$ of $`B_0`$, then either $`a_0<a_i<b_1`$ or $`a_i<a_0<b_1`$, contradicting our knowledge that $`\{a_0,b_1\}`$ and $`\{a_i,b_1\}`$ are edges in the cover graph of $`𝒜`$. Thus $`a_0`$ is incomparable with each of $`a_1,\mathrm{},a_m`$. Similarly we can see that $`b_0`$ is incomparable with each of $`b_1,\mathrm{},b_n`$. If $`a_0`$ is comparable with some $`b_j`$ of $`B_0`$ with $`j2`$, then $`a_0<b_j`$; otherwise, $`b_j<a_0<b_1`$, contradicting the assumption that $`b_1`$ and $`b_j`$ are incomparable. Suppose that there is an upward path $`(v_1,\mathrm{},v_l)`$ in $`𝒜`$ from $`v_1=a_0`$ to $`v_l=b_j`$ with $`l3`$. Then it is not hard to see that the cycle $`(a_0,v_2,\mathrm{},v_l,a_1,b_0,a_0)`$ in $`𝒜`$ is an induced subposet of $`𝒜`$ with height $`l`$, contradicting condition (b) of Proposition 5.11. Therefore, $`b_j`$ must cover $`a_0`$ in $`𝒜`$. Note that $`a_0`$ cannot be comparable with all of $`b_2,\mathrm{},b_n`$, since $`B_0\{a_0\}𝒢`$. Similarly, $`b_0`$ cannot be comparable with all of $`a_2,\mathrm{},a_m`$ (but $`b_0`$ may cover some of them). Thus, we can find some elements $`a_i`$, $`b_j`$ of $`B_0`$ so that $`a_0`$ is incomparable with $`b_j`$ and $`b_0`$ is incomparable with $`a_i`$. We have now found the subposet of $`𝒜`$ induced by $`\{a_0,a_1,a_i,b_0,b_1,b_j\}`$ to be poset-isomorphic to the poset (5.12). This contradicts condition (b) of Proposition 5.11.
We now give the proof for (b) $``$ (a) in Proposition 5.11, by induction over the cardinality of $`A`$. The idea of the proof is to build up an acyclic poset by using Lemma 5.12.
###### Proof of (b) $``$ (a) in Proposition 5.11..
Suppose that a poset $`𝒜`$ satisfies the condition (b). We will make an inductive argument over the cardinality of $`A`$. But if $`𝒜`$ is acyclic, then the argument is vacuous. In particular, $`𝒜`$ with cardinality at most $`3`$ is acyclic. Now let $`𝒜`$ be a connected non-acyclic poset with cardinality $`n4`$. By Lemma 5.12, there exists a pair $`𝒜_0^{}`$ and $`𝒜_0^{\prime \prime }`$ of connected posets satisfying (i)–(iii) in Lemma 5.12. Then, by the induction hypothesis and (i)–(ii) in Lemma 5.12, $`𝒜_0^{}`$ and $`𝒜_0^{\prime \prime }`$ can be enlarged to acyclic posets $`\stackrel{~}{𝒜}^{}`$ and $`\stackrel{~}{𝒜}^{\prime \prime }`$, respectively.
Since $`A_0^{}A_0^{\prime \prime }=\{c\}`$, the ground sets $`\stackrel{~}{A}^{}`$ and $`\stackrel{~}{A}^{\prime \prime }`$ can be given so that $`\stackrel{~}{A}^{}\stackrel{~}{A}^{\prime \prime }=\{c\}`$. Let $`\stackrel{~}{𝒜}:=\stackrel{~}{𝒜}^{}\stackrel{c}{}\stackrel{~}{𝒜}^{\prime \prime }`$. Then $`\stackrel{~}{𝒜}`$ is acyclic. Furthermore, $`𝒜_0=𝒜_0^{}\stackrel{c}{}𝒜_0^{\prime \prime }`$ is an induced subposet of $`\stackrel{~}{𝒜}`$. By (iii) in Lemma 5.12, $`𝒜`$ is an induced subposet of $`\stackrel{~}{𝒜}`$, as desired.
Now we turn to the proof of Theorem 5.1. The proof will parallel that of (b) $``$ (a) in Proposition 5.11 somewhat. Let $`𝒜`$ be a connected poset and let $`𝒮`$ be a poset of Class Y.
###### Proof of Theorem 5.1..
Suppose first that a non-acyclic poset $`𝒜`$ is not enlargeable to an acyclic poset. Then, by Proposition 5.11, $`𝒜`$ has an induced subposet $``$ which is one of the posets (i)–(iv) in Proposition 2.7. If $``$ is the diamond, then by Lemma 2.5 $`𝒜`$ has a cycle with height at least $`3`$. Thus Example 5.8 implies that monotonicity equivalence fails for $`(𝒜,𝒮)`$. If $``$ is either the $`k`$-crown for some $`k3`$ or the double-bowtie poset (5.12), then by Examples 5.95.10 monotonicity equivalence fails for $`(𝒜,𝒮)`$. Suppose now that $``$ is a subdivision of the $`k`$-crown as displayed and labeled in (4.6) and has height at least $`3`$. Then we may assume that there exists $`c_0B`$ such that $`a_0<c_0<b_0`$ in $``$. So we find an upward path $`(a_0,\mathrm{},c^{},c_0,c^{\prime \prime },\mathrm{},b_0)`$ in $`𝒜`$ from $`a_0`$ to $`b_0`$ with height at least $`3`$. Since $`c_0`$ is incomparable in $`𝒜`$ with each of $`a_1,\mathrm{},a_{k1}`$ and each of $`b_1,\mathrm{},b_{k1}`$, no upward path in $`𝒜`$ from any $`a_i`$ to any $`b_j`$ contains either $`\{c^{},c_0\}`$ or $`\{c_0,c^{\prime \prime }\}`$ as an edge unless $`(i,j)=(0,0)`$. Let $`(A,_𝒜)`$ be the cover graph of $`𝒜`$. Then there exists a path $`(u_1,\mathrm{},u_k)`$ from $`u_1=c^{\prime \prime }`$ to $`u_k=c^{}`$ in the graph $`(A,_𝒜\{\{c^{},c_0\},\{c_0,c^{\prime \prime }\}\})`$. Thus, $`𝒜`$ has a cycle $`(c^{},c_0,c^{\prime \prime }=u_1,\mathrm{},u_k)`$ with height at least $`3`$. Therefore, by Example 5.8 monotonicity equivalence fails for $`(𝒜,𝒮)`$.
Suppose now that a poset $`𝒜`$ is enlargeable to an acyclic poset. We will prove that monotonicity equivalence holds for $`(𝒜,𝒮)`$ by induction over the cardinality of $`A`$. If $`𝒜`$ is acyclic, then, by Theorem 4.1, $`𝒜`$ is a poset of monotonicity equivalence. Thus, if $`|A|3`$, then $`𝒜`$ is acyclic and therefore a poset of monotonicity equivalence. Now let $`𝒜`$ be a non-acyclic poset with cardinality $`n4`$ and let $`(P_\alpha :\alpha A)`$ be a stochastically monotone system of probability measures on $`S`$. By Lemma 5.12, there exists a pair $`𝒜_0^{}`$ and $`𝒜_0^{\prime \prime }`$ of posets satisfying (i)–(iii) in Lemma 5.12. Let $`a_1,\mathrm{},a_m`$ be all the elements covered by $`c`$, and let $`b_1,\mathrm{},b_n`$ be all the elements covering $`c`$ in $`𝒜_0=𝒜_0^{}\stackrel{c}{}𝒜_0^{\prime \prime }`$. Since $`𝒜`$ is an induced subposet of $`𝒜_0`$, we have $`P_{a_i}P_{b_j}`$ for all $`i=1,\mathrm{},m`$ and all $`j=1,\mathrm{},n`$. By Proposition 5.2, we can find a probability measure $`P_0`$ on $`S`$ such that $`P_{a_i}P_0P_{b_j}`$ for all $`i=1,\mathrm{},m`$ and all $`j=1,\mathrm{},n`$.
Let $`P_c:=P_0`$. Then we can enlarge the system $`(P_\alpha :\alpha A)`$ to a system $`(P_\alpha :\alpha A_0)`$, maintaining stochastic monotonicity. Note that the subsystems $`(P_\alpha :\alpha A_0^{})`$ and $`(P_\alpha :\alpha A_0^{\prime \prime })`$ are also stochastically monotone. Since \[by Lemma 5.12(i) and Proposition 5.11\] $`𝒜_0^{}`$ \[respectively, $`𝒜_0^{\prime \prime }`$\] is enlargeable to an acyclic poset, by the induction hypothesis there is a system $`(𝐗_\alpha ^{}:\alpha A_0^{})`$ \[respectively, $`(𝐗_\alpha ^{\prime \prime }:\alpha A_0^{\prime \prime })`$\] of $`S`$-valued random variables which realizes the monotonicity of $`(P_\alpha :\alpha A_0^{})`$ \[respectively, of $`(P_\alpha :\alpha A_0^{\prime \prime })`$\]. Let $`A^{}:=A_0^{}\{c\}`$ and $`A^{\prime \prime }:=A_0^{\prime \prime }\{c\}`$. We can define a probability measure $`Q`$ on $`S^{A_0}=S^A^{}\times S^{\{c\}}\times S^{A^{\prime \prime }}`$ by
$$Q(\{(𝐱^{},\xi ,𝐱^{\prime \prime })\}):=Q^{}(𝐱^{}|\xi )Q^{\prime \prime }(\xi ,𝐱^{\prime \prime })\text{for }(𝐱^{},\xi ,𝐱^{\prime \prime })S^A^{}\times S^{\{c\}}\times S^{A^{\prime \prime }}\text{,}$$
where
$`Q^{}(𝐱^{}|\xi )`$ $`:=P(𝐗_\alpha ^{}=\pi _\alpha (𝐱^{})\alpha A^{}|𝐗_c^{}=\xi );`$
$`Q^{\prime \prime }(\xi ,𝐱^{\prime \prime })`$ $`:=P(𝐗_c^{\prime \prime }=\xi ,𝐗_\alpha ^{\prime \prime }=\pi _\alpha (𝐱^{\prime \prime })\alpha A^{\prime \prime }).`$
Then $`Q`$ realizes the monotonicity of $`(P_\alpha :\alpha A_0)`$. By Lemma 2.12, $`(P_\alpha :\alpha A)`$ is realizably monotone; thus, monotonicity equivalence holds for $`(𝒜,𝒮)`$.
###### Remark 5.13.
In the second part of the proof of Theorem 5.1, we invoke Proposition 5.2, which requires only $`𝒮\mathrm{B}`$. Thus, we have actually proved that if $`𝒮\mathrm{B}`$ and $`𝒜`$ is enlargeable to an acyclic poset then monotonicity equivalence holds for $`(𝒜,𝒮)`$.
## 6 Probability measures on a path
In Section 5 we have seen that, when $`𝒮\mathrm{B}`$, stochastic ordering can be decided from the distribution function (5.3). In this Section 6 we establish that the inverse probability transform (6.1) can be used to realize monotonicity when $`𝒮\mathrm{Z}`$; this result extends Example 1.2. As a result, we will obtain
###### Theorem 6.1.
Let $`𝒜`$ be any poset and let $`𝒮`$ be a poset of Class Z. Then monotonicity equivalence holds for $`(𝒜,𝒮)`$.
Let $`𝒮`$ be a poset of Class Z. As we observed in Section 3, $`𝒮`$ is poset-isomorphic to a path, say $`(x_1,\mathrm{},x_n)`$. So a natural linear order $`_n`$ of the path is introduced by declaring $`x_i_nx_j`$ if and only if $`ij`$. In other words, $`(S,_n)`$ is a rooted tree with root $`x_n`$ (see Section 5.1). Note that such a linear order $`_n`$ is not consistent in general with the partial order $``$ of the poset $`𝒮`$. In Figure 6.1 we give an example of (a) a poset $`𝒮`$ of Class Z and (b) its linear order $`_n`$. For every $`x_iS`$, the section $`(,x_i]:=\{x_jS:x_j_nx_i\}`$ of the path is either an up-set or a down-set in $`𝒮`$, which is obvious pictorially in Figure 6.1. In fact, the linearly ordered set $`(S,_n)`$ is a rooted tree with root $`x_n`$; thus, Lemma 5.3 applies.
For a probability measure $`P`$ on $`S`$, the distribution function $`F`$ of $`P`$ is given by (5.3), that is, $`F(x_i)=P((,x_i])`$ for each $`x_iS`$. Furthermore, we can define the inverse probability transform $`P^1`$ from $`[0,1)`$ to $`S`$ by
(6.1)
$$P^1(u):=\mathrm{min}\{x_k:u<F(x_k)\}\text{ for }u[0,1)\text{,}$$
where the minimum is given in terms of the linear order $`_n`$. Then we can state an equivalent condition for stochastic ordering as the following lemma.
###### Lemma 6.2.
Let $`(P_1,P_2)`$ be a pair of probability measures on $`S\mathrm{Z}`$. Then $`P_1P_2`$ if and only if
(6.2) $`P_1^1(u)P_2^1(u)`$ in $`𝒮`$ for all $`u[0,1)`$.
###### Proof..
Suppose first that $`P_1P_2`$. Let $`F_1`$ and $`F_2`$ denote the distribution functions of $`P_1`$ and $`P_2`$, respectively. Let $`u[0,1)`$ be fixed, $`x_i:=P_1^1(u)`$, and $`x_j:=P_2^1(u)`$. If $`x_i=x_j`$, then (6.2) obviously holds. If $`x_i<_nx_j`$, then we have
$$F_2(x_k)F_2(x_{j1})u<F_1(x_i)F_1(x_k)$$
for all $`x_k`$ such that $`x_i_nx_k<_nx_j`$. By Lemma 5.5, the section $`(,x_k]`$ is a down-set for every $`k=i,i+1,\mathrm{},j1`$, which implies, by Lemma 5.3, that $`x_i<x_{i+1}<\mathrm{}<x_{j1}<x_j`$ in $`𝒮`$. If $`x_j<_nx_i`$, then we have
$$F_1(x_k)F_1(x_{i1})u<F_2(x_j)F_2(x_k)$$
for all $`x_k`$ such that $`x_j_nx_k<_nx_i`$. Again by applying Lemmas 5.5 and 5.3, we obtain $`x_j>x_{j+1}>\mathrm{}>x_{i1}>x_i`$ in $`𝒮`$. In any case, (6.2) holds.
Now suppose that (6.2) holds. Then we can construct a pair $`(𝐗_1,𝐗_2)`$ of $`S`$-valued random variables satisfying (1.3)–(1.4) via
$$𝐗_i:=P_i^1(𝐔)\text{ for each }i=1,2$$
with a single random variable $`𝐔`$ uniformly distributed on $`[0,1)`$. By Proposition 4.4, we have $`P_1P_2`$. This completes the proof.
Lemma 6.2 is exactly the property needed to generalize Example 1.2 to the case where $`𝒮`$ is a poset of Class Z. We complete the
###### Proof of Theorem 6.1..
Let $`(P_\alpha :\alpha A)`$ be a stochastically monotone system of probability measures on $`S`$. Let $`𝐔`$ be a random variable uniformly distributed on $`[0,1)`$. Then we can construct a system $`(𝐗_\alpha :\alpha A)`$ of $`S`$-valued random variables satisfying (1.6) via
$$𝐗_\alpha :=P_\alpha ^1(𝐔)\text{ for }\alpha A\text{.}$$
By Lemma 6.2, the system $`(𝐗_\alpha :\alpha A)`$ satisfies (1.5); thus, $`(P_\alpha :\alpha A)`$ is realizably monotone.
Acknowledgments. We thank Keith Crank, Persi Diaconis, Alan Goldman, Leslie Hall, Robin Pemantle, and Edward Scheinerman for providing helpful comments.
James Allen Fill
Department of Mathematical Sciences
The Johns Hopkins University
Baltimore, MD 21218-2682
jimfill@jhu.edu
Motoya Machida
Department of Mathematical Sciences
The Johns Hopkins University
Baltimore, MD 21218-2682
machida@mts.jhu.edu
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# Quantum Remote Control: Teleportation of Unitary Operations
## Abstract
We consider the implementation of an arbitrary unitary operation U upon a distant quantum system. This teleportation of U can be viewed as a quantum remote control. We investigate protocols which achieve this using local operations, classical communication and shared entanglement (LOCCSE). Lower bounds on the necessary entanglement and classical communication are determined using causality and the linearity of quantum mechanics. We examine in particular detail the resources required if the remote control is to be implemented as a classical black box. Under these circumstances, we prove that the required resources are, necessarily, those needed for implementation by bidirectional state teleportation.
Much of the current fascination with quantum information processing derives from the properties of entanglement . On one hand, entanglement can give rise to nonlocal correlations which defy explanation in terms of local, realistic theories, but on the other, it can also be used as a resource. While it is impossible, for example, to determine the state of a quantum system, entanglement makes it possible to transmit an unknown state. This process is known as quantum state teleportation . Quantum state teleportation can be linked directly to various interrelated principles of quantum information processing, such as the impossibility of superluminal communication, the non-increasing of entanglement under local operations and classical communication and the no-cloning theorem . However, the information contained in the state of a quantum system is only one kind of information which is important in quantum mechanics. Another is the information which describes quantum operations. In this paper, we examine the issue of teleporting, not a quantum state, but rather a quantum operation. In particular, we examine the teleportation of an unknown unitary operation on a qubit. This procedure would function in a manner similar to that of a remote control apparatus, and so we shall also refer to it as quantum remote control. We will first pose the problem in a completely general theoretical framework and focus later on an experimentally feasible scenario where entanglement resources are limited.
The most general scenario for the teleportation of an arbitrary unitary operation is depicted in Figure 1. One party, Alice, possesses a physical system, C, which we shall refer to as the control. The control contains information describing a unitary operation U upon the state of a qubit, and is itself a quantum system. The control state corresponding to the unitary operation U will be denoted by $`|U_C`$. Her colleague Bob has a qubit $`\beta `$ prepared in the state $`|\psi _\beta `$. The aim is to devise a physical procedure which effects the transformation $`|\psi _\beta U|\psi _\beta `$, for every initial state $`|\psi _\beta `$ and every unitary operation $`U`$. The most general such procedure can be represented by a completely positive, linear, trace preserving map on the set of density operators for the combined $`C\beta `$ system. Any such map has a unitary representation $`𝒯`$ involving ancillary systems. We shall denote the state of the ancilla at Alice’s and Bob’s laboratories by $`|\chi _{AB}`$. Then the teleportation operation has the general form
$$𝒯\left[|\chi _{AB}|U_C|\psi _\beta \right]=|\mathrm{\Phi }(U,\chi )_{ABC}\left(U|\psi _\beta \right).$$
(1)
In the following we investigate some of the properties of $`𝒯`$. In particular, we derive lower bounds on the amount of non-local resources that are needed to implement $`𝒯`$ using only local operations and classical communication. The unitary teleportation operator $`𝒯`$ is independent of both $`U`$ and $`|\psi _\beta `$. The final state of the ancilla+control, $`|\mathrm{\Phi }(U,\chi )_{ABC}`$, must be independent of $`|\psi _\beta `$. To see why, let us suppose that it isn’t, in which case there will be at least one U, and two states, $`|\psi _\beta `$ and $`|\psi ^{^{}}_\beta `$, for which $`|\mathrm{\Phi }(U,\chi ,\psi )_{ABC}|\mathrm{\Phi }(U,\chi ,\psi ^{^{}})_{ABC}`$. We imagine that U is successfully teleported for the states $`|\psi _\beta `$ and $`|\psi ^{^{}}_\beta `$. Suppose now that Bob’s qubit is prepared in a superposition of these states, $`(c_1|\psi +c_2|\psi ^{^{}})_\beta `$. The linearity of $`𝒯`$ implies that the final total state will be
$$(𝟙_{ABC}U_\beta )\left[c_1|\mathrm{\Phi }(U,\chi ,\psi )_{ABC}|\psi _\beta +c_2|\mathrm{\Phi }(U,\chi ,\psi ^{^{}})_{ABC}|\psi ^{^{}}_\beta \right].$$
(2)
The requirement that Bob’s qubit undergoes a unitary evolution implies that it must not be entangled with the remaining systems. However, one can see that it is entangled with ABC whenever $`c_1c_20`$. Thus, the final state of ABC must be independent of $`|\psi _\beta `$.
The set of all unitary operations U is infinite. This implies that if the dimension of the control system is to be finite, then the control states $`|U_C`$ must, in general, be non-orthogonal. However, Nielsen and Chuang showed, in a slightly different context, that this cannot be the case . The problem investigated by these authors was whether or not one could devise a universal programmable quantum gate array, which could be used to store and execute any program upon a quantum register. They showed that no such finite array can be constructed. Their method of proof can readily be transferred to this context, making use of the correspondences between programmable gate array/control, and register/Bob’s qubit. Following their reasoning, we note that Eq. (1) and the unitarity of $`𝒯`$ imply that, for any two different unitary transformations $`U`$ and $`U^{}`$,
$$\frac{{}_{C}{}^{}U^{}|U_{C}^{}}{{}_{ABC}{}^{}\mathrm{\Phi }(U^{},\chi )|\mathrm{\Phi }(U,\chi )_{ABC}^{}}={}_{\beta }{}^{}\psi |U^{{}_{}{}^{}}U|\psi _{\beta }^{}.$$
(3)
The left hand side is independent of $`|\psi _\beta `$, and this equality is true for all $`|\psi _\beta `$. It follows that $`U^{{}_{}{}^{}}U=\gamma 𝟙`$, for some constant $`\gamma `$, leading to the conclusion that $`U`$ and $`U^{}`$ are identical up to a multiplicative constant. This conclusion, however, is valid only when the denominator on the left hand side is non-zero. If it is zero, then $`{}_{C}{}^{}U^{}|U_{C}^{}=0`$, by the unitarity of $`𝒯`$. Control states corresponding to different unitary transformations are orthogonal, so that no finite-dimensional control system can be used to teleport an arbitrary unitary operation. For the remainder of this paper, when we speak of an arbitrary unitary operation, we will mean one which belongs to some arbitrarily large, but finite, set. We will also assume that this set contains the identity $`\sigma ^0=𝟙`$ and the 3 Pauli operators $`\sigma ^i`$. Note that the orthogonality of the control states opens the possibility that different operations can, at least in principle, be distinguished by Alice. This will only be possible though if Alice knows the basis $`\{|U_C\}`$ in which the information is encoded.
The teleportation of $`U`$ is a collective operation on spatially separated systems, which we wish to carry out using shared entanglement and classical communication. In the derivation of lower bounds on the amount of non-local resources that are required to implement the teleportation of $`U`$ locally, two guiding principles will be very useful :
(i) The amount of classical information able to be communicated by an operation in a given direction across some partition between subsystems cannot exceed the amount of information that must be sent in this direction across the same partition to complete the operation.
(ii) The amount of bipartite entanglement that an operation can establish across some partition between subsystems cannot exceed the amount of prior entanglement across the partition that must be consumed in order to complete the operation.
We now use principle (i) to establish the fact that at least two classical bits must be sent from Alice to Bob to complete the teleportation of an arbitrary U. Suppose that, rather than being prepared in a pure state, Bob’s qubit is initially maximally entangled with some other qubit, $`\beta ^{^{}}`$, which is also in Bob’s laboratory. Let us denote the four Bell states for a pair of qubits by $`|B^\mu `$, where $`\mu =0,\mathrm{},3`$. Using the technique of super-dense coding, any of the four Bell states can be transformed into any other by application of one of the Pauli operators $`\sigma ^i`$ on one of the qubits. We take this qubit to be $`\beta `$, and notice that the $`|B^\mu `$ can be ordered in such a way that $`(\sigma _\beta ^\mu 𝟙_\beta ^{})|B^0_{\beta \beta ^{}}=|B^\mu _{\beta \beta ^{}}`$. Alice can easily transmit two bits of information to Bob if he prepares the $`\beta \beta ^{}`$ system in the state $`|B^0_{\beta \beta ^{}}`$. She chooses the control system to be in one of the states $`|\sigma ^\mu _C`$. Following the action of $`𝒯`$, Bob will be in possession of the corresponding Bell state $`|B^\mu _{\beta \beta ^{}}`$. If he subsequently performs a Bell measurement on $`\beta \beta ^{}`$, then he will be able to determine the value of $`\mu `$, and hence the control state which Alice prepared, revealing 2 bits of classical information.
We now show that, by teleporting an arbitrary U according to the general prescription in Eq. (1), Alice and Bob can establish 2 ebits of shared entanglement. Imagine that, in addition to the systems we have already introduced, Alice has a further 4-dimensional ancilla, which we shall label R. Let the states $`|\mu _R`$ be a particular orthonormal basis for R. Suppose now that Alice initially prepares R and the control C in the maximally entangled state $`(1/2)_\mu |\mu _R|\sigma ^\mu _C`$. Bob once more prepares the Bell state $`|B^0_{\beta \beta ^{}}`$. The teleportation operation $`𝒯`$ is then carried out according to Eq. (1). It is more convenient here, however, to work with a form of this equation that represents, explicitly, any local measurements made by Alice and Bob and any classical communication between them. In this case $`𝒯`$ in Eq. (1) is replaced by a pair of classically-correlated local CP maps, one in each laboratory. Classical information is revealed by measurements, and we let the index i denote each measurement outcome. The final state corresponding to the ith outcome is
$$|\psi _F_i=\frac{1}{2}\underset{\mu }{}|\mu _R|\mathrm{\Phi }_i(\sigma ^\mu ,\chi )_{ABC}|B^\mu _{\beta \beta ^{}}$$
(4)
We now calculate the entanglement shared by Alice and Bob. Alice is in possession of the compound system RAC, while Bob has the system $`B\beta \beta ^{}`$. For each outcome, these subsystems have respective density operators $`\rho _{RAC}^i`$ and $`\rho _{B\beta \beta ^{}}^i`$. Since $`|\psi _F_i`$ is a pure state, it follows that the entanglement shared by Alice and Bob is simply the (base 2) von Neumann entropy of either of these density operators. Fortunately, we can calculate this explicitly. To do so, we notice that the states $`|\mathrm{\Phi }_i(\sigma ^\mu ,\chi )_{ABC}`$ will generally contain entanglement between B and AC. Let us write $`\rho _B^{i\mu }=\mathrm{T}r_{AC}(|\mathrm{\Phi }_i(\sigma ^\mu ,\chi )\mathrm{\Phi }_i(\sigma ^\mu ,\chi )|)`$. We find that
$$\rho _{B\beta \beta ^{}}=\frac{1}{4}\underset{\mu }{}(|B^\mu B^\mu |)_{\beta \beta ^{}}\rho _B^{i\mu }.$$
(5)
Making use of the orthogonality of the $`|B^\mu `$, we find that the total entropy of entanglement shared by Alice and Bob is simply
$$E(|\psi _F)=S(\rho _{B\beta \beta ^{}})=2+\frac{1}{4}\underset{\mu }{}S(\rho _B^{i\mu })\mathrm{\hspace{0.17em}2}.$$
(6)
It follows from principle (ii) that at least 2 ebits of entanglement need to be consumed to implement $`𝒯`$ locally, i.e. to teleport an arbitrary unitary operation.
We can summarize the results obtained so far as follows. The resources required to perform quantum remote control can be classified into shared entanglement, classical information transmission from Alice to Bob, and from Bob to Alice. We have established absolute lower bounds on the first two of these resources. Alice and Bob have to share at least two ebits and Alice needs to transmit to Bob, at least, two bits of classical information.
These bounds can be attained by a procedure in which Bob teleports the state of his particle to Alice who, after applying the unitary transformation, teleports it back to him. We will call this the “bidirectional state teleportation” scheme. The scheme requires sending 2 classical bits in each direction, and using 2 ebits of shared entanglement. It would also be conceivable to adopt a different strategy – teleporting the state of the control system from Alice to Bob who would then implement the control directly onto $`\beta `$. We call this the “control-state teleportation” scheme.
Control-state teleportation is a unidirectional communication scheme from Alice to Bob, so the absolute lower bound for the communication exchange from Bob to Alice is zero. Obviously, the overall resources will depend on the dimensionality of the control system $`C`$. We cannot say anything about the optimality of this procedure; whether there exists another unidirectional protocol which uses less resources is an open problem.
On the other hand, bidirectional state teleportation saturates the lower bounds for the amount of shared ebits and classical bits transmitted from Alice to Bob and additionally uses two bits of classical communication from Bob to Alice. This scheme allows the faithful implementation of $`U`$ independently of the dimension of the control system. To be more efficient overall, any other scheme would need less resources than bidirectional state teleportation. This establishes an upper bound in the overall amount of resources required for the efficient remote implementation of an arbitrary $`U`$ as 4 classical bits and 2 ebits.
We now consider an experimental scenario where the black box implementing an arbitrary transformation $`U`$ is a macroscopic object, involving a (very) large number of degrees of freedom. The option of teleporting the control apparatus is then unfeasible, given that it would consume an infinite amount of entanglement and classical communication resources. However, the question remains whether there exists a more economical protocol than bidirectional state teleportation. We will prove in the following that this is not possible and bidirectional state teleportation is an unconditional optimal way to remotely implement an arbitrary $`U`$.
Discarding the possibility of control-state teleportation allows us to replace the transformation given by Eq. (1) with
$$G_2UG_1(|\chi _{\alpha AB}|\psi _\beta )=|\mathrm{\Phi }(U,\chi )_{\alpha AB}U|\psi _\beta ,$$
(7)
where certain fixed operations $`G_1`$ and $`G_2`$ are performed, respectively, prior to and following the action of the arbitrary $`U`$ on a qubit $`\alpha `$ on Alice’s side. We assume that Alice and Bob share initially some entanglement, represented by the state $`|\chi _{\alpha AB}`$. As before, the purpose of the transformation is to perform the operation $`U`$ on Bob’s qubit $`\beta `$. We continue to use a nonlocal unitary representation of the transformation where $`G_1`$ and $`G_2`$ are unitary operators acting on possibly all subsystems. A pictorial scheme of the situation using a quantum circuit is given in Figure 2. The two upper wires refer to Alice’s subsystems and the two lower ones to Bob’s. Note that operations $`G_i`$ are represented by non-local gates while the action of $`U`$ takes place locally on Alice’s side.
We prove in the following that the only way that Eq. (7) can be implemented (locally) is by teleporting the state $`|\psi _\beta `$ from Bob to Alice, and then teleporting back the transformed state $`U|\psi _\beta `$ from Alice to Bob.
We begin by noting that linearity forces the transformed state of systems $`\alpha AB`$ to be independent of the particular input state $`|\psi _\beta `$. In addition, linearity imposes the condition that the state $`|\mathrm{\Phi }(U,\chi )_{\alpha AB}`$ has to be independent of U itself. To see this, consider the case where the transformation $`U`$ is one of the four Pauli operators $`\sigma ^\mu `$ and assume that the global state of $`\alpha AB`$ after completing the protocol may depend on the choice of $`U`$. According to Eq. (7), the combined action of the operations $`G_i`$ has to be such that
$$G_2\sigma ^\mu G_1\left(|\chi _{\alpha AB}|\psi _\beta \right)=|\mathrm{\Phi }(\sigma ^\mu ,\chi )_{\alpha AB}(\sigma ^\mu |\psi _\beta ).$$
(8)
On the other hand, an arbitrary one-qubit unitary transformation $`U`$ can always be decomposed in terms of the Pauli operators, $`U=_{\mu =0}^4\alpha _\mu \sigma ^\mu `$, and it must hold that
$`G_2UG_1\left(|\chi _{\alpha AB}|\psi _\beta \right)`$ $`=`$ $`{\displaystyle \underset{\mu }{}}\alpha _\mu |\mathrm{\Phi }(\sigma ^\mu ,\chi )_{\alpha AB}(\sigma ^\mu |\psi _\beta ).`$ (9)
For the RHS to be a product state, as is required by Eq. (7), we must have $`|\mathrm{\Phi }(\sigma ^\mu ,\chi )_{\alpha AB}=|\mathrm{\Phi }(\chi )_{\alpha AB}`$, independent of the operator $`\sigma ^\mu `$. This is true for any basis set of operators, and so the final state of the ancillas $`\alpha AB`$ on the RHS of Eq. (7) is independent of $`U`$.
We can now show that the operation $`G_1`$ necessarily has to be non-trivial. We do this by first assuming the contrary that $`G_1=𝟙`$, and considering two input states, $`|\psi _\beta `$ and $`|\psi ^{}_\beta `$ such that $`{}_{\beta }{}^{}\psi ^{^{}}|\psi _{\beta }^{}=0`$, and two unitary transformations $`\mathrm{U}`$ and $`\mathrm{U}^{}`$ which bring these two states to the same state $`|\gamma _\beta `$. Using Eq. (7), this implies that
$`G_2\left(U|\chi _{\alpha AB}|\psi _\beta \right)`$ $`=`$ $`|\mathrm{\Phi }(\chi )_{\alpha AB}|\gamma _\beta `$ (10)
$`G_2\left(U^{}|\chi _{\alpha AB}|\psi ^{}_\beta \right)`$ $`=`$ $`|\mathrm{\Phi }(\chi )_{\alpha AB}|\gamma _\beta .`$ (11)
No universal unitary action $`G_2`$ can be found to satisfy Eq. (11), as this would require the mapping of orthogonal states onto the same state. This shows that no universal operation $`G_2`$ that satisfies Eq. (11) can exist and therefore, for the $`U`$-teleportation to succeed, $`G_1`$ has to be non-trivial.
The final step in our proof is to rewrite Eq. (7) as
$$UG_1(|\chi _{\alpha AB}|\psi _\beta )=G_2^{}(|\mathrm{\Phi }(\chi )_{\alpha AB}U|\psi _\beta ).$$
(12)
Since $`G_1`$ and $`G_2`$ are universal gates, we may choose $`U`$ and $`|\psi _\beta `$ freely. For each $`|\psi _\beta `$ let the operator $`U_\psi `$ be such that $`U_\psi |\psi =|\mathrm{\hspace{0.17em}0}`$ where $`\sigma _z|\mathrm{\hspace{0.17em}0}=|\mathrm{\hspace{0.17em}0}`$. If $`U=\sigma _zU_\psi `$, then
$`\left(\sigma _zU_\psi \right)G_1\left(|\chi _{\alpha AB}|\psi _\beta \right)`$ $`=`$ $`G_2^{}\left(|\mathrm{\Phi }(\chi )_{\alpha AB}\sigma _zU_\psi |\psi _\beta \right)`$
$`=`$ $`G_2^{}\left(|\mathrm{\Phi }(\chi )_{\alpha AB}|\mathrm{\hspace{0.17em}0}_\beta \right).`$
The RHS is simply $`(U_\psi )G_1\left(|\chi _{\alpha AB}|\psi _\beta \right)`$ and so, necessarily, $`(U_\psi )G_1\left(|\chi _{\alpha AB}|\psi _\beta \right)`$ is the eigenstate $`|\mathrm{\hspace{0.17em}0}_\alpha |\varphi _{AB\beta }`$ of $`\left(\sigma _z\right)_\alpha 𝟙_{𝔸𝔹\beta }`$. Equivalently,
$`G_1\left(|\chi _{\alpha AB}|\psi _\beta \right)`$ $`=`$ $`\left(U_\psi ^{}|\mathrm{\hspace{0.17em}0}_\alpha \right)|\varphi _{AB\beta }`$ (13)
$`=`$ $`|\psi _\alpha |\varphi _{AB\beta }.`$ (14)
In other words, the operation $`G_1`$ necessarily transfers Bob’s state $`|\psi `$ to Alice’s qubit $`\alpha `$. Substituting Eq. (14) into Eq. (7) then shows that $`G_2`$ necessarily transfers $`U|\psi `$ back to Bob’s qubit $`\beta `$. From these results and the fact that quantum state teleportation is an optimal procedure for local state transfer, we conclude that the optimal procedure for implementing locally a universal U-teleportation scheme is by means of bidirectional state teleportation.
In this paper we have investigated the potential use of LOCCSE for the remote control of a quantum system. We have determined requirements that must be satisfied by any method that implements this task by LOCCSE means. In particular, we have shown that, if Alice can teleport an arbitrary unitary operation to a qubit in her colleague Bob’s laboratory, then she must communicate at least two bits of classical information to him, and they must share at least 2 ebits of entanglement. If the unitary operation is remotely implemented by a classical apparatus, then to effect the teleportation at least 2 classical bits must also be transmitted from Bob to Alice. These resources can be used to perform the teleportation of U using bidirectional state teleportation. Remarkably, no protocol employing a smaller amount of resources is possible.
We believe that this work will stimulate further research into ways in which LOCCSE can be used to control remotely the properties of other quantum system, with potential applications ranging from remotely synchronized time evolutions to distributed quantum computing.
Acknowledgements: The authors thank O. Steuernagel and S. M. Barnett for discussions and D. Jonathan and S. Virmani for critically reading the manuscript. This work has been supported by The Leverhulme Trust, the EQUIP project of the European Union, the Engineering and Physical Sciences Research Council (EPSRC) and DGICYT Project No. PB-98-0191 (Spain).
$`()`$ Permanent Address: Departamento de Física. Universidad de Oviedo. Calvo Sotelo s/n. 33007 Oviedo. Spain.
Figure 1. Caption.
Initial setup involved in the teleportation of an arbitrary unitary operation. The control system C in Alice’s laboratory is initially prepared in the state $`|U_C`$, corresponding to the unitary operation U. This operation is to be remotely carried out on Bob’s qubit $`\beta `$, which is initially prepared in an arbitrary pure state $`|\psi _\beta `$. This will be achieved by local operations in the individual laboratories, involving a collective ancilla initially prepared in the state $`|\chi _{AB}`$, supplemented by the exchange of classical communication, represented in the diagram by the arrow lines.
Figure 2. Caption.
Quantum circuit representation of the process of teleporting an arbitrary one-qubit transformation. The two upper wires belong to Alice and the lower ones to Bob. Initially Alice and Bob share some entanglement, represented by the joint state $`|\chi _{\alpha AB}`$. Operations $`G_1`$ and $`G_2`$ are modeled in terms of non-local unitary transformations.
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# Description of Multi Quasi Particle Bands by the Tilted Axis Cranking Model
## I Introduction
Since its introduction in ref. , the Tilted Axis Cranking (TAC) approach has turned out to be quite successful in describing $`\mathrm{\Delta }I=1`$ rotational bands . In particular it has led to the understanding of the appearance of regular magnetic dipole bands in nearly spherical nuclei . The physical aspects of these investigations have been reviewed in . Though different aspects of the actual calculations were touched in these papers, a comprehensive presentation of the theory, the calculational methods, and the practical application of the TAC approach is still missing. In the present paper we provide it by using the rotational bands in the nuclides with $`N=102,103`$ and $`Z=71,72,73`$ as examples. The TAC approach has been applied so far only for potentials of the Nilsson type, which are combined with a pairing plus quadrupole model interaction or the shell correction method for finding the deformation. A systematic exposure of the applied techniques, the experiences gathered as well as the successes and limitations of the TAC approach as it stands seems to be timely for two reasons. On the one hand hand it is meant as guideline for application of the existing program system, which has turned out quite useful in the data analysis. One the other hand it may serve as a starting point for more sophisticated versions of the TAC mean field approximation, as the up to date versions of the Hartree-Fock approximation or the Relativistic Mean Field approach.
The earliest invocations of cranking about a non-principal axis were in the context of wobbling motion . Kerman and Onishi pointed out the possibility of uniform rotation about a non-principal axis. Frisk and Bengtsson demonstrated the existence of such solutions for realistic nuclei and discussed conditions where to expect them . Goodman demonstrated that the moments of inertia may strongly depend on the orientation of the rotational axis, which implies the possibility of uniform rotation about a tilted axis. However, these studies did not give the physical interpretation of the TAC solutions and left open the question if taking into account the self-consistency with respect to the shape degrees of freedom would not result in rotation about a principal axis. In fact, the investigation of the rotating harmonic oscillator by Cuypers and a few level model by Nazarewicz and Szymanski seemed to support the latter possibility. Frauendorf found the first fully self-consistent TAC solutions and gave their interpretation in terms of $`\mathrm{\Delta }I=1`$ rotational bands. This marked the origin of the fully fledged Tilted Axis Cranking (TAC) approach.
Marshalek studied the possibility of tilted rotation generated by superpositions of collective vibrations, while Alhassid and Bush , Goodman , and Dodaro and Goodman included the tilt of the rotational axis in their analysis of nuclei at nonzero temperature. A recent reinvestigation of the rotating harmonic oscillator by Heiss and Nazmitdinov claims the existence of TAC solutions within this model, in contrast to . Horibata and Onishi , Horibata et al. and Dönau et al. have started to investigate the dynamics of the orientation angles in the frame of the Generator Coordinate Method.
Section II presents the relevant expressions for the energies and electro-magnetic transition matrix elements. It discusses the interpretation of the cranking solutions, important technical aspects and approximations that help to find the solution of the TAC equations in an efficient way. It investigates the relation of TAC to the treatment of $`\mathrm{\Delta }I=1`$ bands in the framework of the standard Principal Axis Cranking (PAC) approach. It explains how to read the quasi particle diagrams. Section III analyses the rotational bands in the yrast region of the nuclides with $`N=102,103`$ and $`Z=71,72,73`$. The main purpose is to illustrate how to construct the multi quasi particle configurations and how to relate them to the experimental rotational bands. Merits and limitations of the method will be exposed and compared with the standard Cranking approach. We are not going to optimize all parameters of the mean field for each configuration. In the spirit of the Cranked Shell Model approach only semi quantitative agreement with the data is sought, the focus being the qualitative structure of the band spectrum. Well deformed nuclei are taken as examples because the assumption of one and the same deformation for the various quasi particle configurations is realistic. The specific nucleon numbers are chosen because a large number of high $`K`$ bands and low $`K`$ bands have been found in these nuclides. This makes them an appropriate test ground for the TAC model. This paper is restricted to the HFB approximation for pairing. A more sophisticated version of TAC based on particle number projection will be presented separately . Since the change of the pair field is not in the concern of this paper but rather an unwanted complication, the self-consistency of this degree of freedom is treated in a schematic way.
Section IV provides a set of rules for the analysis of rotational bands in terms TAC, which is meant as a reference for potential users of the method. The conclusion are given in section V.
## II Tilted axis cranking
### A General layout
Two versions of the TAC have been developed and applied
i) The pairing plus quadrupole model (PQTAC)
ii) Shell correction method (SCTAC)
In the subsections II B \- II I we present the PQTAC in detail. Section II E describes the differences between SCTAC and PQTAC. The PQTAC it more appropriate for small deformations, whereas SCTAC is better suited for large deformation. Subsection II J discusses the schematic treatment of pairing and subsection II K explains the how to read the quasi particle diagrams.
### B The pairing plus quadrupole model (PQTAC)
We assume that the rotational axis is the z-axis and start with the two-body Routhian
$$H^{}=H\omega \widehat{J}_z.$$
(1)
It consists of the rotationally invariant two-body Hamiltonian $`H`$ and the constraint $`\omega \widehat{J}_z`$ which ensures that the low-lying states have a finite angular momentum projection $`J=<\widehat{J}_z>`$. As a two body Hamiltonian, the pairing plus quadrupole interaction is used,
$$H=H_{sph}\frac{\chi }{2}\underset{\mu =2}{\overset{2}{}}Q_\mu ^+Q_\mu GP^+P\lambda N.$$
(2)
The model and its justification are described in the textbooks (see, for example, Ring and Schuck ). We use a slightly modified version, which is constructed such that the derived mean field Hamiltonian coincides with the popular Nilsson Hamiltonian (see, for example, Ring and Schuck and Nilsson and Ragnarsson ). The motivation is that the parameters of the Nilsson Hamiltonian are carefully adjusted and that it is useful to have the same standard mean field for nuclei with large deformation, where the shell correction method is more appropriate. Thus, the spherical part
$$H_{sph}=\underset{k}{}\epsilon _kc_k^+c_k$$
(3)
is parameterized in the same way as the Nilsson Hamiltonian. For the calculations in this paper we use the set of parameters given in .
The pairing interaction is defined by the monopole pair field
$$P^+=\underset{k>0}{}c_k^+c_{\overline{k}}^+.$$
(4)
Here $`\overline{k}`$ is the time reversed state of $`k`$. The quadrupole interaction is defined by the operators <sup>*</sup><sup>*</sup>*This definition of the quadrupole operators corresponds to $`Q_0=r^2P_2(\mathrm{cos}\vartheta )`$, with $`P_2`$ being the Legendre polynomial.
$$Q_\mu =\underset{kk^{}}{}\sqrt{\frac{4\pi }{5}}<k|r^2Y_{2\mu }|k^{}>c_k^+c_k^{}.$$
(5)
In order to simplify the notation all expressions are written only for one kind of particles. They are understood as sums of a proton and a neutron part.
The wave function is approximated by the Hartree - Fock - Bogoljubov (HFB) mean field expression $`|>`$. Neglecting exchange terms, the HFB - Routhian becomes
$$h^{}=h_{sph}\underset{\mu =2}{\overset{2}{}}q_\mu (Q_\mu ^++Q_\mu )\mathrm{\Delta }(P^++P)\lambda N\omega \widehat{J}_z.$$
(6)
The self-consistency equations determine the deformed part of the potential
$$q_\mu =\chi <Q_\mu >$$
(7)
and the pair potential
$$\mathrm{\Delta }=G<P>.$$
(8)
The chemical potential $`\lambda `$ is fixed by the standard condition
$$N=<\widehat{N}>.$$
(9)
The quasi particle operators
$$\alpha _i^+=\underset{k}{}(U_{ki}c_k^++V_{ki}c_k)$$
(10)
obey the equations of motion
$$[h^{},\alpha _i^+]=e_i^{}\alpha _i^+,$$
(11)
which define the well known HFB eigenvalue equations for the quasi particle amplitudes $`U_{ki}`$ and $`V_{ki}`$. The explicit form of these equations can be found in . They have the familiar symmetry under particle hole conjugation, which has the consequence that for each quasi particle solution $`i`$ there is a conjugate $`i^+`$ with
$$e_{i^+}^{}=e_i^{},U_{ki^+}=V_{ki},V_{ki^+}=U_{ki}.$$
(12)
The quasi particles have good parity, but in general no good signature. The consequences of good or broken signature will be discussed in subsection II C.
The quasi particle operators refer to the vacuum state $`|0>`$, which is defined by the condition
$$\alpha _i|\mathrm{\hspace{0.17em}0}>=0i.$$
(13)
They define the excited quasi particle configurations
$$|i_1,i_2,\mathrm{}>=\alpha _{i_1}^+\alpha _{i_2}^+\mathrm{}|0>.$$
(14)
The rules and strategies for constructing quasi particle configurations from them will be be discussed below by means of concrete examples.
The set of HFB eq. (6) - (13) can be solved for any configuration $`|i_1,i_2\mathrm{}>`$. For the self-consistent solution, the total Routhian
$$E^{}=<H^{}>$$
(15)
has an extremum
$$\frac{E^{}}{_{q_\mu }}|_\omega =0,\frac{E^{}}{\mathrm{\Delta }}|_\omega =0.$$
(16)
The total energy as function of the angular momentum
$$J=<\widehat{J}_z>$$
(17)
is given by
$$E(J)=E^{}(\omega )+\omega J(\omega )$$
(18)
where the eq. (17) implicitly fixes $`\omega (J)`$. The total energy is also extremal for the self-consistent solution
$$\frac{E}{_{q_\mu }}|_J=0,\frac{E}{\mathrm{\Delta }}|_J=0,$$
(19)
where the derivatives have to be taken at a fixed value of $`J`$. For a family of self-consistent solutions $`|\omega >`$, found for different values of $`\omega `$, there hold the canonical relations
$$\frac{dE^{}}{d\omega }=J,\frac{dE}{dJ}=\omega .$$
(20)
In ref. it was shown that for a self-consistent solution the vector of the angular velocity
$$\stackrel{}{\omega }=(\omega _x,\omega _y,\omega _z)=(0,0,\omega )$$
(21)
and the vector of the expectation values of the angular momentum components
$$\stackrel{}{J}=(\widehat{J}_x,\widehat{J}_y,\widehat{J}_z)$$
(22)
must be parallel. The argument is as follows. Since the interaction is rotational invariant, one has
$$[H^{},\widehat{J}_x]=i\omega \widehat{J}_y,$$
(23)
$$[H^{},\widehat{J}_y]=i\omega \widehat{J}_x.$$
(24)
Since the left-hand sides are small variations of $`E^{}`$, the stationarity of $`E^{}`$ implies
$$<\widehat{J}_x>=<\widehat{J}_y>=0,$$
(25)
i. e.
$$\stackrel{}{\omega }||\stackrel{}{J}.$$
(26)
This holds also in the intrinsic frame of reference, which will be discussed in sect. II D.
### C Tilted Solutions
The formalism presented above is the well known ”Cranking Model” as laid out in the textbooks . The ”Tilted-Axis Cranking” (TAC) version accounts for the possibility that the principal axes (PA) of the quadrupole tensor $`q_\mu `$ need not to coincide with the rotational axis (z). Hence, one has to distinguish between two possibilities:
* Principal Axis Cranking (PAC). The z - axis (rotational or cranking axis) coincides with one of the PA. Then, the signature $`r`$ is a good quantum number, i. e.
$$e^{i\pi \widehat{J}_z}|\pi ,\alpha ,\omega >=r|\pi ,\alpha ,\omega >.$$
(27)
Following we indicate the signature $`r=e^{i\pi \alpha }`$ by the signature exponent $`\alpha `$. The quasi particle configuration $`|\pi ,\alpha ,\omega >`$ describes a $`\mathrm{\Delta }I=2`$ rotational band, the spins of which take the values
$$I=\alpha +evennumber.$$
(28)
* Tilted Axis Cranking (TAC). The z - axis does not coincides with one of the PA, i. e. it is tilted away from the PA. Then,
$$e^{i\pi \widehat{J}_z}|\pi ,\omega >e^{i\pi \alpha }|\pi \omega >.$$
(29)
The signature is no longer a good quantum number. The quasi particle configuration $`|\pi ,\omega >`$ describes a $`\mathrm{\Delta }I=1`$ rotational band of given parity.
The different interpretation of solutions with different symmetry is characteristic for spontaneous symmetry breaking in the mean field approximation. It makes it necessary to eliminate spurious configurations and will lead to discontinuities when the symmetry changes as a function of the frequency $`\omega `$. These problems which will be discussed in the subsections III B by means of concrete examples.
### D Intrinsic coordinates
It is useful to reformulate the TAC approach in the frame of the PA of the quadrupole tensor. This ”intrinsic” coordinate system is defined by demanding that the components of the quadrupole tensor satisfy the conditions
$$q_1^{}=q_1^{}=0,q_2^{}=q_2^{}.$$
(30)
The orientation of the PA, which are denoted by 1, 2, and 3, with respect to the lab frame is fixed by the three Euler angles $`\psi `$, $`\vartheta `$ and $`\phi `$, the meaning of which is illustrated in fig. 1. The two intrinsic quadrupole moments $`q_0^{}`$ and $`q_2^{}`$ specify the deformation of the potential. The quadrupole moments in the lab frame are related to them by
$$q_\mu =D_{\mu 0}^2(\psi ,\vartheta ,\phi )q_0^{}+(D_{\mu 2}^2(\psi ,\vartheta ,\phi )+D_{\mu 2}^2(\psi ,\vartheta ,\phi ))q_2^{},$$
(31)
where $`D_{\nu \mu }^2(\psi ,\vartheta ,\phi )`$ are the Wigner $`D`$-functions The convention of is used..
The different angles $`\psi `$ correspond to one and the same intrinsic state. They are degenerate. We choose the one with $`\psi =0`$. We restrict the consideration to $`planar`$ TAC solutions, which is the case when the z-axis lies in one of the three principal planes defined by the PA. We assume that it lies in the 1 - 3 plane, i. e. choose $`\phi =0`$. In the case of axial shapes, this is one choice from the equivalent solutions differing by the angle $`\phi `$. For triaxial shapes one may always relabel the PA by means of the shape parameterization, letting the triaxiality parameter $`\gamma `$ vary within an interval of 180<sup>o</sup> . Which axes of the triaxial potential lie in the 1-3 plane can be found in table I. It is seen that all three possibilities appear in the half-plane $`60^o\gamma 120^o`$. The other half-plane is a repetition with the axes 1 and 3 exchanged.
With the above mentioned restrictions and conventions the deformed potential is fixed by the two intrinsic quadrupole moments $`q_0^{}`$ and $`q_2^{}`$ and the orientation (”tilt”) angle $`\vartheta `$ between the 3 - and the z - axis, which is the direction of the rotational axis. In the intrinsic frame the HFB Routhian reads
$`h^{}=h_{sph}q_0^{}Q_0^{}q_2^{}(Q_2^{}+Q_2^{})\mathrm{\Delta }(P^++P)`$ (32)
$`\lambda N\omega (\mathrm{sin}\vartheta J_1+\mathrm{cos}\vartheta J_3).`$ (33)
Figs. 2 and 3 show examples of the the quasi particle levels $`e_i^{}`$ as functions of the rotational frequency $`\omega `$ and the orientation angle $`\vartheta `$.
The shape is fixed by the two equations
$$q_0^{}=\kappa <Q_0^{}>,q_2^{}=\kappa <Q_2^{}>$$
(34)
and the orientation angle $`\vartheta `$ by the condition that the expectation value of the angular momentum and the angular velocity must have the same direction, i. e.
$$\stackrel{}{J}=(<\widehat{J}_1>,\mathrm{\hspace{0.17em}0},<\widehat{J}_3>)||\stackrel{}{\omega }=(\omega \mathrm{sin}\vartheta ,\mathrm{\hspace{0.17em}0},\omega \mathrm{cos}\vartheta ),$$
(35)
respectively. These parameters correspond to extrema of total Routhian, that is
$$\frac{E^{}}{_{q_0^{}}}|_\omega =0,\frac{E^{}}{_{q_2^{}}}|_\omega =0,\frac{E^{}}{\vartheta }|_\omega =0.$$
(36)
Of course only the minima are interpreted as bands.
In praxis it is convenient to solve the equation (35) for each combination of $`q_0^{}`$ and $`q_2^{}`$ which is needed to obtain the shape from the equations (34) with the desired accuracy. Very often it is enough to determine the shape for one value of $`\omega `$ and then keep it fixed for other values, only calculating the orientation angle $`\vartheta `$ by means of the condition (35).
Using the Cartesian representation of the quadrupole moments, the HFB Routhian (21) becomes the modified oscillator potential
$`h^{}={\displaystyle \underset{\nu =1}{\overset{3}{}}}{\displaystyle \frac{(p_\nu ^2+\omega _\nu ^2x_\nu ^2)}{2M}}+\kappa l_\nu s_\nu +\mu (l_\nu ^2<l_\nu ^2>)`$ (37)
$`\mathrm{\Delta }(P^++P)\lambda N\omega (sin\vartheta J_1+\mathrm{cos}\vartheta J_3),`$ (38)
where the oscillator frequencies are parameterized by means of Nilsson’s deformation parameters $`\delta `$ and $`\gamma `$ (cf. e. g. ),
$$\omega _\nu ^2=\omega _{00}^2\left[1\frac{2}{3}\delta \mathrm{cos}(\gamma \frac{2\pi \nu }{3})\right].$$
(39)
The only difference to the standard modified oscillator model is that there is no volume conservation in the pairing plus quadrupole model. Since we are only interested in small deformation the coupling between the oscillator shells is not taken into account when diagonalizing the HFB Routhian (37). Solving the self-consistency equation (34) and calculating the total energy, the coupling between the oscillator shells is also neglected.
### E Strutinsky Renormalization (SCTAC)
An alternative version of TAC starts with the modified oscillator Routhian (37). As e. g. described in , stretched coordinates are introduced and the matrix elements $`<N|J_\mu |N\pm 2>`$ are neglected in the stretched basis. This is a standard procedure which takes into account most of the couplings between the oscillator shells. The oscillator frequencies are parameterized by means of Nilsson’s alternative set of deformation parameters $`\epsilon `$ and $`\gamma `$ ,
$$\omega _\nu =\omega _0\left[1\frac{2}{3}\epsilon \mathrm{cos}(\gamma \frac{2\pi \nu }{3})\right],$$
(40)
where the condition of volume conservation $`\omega _0^3=\omega _1\omega _2\omega _3`$ fixes $`\omega _0`$. The total Routhian is obtained by applying the Strutinsky renormalization to the energy of the non-rotating system $`E_0`$. This kind of approach has turned out to be a quite reliable calculation scheme in the case of standard PAC . One minimizes the total RouthianFor the treatment of the term $`\lambda <\widehat{N}>`$ see sect. II J.
$`E^{}(\omega ,\vartheta ,\epsilon ,\epsilon _4,\gamma ,\mathrm{\Delta },\lambda )=E_{LD}(\epsilon ,\epsilon _4,\gamma )\stackrel{~}{E}(\epsilon ,\epsilon _4,\gamma )`$ (41)
$`+<h^{}>+(2\mathrm{\Delta }G<P>)<P>,`$ (42)
where $`|>=|\omega ,\vartheta ,\epsilon ,\epsilon _4,\gamma ,\mathrm{\Delta },\lambda >`$ is a quasi particle configuration belonging to the mean field Routhian $`h^{}(\omega ,\vartheta ,\epsilon ,\epsilon _4,\gamma ,\mathrm{\Delta },\lambda )`$ as defined above. The smooth energy $`\stackrel{~}{E}`$ is calculated from the single-particle energies, which are the eigenvalues of $`h^{}(\omega =0,\vartheta =0,\epsilon ,\epsilon _4,\gamma ,\mathrm{\Delta }=0)`$ by means of Strutinsky averaging . The expressions for the liquid drop energy $`E_{LD}(\epsilon ,\epsilon _4,\gamma )`$ are given, for example, in , where also the averaging procedure is described. For given $`\epsilon ,\epsilon _4`$ and $`\gamma `$, the tilt angle is determined by means of the condition (35). Then, the minimum of $`E^{}(\omega ,\epsilon ,\epsilon _4,\gamma )`$ with respect to the deformation parameters is found. Since $`|\omega >`$ is an eigenfunction of $`h^{}(\omega ,\vartheta )`$ the Routhian $`<\omega ,\vartheta |h^{}|\omega ,\vartheta >`$ is stationary at the angle where the condition (35) is fulfilled and so is $`E^{}`$ because the other terms do not depend on $`\vartheta `$. Hence, the procedure determines a stationary point with respect to the mean field parameters and the canonical relations (20) are satisfied.
The SCTAC approach is preferred to the PQTAC version for well deformed nuclei, because it is a reliable standard method for determining large deformations. In the calculations of well deformed nuclei it is usually a good approximation to keep the deformations fixed within a rotational band. However this is a matter of the needed accuracy and of how much effort one is willing to invest.
### F Electro - magnetic matrix elements
The intra band M1 - transition matrix element is calculated by means of the semiclassical expression
$`<I1I1|_1(M1)|II>`$ $`=<_1(M1)>=`$ (43)
$`=\sqrt{{\displaystyle \frac{3}{8\pi }}}\left[\mu _3\mathrm{sin}\vartheta \mu _1\mathrm{cos}\vartheta \right].`$ (44)
The components of the transition operator $`_\nu `$ refer the the lab system. The expectation value is taken with the TAC configuration $`|>`$. In the second line $`_\nu `$ is expressed by the components of the magnetic moment in the intrinsic frame. The reduced M1-transition probability becomes
$$B(M1,II1)=<_1(M1)>^2.$$
(45)
The spectroscopic magnetic moment is given by
$`\mu =<II|\mu _z|II>={\displaystyle \frac{I}{I+1/2}}<\mu _z>`$ (46)
$`={\displaystyle \frac{I}{I+1/2}}[\mu _1\mathrm{sin}\vartheta +\mu _3\mathrm{cos}\vartheta ].`$ (47)
The factor $`\frac{I}{I+1/2}`$ is a quantal correction which is close to one for high spin. The components of the magnetic moment with respect to the PA are calculated by means of
$`\mu _1=\mu _N(J_{1,p}+(\eta 5.581)S_{1,p}\eta 3.82S_{1,n}),`$ (48)
$`\mu _3=\mu _N(J_{3,p}+(\eta 5.581)S_{3,p}\eta 3.82S_{3,n}),`$ (49)
where the components of the vectors of angular momentum $`\stackrel{}{J}`$ and of the spin $`\stackrel{}{S}=<\stackrel{}{s}>`$ are the expectation values with the TAC configuration $`|>`$. The free spin magnetic moments are attenuated by a factor $`\eta =0.7`$. For mass other mass regions a somewhat different attenuation may be taken, which needs not be the same for protons and neutrons.
The intra band E2 - transition matrix elements are calculated by means of the semiclassical expressions<sup>§</sup><sup>§</sup>§Refs. contain some unfortunate inconsistency between the quadrupole moments of the quadrupole interaction and the electric transition matrix elements. These concern only the written formulae, the results of the calculations quoted are correct and consistent with the ones given here.
$`<I2I2|_2(E2)|II>=<_2(E2)>=`$ (50)
$`=\sqrt{{\displaystyle \frac{5}{4\pi }}}\left({\displaystyle \frac{eZ}{A}}\right)[\sqrt{{\displaystyle \frac{3}{8}}}<Q_0^{}>(\mathrm{sin}\vartheta )^2`$ (51)
$`+{\displaystyle \frac{1}{4}}<Q_2^{}+Q_2^{}>(1+(\mathrm{cos}\vartheta )^2)],`$ (52)
$`<I1I1|_1(E2)|II>=<_1(E2)>=`$ (53)
$`=\sqrt{{\displaystyle \frac{5}{4\pi }}}\left({\displaystyle \frac{eZ}{A}}\right)[\mathrm{sin}\vartheta \mathrm{cos}\vartheta (\sqrt{{\displaystyle \frac{3}{2}}}<Q_0^{}>`$ (54)
$`{\displaystyle \frac{1}{2}}<Q_2^{}+Q_2^{}>)],`$ (55)
and the spectroscopic quadrupole moment by
$`Q=<II|Q_0^{(BM)}|II>={\displaystyle \frac{I}{I+2/3}}<Q_0^{(BM)}>=`$ (56)
$`{\displaystyle \frac{I}{I+2/3}}{\displaystyle \frac{2eZ}{A}}[<Q_0^{}>((\mathrm{cos}\vartheta )^2{\displaystyle \frac{1}{2}}(\mathrm{sin}\vartheta )^2)`$ (57)
$`+\sqrt{{\displaystyle \frac{3}{8}}}<Q_2^{}+Q_2^{}>(\mathrm{sin}\vartheta )^2].`$ (58)
We use the conventional definition of the static quadrupole moment as given in ref. , which differs by a factor of 2 from our quadrupole moments in the lab frame. There is a similar quantal correction factor as for the magnetic moment.
The reduced E2-transition probabilities are
$$B(E2,II2)=<_2(E2)>^2$$
(59)
and
$$B(E2,II1)=<_1(E2)>^2.$$
(60)
The mixing ratio is
$$\delta =\frac{<_1(E2)>}{<_1(M1)>}.$$
(61)
The mass quadrupole moments consist two terms. The first one contains the microscopic expectation values $`<Q_0^{}>_N`$ and $`<Q_{\pm 2}^{}>_N`$, where the subscript $`N`$ indicates that only the $`\mathrm{\Delta }N=0`$ matrix elements of the quadrupole operator are taken. The second term takes care of the coupling between the oscillator shells.
In the case of SCTAC the stretched coordinates are introduced to approximately take the coupling between the oscillator shells into account. The expectation values needed in eqs. (59, 60, 56), are the quadrupole moments in unstretched coordinates, which are given by
$`<Q_0^{}>=`$ (62)
$`{\displaystyle \frac{1}{6}}({\displaystyle \frac{2\omega _o}{\omega _3}}{\displaystyle \frac{\omega _o}{\omega _1}}{\displaystyle \frac{\omega _o}{\omega _2}})<r^2>_N`$ (63)
$`+{\displaystyle \frac{1}{6}}({\displaystyle \frac{4\omega _o}{\omega _3}}+{\displaystyle \frac{\omega _o}{\omega _1}}+{\displaystyle \frac{\omega _o}{\omega _2}})<Q_0^{}>_N`$ (64)
$`{\displaystyle \frac{1}{2\sqrt{6}}}({\displaystyle \frac{\omega _o}{\omega _1}}{\displaystyle \frac{\omega _o}{\omega _2}})<Q_2^{}+Q_2^{}>_N`$ (65)
$`<Q_2^{}+Q_2^{}>=`$ (66)
$`{\displaystyle \frac{1}{\sqrt{6}}}({\displaystyle \frac{\omega _o}{\omega _1}}{\displaystyle \frac{\omega _o}{\omega _2}})<r^2>_N`$ (67)
$`{\displaystyle \frac{1}{\sqrt{6}}}({\displaystyle \frac{\omega _o}{\omega _1}}{\displaystyle \frac{\omega _o}{\omega _2}})<Q_0^{}>_N`$ (68)
$`+{\displaystyle \frac{1}{2}}({\displaystyle \frac{\omega _o}{\omega _1}}+{\displaystyle \frac{\omega _o}{\omega _2}})<Q_2^{}+Q_2^{}>_N`$ (69)
where the semiclassical value $`<r^2>_N1.2ZA^{2/3}fm`$ is used.
In the case of the PQTAC the coupling between the oscillator shells is neglected. This is a reasonable approximation for the rotational response of the valence particles. However, when calculating electric quadrupole moments it cannot not be neglected, because it accounts for the polarization of the core by the valence nucleons. We describe the polarization by means of eqs (65,66), setting $`\epsilon =\delta `$. This prescription satisfies the consistency condition that the deformations of the potential and the density should be the same . It corresponds to a polarization charge close to 1, as estimated for the isoscalar quadrupole mode . This choice of the polarization charge makes to PQTAC and the SCTAC as similar as possible.
In the above described methods one could also use the proton part of the quadrupole moments instead of $`Z/A`$ times the mass quadrupole moments.
### G Quantization
Due to leading quantal correction (cf. e. g. ) one must associate the total angular momentum $`J`$ calculated in TAC with $`I+1/2`$, where $`I`$ is the quantum number of the angular momentum. This prescription permits us to compare the TAC calculations with the experimental energies and the static moments. Genuine TAC solutions $`(\vartheta 0^o,\mathrm{\hspace{0.17em}90}^o)`$ represent $`\mathrm{\Delta }I=1`$ bands. In this case, the experimental rotational frequency $`\omega `$ is introduced by
$$J=I,\omega (II1)=E(I)E(I1),$$
(70)
and the experimental Routhian by
$$E^{}(II1)=\frac{1}{2}[E(I)+E(I1)]\omega (II1)J.$$
(71)
Here, the canonical relations (20) are approximated by quotients of finite differences. The data define a discrete sets of points $`J(\omega )`$ and $`E^{}(\omega )`$, which are connected by interpolation. If the axis of rotation coincides with one of the principal axes $`(\vartheta =0^o,\mathrm{\hspace{0.17em}90}^o)`$, states differing by two units of angular momentum arrange into a $`\mathrm{\Delta }I=2`$ band of given signature $`\alpha `$. In this case the frequency is calculated by
$$J=I1/2,\omega (II2)=\frac{1}{2}[E(I)E(I2)],$$
(72)
and the experimental Routhian by
$$E^{}(II2)=\frac{1}{2}[E(I)+E(I2)]\omega (II2)J.$$
(73)
For the transition probabilities, $`J`$ is associated with the mean value of $`I+1/2`$ of the transition, i. e.
$`B(M1,II1)=B(M1,J=I)`$ (74)
$`B(E2,II1)=B(E2,J=I)`$ (75)
$`B(E2,II2)=B(E2,J=I1/2),`$ (76)
where the rhs denotes the result of the TAC calculation taken for the indicated value of $`J`$. Another possibility is to compare the experimental transition probabilities with the ones calculated at the experimental frequency of the transition (70,72). As long as the experimental and calculated functions $`J(\omega )`$ agree well, both ways will give about the same result.
Of course, one may also use the relations (72) and (73) for a $`\mathrm{\Delta }I=1`$ band. Then, the two signature branches will lie nearly on top of each other if the discrete points are connected by smooth interpolation. This choice has the disadvantage that the distance between the discrete frequency points is doubled. It has the advantage to give smooth curves when the splitting between the two signature branches gradually develops with increasing frequency. In such a case (72) and (73) should be used.
### H Diabatic tracing
The goal of the calculation is to describe a rotational band, which corresponds to the ”same quasi particle configuration” for a set of increasing values of $`\omega `$. This means that one should keep fixed the occupation of the quasi particle states with similar structure. Usually one band does not correspond to the same configuration if the quasi particle levels are labeled according to their energy, because the quasi particle trajectories cross each other as functions of $`\omega `$ and $`\vartheta `$. In order to find the equilibrium angle, one has to calculate the functions $`J_1(\vartheta )`$ and $`J_3(\vartheta )`$. This becomes very tedious if the configurations are assigned manually by identifying the crossings from quasi particle diagrams like figs. 2 and 3. The task is greatly facilitated by tracing the structure of the quasi particle wave functions. The calculations are run changing $`\vartheta `$ or $`\omega `$ in finite steps. For a given grid point the overlaps of each quasi particle state with all states of the previous grid point are calculated. The pair with the maximal overlap continues one quasi particle level from the previous to the present grid point. The pair with the next lower overlap continues the second quasi particle trajectory. This procedure is repeated until all quasi particle trajectories are continued. For all the single particle and quasi particle diagrams shown in this paper the grid points are connected by means of this diabatic tracing.
In a practical calculation, the configurations are assigned manually for the first grid point in a loop. The following strategy has turned out to be quite efficient: First a typical angle $`\vartheta _s`$ is chosen and the quasi particle diagram $`e^{}(\omega _i,\vartheta _s)`$ is generated. The step size $`\mathrm{\Delta }\omega =0.05MeV`$ has turned out to be a good choice. Configurations are assigned for a typical frequency. The occupation numbers for the other grid points $`\omega _i`$ are found by means of the quasi particle diagrams or, if the crossing pattern is complex, using the tracing facility of the code. These occupation numbers are used to set the configurations in a $`\vartheta `$-loop starting at $`\omega _i`$ and $`\vartheta _s`$. The configurations of the other grid points in the loop are determined by means of diabatic tracing. Then, the code finds the orientation angle $`\vartheta `$ for each $`\omega _i`$ by means of the self-consistency condition (35) and calculates the interesting quantities. The step size $`\mathrm{\Delta }\vartheta =5^o`$ has turned out to be a good choice. At which $`\vartheta _s`$ the loop is started depends on the type of the band and will be discussed below.
Problems are encountered when the quasi particle levels do not cross sharply when $`\vartheta `$ or $`\omega `$ are changing. If the grid point happens to be located in the middle of the region where the levels strongly mix and repel each other, the diabatic tracing does not always follow the desired structure. Such cases necessitate human interference in order to continue the correct structure. One reruns the calculation with the complementary configuration and puts the parts with the correct configuration together. The grid point itself is problematic because the cranking model becomes a bad approximation due to the unphysical mixing of states with different angular momentum. These problems have been investigated for the standard cranking model . We restrict ourselves to the most simple solution advocated in : We discard such grid points and bridge the crossing region by means of interpolation.
The results of the diabatic tracing depend on the step size. It should not be too small. If the step size is much smaller than the mixing region, the procedure follows the levels adiabatically, i. e. it connects the levels (of the same parity) according to their energy. On the other hand it should not be too large in order to preserve a reasonable precision. As mentioned above, step sizes of $`\mathrm{\Delta }\vartheta =5^o`$ and $`\mathrm{\Delta }\omega =0.05MeV`$ have turned out to be good choices.
For low $`K`$ bands it is usually convenient to choose $`\vartheta _s=85^o`$ for the manual assignment of the configurations. The reason is that with increasing $`\omega `$ the equilibrium angle $`\vartheta _o`$ changes quickly from zero to values close to $`90^o`$. As seen in figs. 2 and 3, the number of avoided crossings, which cause problems, is small at low frequency. Therefore the diabatic tracing works well in most cases and permits calculating the interesting range of $`\omega `$ without human interference. The configuration assignment should not be done at 90<sup>o</sup>, where the signature is good and the levels are often degenerated. The configurations discussed in subsections III D-III H are calculated by assigning configurations at $`\vartheta _s=85^o`$.
For high $`K`$ bands, $`\vartheta _o`$ remains relatively small up to rather large values of $`\omega `$. Then, starting at $`\vartheta _s=85^o`$ becomes less efficient because the number of avoided crossings increases. A smaller value of $`\vartheta _s`$ closer to the equilibrium angle is preferable. For the configurations discussed in subsections III I we used $`\vartheta _s=45^o`$. This choice has the disadvantage that one has to run the $`\vartheta `$ loop two times, for $`\vartheta <\vartheta _s`$ and $`\vartheta >\vartheta _s`$. The choice $`\vartheta _s=0^o`$ has turned out to be quite efficient in other applications of TAC to high $`K`$ bands.
Diabatic tracing is also used when the other parameters of the mean field Hamiltonian are changed in order to solve the complete set of self-consistency equations. Approaching the minimum on the multi dimensional surface $`E^{}(\omega ,\vartheta ,\epsilon ,\epsilon _4,\gamma ,\mathrm{\Delta })`$ it is applied for each step in one of the parameters.
### I Choice of the QQ coupling constant
Using the PQTAC version, the coupling constant $`\chi `$ of the QQ - interaction must be fixed. So far it has been adjusted such that the quadrupole deformations $`\epsilon `$ and $`\gamma `$ calculated for PAC solutions come as close as possible to the ones obtained by means of the shell correction method, which has a considerable predictive potential concerning the nuclear shapes (cf. e.g. ). The adjustment has been carried out for selected nuclei. The QQ coupling constant scales with $`A^{5/3}r_{osc}^4`$, where $`r_{osc}`$ is the oscillator length . This scaling has been used to determine $`\chi `$ in neighboring nuclei.
In the first TAC calculations the equilibrium shape was calculated for $`A=170`$ using the standard shell correction method at $`\omega =0`$. The calculation was repeated for PQTAC at the same deformation and $`\omega =0`$. The coupling constant $`\chi `$ was chosen such that the self-consistency equations (34) were fulfilled. This value of $`\chi `$ was kept constant in the full TAC calculation for all values of $`\omega `$. The value $`\chi =0.0174MeVr_{osc}^4`$ was found. Scaling gives $`\chi =91MeVA^{5/3}r_{osc}^4`$ for the rare earth region. When SCTAC became available, it turned out that the results of PQTAC and SCTAC were nearly identical for the nuclei around $`A=170`$.
Extrapolating by means of scaling gives $`\chi =0.0133MeVr_{osc}^4`$ for $`A=200`$. This value (scaled locally) gives a good overall description of the magnetic dipole bands in the Pb- isotopes . A deformation of $`\epsilon 0.11`$ is obtained. SCTAC gives larger deformations of $`\epsilon 0.15`$, which account less well for the data on the Pb - isotopes.
Scaling of the rare earth value gives $`\chi =0.036MeVr_{osc}^4`$ and $`0.024MeVr_{osc}^4`$ for $`A=110`$ and 140, respectively. A new adjustment of $`\chi `$ was carried out for <sup>110</sup>Cd and <sup>139</sup>Sm. The respective values $`\chi =0.036MeVr_{osc}^4`$ and $`0.022MeVr_{osc}^4`$ were determined by making equal the deformation obtained by means of PQTAC and the shell correction method for zero pairing at finite $`\omega `$. The latter values (including local scaling) gave a good description of the magnetic dipole bands in a number of nuclides of the two regions . The data on electro-magnetic transitions in <sup>105,106,108</sup>Sn , which according to the calculations have a deformation of $`\epsilon <0.11`$, seem to point to a smaller deformation than calculated. That is a smaller value of $`\chi `$ leading to smaller deformation than predicted by the shell correction method appears to be more appropriate for <sup>105,106,108</sup>Sn. This is similar to the Pb-isotopes.
The shell correction method accounts rather well for the overall tendencies of the shape, in particular for the well deformed nuclei. However, it is not obvious that for the small deformations encountered for magnetic rotation (typically $`|\epsilon |<0.11`$) the shell correction method provides a reliable gauge for $`\chi `$. In such cases it seems preferable to fix the QQ coupling constant in a different way. Since $`\chi `$ controls the quadrupole polarizability, one may adjust it to the static quadrupole moments of high spin states and the $`B(E2)`$ values of transition between them, which are particular sensitive to the quadrupole polarizability. It seems promising to use this experimental information for a fine tuning of $`\chi `$. This approach, which is discussed in more detail in the review , is being investigated .
### J Approximate treatment of self-consistency
The CHFB equations are a complex system of nonlinear equations. TAC adds a new dimension to it, the orientation angle $`\vartheta `$. The fully self-consistent solution of the equations becomes rather tedious, in particular if one tries to describe several non yrast bands. The success of the CSM shows that for a first analysis of the excitation spectrum it is often sufficient, even preferable, to keep fixed the parameters of the mean field. For selected bands, they may be determined self-consistently in subsequent calculations if one is interested in specific properties. But very often the additional effort does not pay off the gain in insight. We shall follow the CSM approach and carry out the calculations assuming that the deformation and the parameters of the pair field, $`\mathrm{\Delta }`$ and $`\lambda `$ do not depend on the rotational frequency $`\omega `$. Only the tilt angle, which usually strongly changes, is determined by means of (35) for each value of $`\omega `$.
For the well deformed nuclei considered in this paper the deformation changes turn out to be moderate. They are negligible for the more qualitative comparison with the data which we are aiming at.
The approximation of a constant pair field needs a more careful discussion. The original assumption of the CSM to keep $`\mathrm{\Delta }`$ at 80% of the experimental odd - even mass difference $`\mathrm{\Delta }_{oe}`$ becomes problematic, because modern data reach rotational frequencies where the static pair field disappears . Since the transition to the unpaired state may substantially change rotational response, self-consistency must be taken into account at least in some rough way. We found the following compromise between accuracy and effort quite satisfying. The TAC calculations are carried out at a few values of $`\mathrm{\Delta }`$, which do not depend on $`\omega `$. The total Routhians $`E^{}(\omega ,\mathrm{\Delta })`$ are plotted. At each frequency one can easily choose the best $`\mathrm{\Delta }`$-value as the one that has the lowest value of $`E^{}`$. The upper panel of fig. 4 shows as an example the yrast band of <sup>174</sup>Hf. The values $`\mathrm{\Delta }_n=0.95,0.69,0MeV`$ and $`\mathrm{\Delta }_p=1.05,0.75,0MeV`$ are used. The first point corresponds to the self-consistent ground state value and the corresponding curve is lowest for small $`\omega `$. At $`\omega =0.24MeV`$ the curve with the reduced values of $`\mathrm{\Delta }_n`$ and $`\mathrm{\Delta }_p`$ takes over. Within the considered frequency range the unpaired solution cannot compete, though the neutron correlation energy is rather small.
For paired configurations the proper Routhian is $`E^{}(\omega ,\lambda )+\lambda N`$, where $`N`$ is the exact particle number. The term $`\lambda N`$ compensates the Lagrangian multiplier introduced in (2). However, exact compensation appears only if the self-consistency condition (9) for $`\lambda `$ is fulfilled. If $`\lambda `$ is approximately determined one can proceed in a similar way as for $`\mathrm{\Delta }`$ by plotting $`E^{}(\omega ,\lambda )+\lambda N`$. Since
$$\frac{}{\lambda }E^{}(\omega ,\lambda )=<\widehat{N}>,\frac{}{\lambda }<\widehat{N}>>0,$$
(77)
the Routhian $`E^{}(\omega ,\lambda )+\lambda N`$ has a maximum at the self-consistent value of $`\lambda `$. Accordingly, TAC calculations are carried out for a few values of $`\lambda `$, which do not depend on $`\omega `$ and $`E^{}(\omega ,\lambda )+\lambda N`$ is plotted. The highest curve corresponds to the best value of $`\lambda `$. The lower panel of fig. 4 shows the three points $`\lambda _n=49.0,49.15,49.3MeV`$. The arrows indicate the frequencies where the self-consistency condition (9) is fulfilled. For $`\lambda =49.0MeV`$, the deviation in particle number is about 2 at $`\omega =0.5MeV`$. The upper envelop of the curves represents the best choice of $`\mathrm{\Delta }`$ and $`\lambda `$ within the restricted set of grid points investigated. For $`\omega >0.3MeV`$, it behaves very similar to the unpaired curve. The small correlation energy of 0.1 - 0.2 $`MeV`$ indicates weak static pairing. We will show this optimized Routhian of the yrast sequence in the figures as a reference. As long as the values of $`\lambda `$ are the same for all configurations one may leave away the term $`\lambda N`$. It is however needed to correctly calculate the relative position of configurations with different $`\lambda `$ or of paired and unpaired configurations .
This method is quite useful because it is simple and it can easily be made as accurate as needed by adding more $`\mathrm{\Delta }`$ and $`\lambda `$ values. At each stage one has a clear idea of the remaining error of the energy. The simplest variant of considering only $`\mathrm{\Delta }0.8\mathrm{\Delta }_{oe}`$ and $`0`$ and choosing $`\lambda _p`$ and $`\lambda _n`$ such that the particle numbers are right for $`\omega =0`$ turns out to be sufficient for a first orientation. It shows the pair correlation energy directly. The discussion of pairing will be restricted to this minimal variant. All figures showing total Routhians display the quantity $`E^{}+\lambda _pZ+\lambda _nN`$. In order to keep the figures simple, the ordinate is labeled with $`E^{}`$ only. The energy of the ground state of the nucleus, which is not the concern of this paper, is not calculated correctly. In all figures, only the Routhians relative to the ground state energy are of relevance.
### K Reading the quasi particle diagrams
Figs. 5 \- 8 show the single particle Routhians $`e_i^{}(\omega ,\vartheta )`$ as functions of the frequency $`\omega `$ and of the tilt angle $`\vartheta `$. In order to demonstrate change of the particle response to with the magnitude and orientation of $`\omega `$ two different frequencies are presented for each kind of particles. The figures with intermediate frequency are relevant for the present day high spin data. The high frequency figures show territory yet to be explored. The side panels of each figure are added for helping the reader to connect to new middle panel with the familiar single particle Routhians for rotation about the PA axes.
The slope of the trajectories gives negative projection of the quasi particle angular momentum on the $`\omega `$ \- axis,
$$\frac{e_i^{}}{\omega }=j_{||}=(j_{1,i}\mathrm{sin}\vartheta +j_{3,i}\mathrm{cos}\vartheta ),$$
(78)
and its perpendicular component,
$$\frac{e_i^{}}{\vartheta }=\omega j_{}=\omega (j_{1,i}\mathrm{cos}\vartheta j_{3,i}\mathrm{sin}\vartheta ),$$
(79)
where $`j_{1,i}`$ and $`j_{3,i}`$ are the expectation values of the angular momentum in the single-particle or quasi particle state $`i`$.
For $`\vartheta =0^o`$, the cranking term $`\omega \widehat{J}_3`$ commutes with the axial symmetric deformed potential and the projection of the angular momentum on the symmetry axis $`K`$ is a good quantum number. In this case, the states coincide with the non-rotating Nilsson states, which are indicated by the labels in the figure. For $`\vartheta =90^o`$, the signature $`\alpha `$ is a good quantum number which is also indicated in the figure. As the signature operator $`e^{i\pi \widehat{J}_1}`$ and $`\widehat{J}_3`$ do not commute, there is a transition from one to the other type of symmetry when $`\vartheta `$ changes from 0<sup>o</sup> to $`90^o`$.
Discussing the features of the quasi particle diagrams, we will refer to the three types of coupling schemes that appear as a consequence of the competition between the deformed potential, the inertial forces and the pair correlations. They are discussed in . Let us start with moderate frequencies $`\omega <0.3MeV`$, which are illustrated in figs. 5, and 7.
The normal parity states with $`K5/2`$ obey the deformation aligned coupling (DAL) scheme. These orbitals are strongly coupled to the deformed potential. In the $`e^{}(\vartheta )`$ plot, they are recognized as the pairs of trajectories, which branch at $`\vartheta =90^o`$. They have a small component $`j_{1,i}`$ but a large component $`j_{3,i}K`$. They approximately behave like $`K\mathrm{cos}\vartheta `$ in an extended region. Near $`\vartheta =90^o`$ there is the very narrow transition region from good signature to almost good $`K`$, where the slope changes from zero to approximately $`K\mathrm{sin}\vartheta `$. The region is too narrow to be discerned in the figure, where it looks like a kink.
There are pairs of parallel trajectories in the $`e^{}(\vartheta )`$ plot, originating from the 1/2 neutron and 1/2 proton orbitals. These are pseudo spin singlets. A discussion of the pseudo spin symmetry in deformed potentials is given in . The projection $`\stackrel{~}{\mathrm{\Lambda }}_3`$ of pseudo orbital momentum is zero. Thus, the pseudo spin is decoupled from the orbital motion. It only reacts to the cranking term $`\stackrel{}{\omega }\stackrel{}{S}`$, where $`\stackrel{}{S}`$ is the pseudo spin. The two parallel, nearly horizontal trajectories with the distance $`\mathrm{\Delta }e_{pso}=\omega `$ correspond to the pseudo spin being aligned or anti aligned with the rotational axis $`\stackrel{}{\omega }`$. The pseudo spin vector follows the tilt of $`\stackrel{}{\omega }`$, remaining parallel to it. Since the pseudo orbital momentum remains small, the two trajectories are almost horizontal. The signature $`\alpha `$ is gradually lost when the pseudo spin vector tilts away from the 1 - axis.
The states with the highest $`K`$ values of the $`h_{11/2}`$ and $`i_{13/2}`$ intruder orbitals obey the DAL coupling. The states with lower $`K`$ have an extended region around $`90^o`$ where the Routhians are relatively flat functions of $`\vartheta `$, what means that $`j_{}`$ is small. In this region the orbitals are rotational aligned (RAL), precessing around the rotational axis $`\stackrel{}{\omega }`$, where the precession cone follows the tilt of the axis. The signature is gradually lost when $`\stackrel{}{\omega }`$ tilts away from the 1 - axis. With decreasing $`\vartheta `$, they make a quasi crossing with other members of the same intruder orbital. These crossings mark the transition to the deformation aligned (DAL) coupling, which is shows up as the $`K\mathrm{cos}\vartheta `$ behavior.
Fig. 3 shows the quasi neutron energies, which are relevant when the pair correlations are important. What has been said about the single particle Routhians also applies to the quasi particle Routhians. As a new type, the Fermi aligned (FAL) coupling appears. It is realized by the lowest $`i_{13/2}`$ trajectory, denoted by A. The FAL coupling appears at some distance from $`90^o`$. It corresponds to a substantial component $`j_3`$ as well as to a substantial $`j_1`$. It is most favored at the minimum of $`e_A^{}(\vartheta )`$ at $`\vartheta =38^o`$, where $`j_{}=0`$, that is $`\stackrel{}{j}||\stackrel{}{\omega }`$. With $`\vartheta 0`$ the non rotating quasi particle state is approached, i. e. $`j_3K`$ and $`j_10`$, corresponding to the maximum. Overall, the lowest $`i_{13/2}`$ trajectory A is rather flat, indicating that the orientation of $`\stackrel{}{j}`$ does never too strongly deviate from $`\stackrel{}{\omega }`$. At larger $`\omega `$, where the negative and positive quasi particle states strongly interact with each other, a complex pattern of avoided crossings emerges, which we have not found a simple interpretation for.
The high frequency regime is illustrated in figs. 6 and 8. It is characterized by many avoided crossing between the orbitals. This indicates the progressive dissolution of the approximate conservation of the $`K`$ quantum number for the DAL orbitals.
### L Relation to the PAC treatment of high-$`K`$ bands
Bands with a finite value of $`K`$ have been studied by means of the standard PAC scheme using the following prescription , which may be considered as an approximation to TAC. A fixed $`K`$ value is ascribed to each band, which is the spin value at the band head. It is taken from experiment or calculated by means of the cranking model choosing $`\stackrel{}{\omega }`$ parallel to the symmetry axis. It is assumed that $`J_3=K`$, independent of $`\omega `$ and $`J_1=<J_1>`$. The configuration $`|>`$ is generated by the quasi particle Routhian (32) assuming $`\vartheta =90^o`$. Only the reduced cranking term $`\omega _1J_1`$ appears, where $`\omega _1`$ is the 1 - component of the angular velocity.
With these assumptions TAC goes over into the CSM scheme: The constraint (17) becomes
$$J=I+1/2=\sqrt{J_1^2+J_3^2}=\sqrt{<J_1>^2+K^2},$$
(80)
Solving for $`<J_1>`$, the standard cranking constraint
$$<J_1>=J_1\sqrt{(I+1/2)^2K^2}$$
(81)
to fix $`\omega _1`$ is obtained. In the CSM one uses $`\omega _1`$ as the independent variable. Experimental values of $`\omega _1`$ are derived by means of the expression
$$\omega _1=\frac{E(I)E(I2)}{J_1(I)J_1(I2)},$$
(82)
with $`J_1`$ being the rhs. of expression (81). The TAC condition $`\stackrel{}{J}||\stackrel{}{\omega }`$ implies
$$\omega _1=\frac{J_1}{J}\omega =\frac{\sqrt{(I1/2)^2K^2}}{2(I1/2)}(E(I)E(I2)),$$
(83)
where the expression (72) for the experimental frequency is used. The $`\omega _1`$ values obtained by expressions (82) and (83) almost coincide, except near the band head. As demonstrated in the model study , expression (83) reproduces the quantal results slighly better.
With the fixed-$`K`$ assumption the expressions for the electro-magnetic matrix elements given in sect. II F become the ones of the semiclassical vector model of ref. . For the magnetic moments, the vector model additionally assumes that each quasi particle has a fixed value of $`j_{3,i}`$, which is $`\omega `$ independent and given by the $`K_i`$ of the Nilsson label. The individual $`j_{1,i}`$ values of the quasi particles are either calculated or extracted from differences between the experimental curves $`J_1(\omega _1)`$ ( the quasi particle alignments ). The magnetic moments are approximated by
$$\stackrel{}{\mu }=g_K\stackrel{}{j},$$
(84)
where the gyromagnetic ratio $`g_K`$ is either calculated by means of the Nilsson model or taken from experiment. In TAC expression (45,46) the components of the magnetic moments are calculated, attenuating the free spin magnetic moment of the proton and neutron by a constant factor.
In ref. an additional term is introduced which permits the calculation of the signature dependence of the $`B(M1)`$ values for the case that the $`J_3`$ component is generated by only one quasi particle . It is not expected that this correction is also applicable if $`J_3`$ is generated by many quasi particles , whereas the vector model without the signature term also applies to this more general case.
As discussed, the standard CSM becomes a good approximation of TAC if for the active quasi particles
i) $`j_{3,i}(\omega ,\vartheta )`$ can be approximated by the constant $`K_i`$ and
ii) $`j_{1,i}(\omega ,\vartheta )`$ can be approximated by $`j_{1,i}(\omega _1,90^o)`$.
Then
$$e_i^{}(\omega ,\vartheta )e_i^{}(\omega _1,90^o)\omega _3K_i,$$
(85)
where $`\omega _1=\omega \mathrm{sin}\vartheta `$ and $`\omega _3=\omega \mathrm{cos}\vartheta `$ (cf. eqs. (35) and (37)). The assumption $`J_3K`$ is justified and the tilt angle is given by $`\mathrm{cos}\vartheta =K/J_1`$, which leads to the expressions of the vector model .
Fig. 9 shows the angular momentum components $`j_1(\omega ,\vartheta )`$ and $`j_3(\omega ,\vartheta )`$ of some representative quasi particles. They are compared with the CSM values $`j_{1,i}(\omega _1,90^o)`$ and $`K_i`$, respectively. The DAL quasi neutron E obeys i) and ii) rather well. The orbital G shows the behavior $`j_1\frac{1}{2}\mathrm{sin}\vartheta `$ and $`j_3\frac{1}{2}\mathrm{cos}\vartheta `$, which is characteristic for the pseudo spin singlet. This is at variance with i) and ii). Although the contribution to the angular momentum is small, the characteristic spacing of $`\mathrm{\Delta }e_{pso}=\omega `$ between the two pseudo spin partners is not obtained when $`\vartheta `$ is substantially below $`90^o`$. The intruder orbital A shows in the range $`40^o<\vartheta <60^o`$ the typical FAL behavior, which is fairly well reproduced by the approximations i) and ii). For larger values of $`\vartheta `$ it changes gradually into a state with a good signature. The transition is accompanied by changes of $`j_1`$ and $`j_3`$ which are at variance with i) and ii). For smaller values of $`\vartheta `$ orbital A keeps its FAL character and remains close to the approximations i) and ii). The orbital B changes dramatically when $`\vartheta `$ decreases from $`90^o`$, because there is the quasi crossing with the down sloping orbital C (cf. fig. 3). For values of $`\omega `$ smaller than displayed this crossing is rather sharp. One can follow the FAL branch of B, which is well approximated by i) and ii), below the crossing. For larger values of $`\omega `$ the two orbitals strongly mix and a new $`\vartheta `$ dependence emerges, which is shown fig. 11. Obviously such changes of the quasi particle structure cannot be described by means of the traditional CSM treatment basing on the assumptions i) and ii).
The use of $`\omega _1`$ as the rotational parameter in the CSM has the advantage that all configurations can be constructed from one and the same quasi particle diagram $`e_i(\omega _1,90^o)`$, which is a great simplification and has lead to the popularity of this approach. Another pleasant feature is that the signature splitting appears in a gradual way. The disadvantage is that it is only an approximation to the TAC mean field solution, the latter being completely self-consistent and more accurate when $`\vartheta `$ substantially deviates from $`90^o`$ . Sometimes the differences are only of quantitative nature, but there are many cases where qualitatively different results are obtained. Magnetic Rotation of weakly deformed nuclei is a conspicuous example, which will not be discussed in this paper. In the following discussion of examples we shall point out the differences between TAC the standard CSM.
High-$`K`$ bands, which are in experiment near yrast, appear in the CSM as relatively high lying configurations, embedded into the back ground of many configurations with low $`K`$. In TAC they are low lying configurations. The reason can be seen in (85). CSM uses $`e_i^{}(\omega _1,90^o)`$, which is shifted up by $`\omega _3K_i`$ with respect to $`e_i^{}(\omega ,\vartheta )`$ used in TAC. It is also noted that only the TAC mean field solutions can be improved by means of RPA corrections in a systematic way. The PAC configurations corresponding to a finite $`K`$ are instable.
It is quite common to present the experimental branching ratios $`\frac{B(M1)}{B(E2)}`$ as effective values of the ratio $`\left|\frac{g_Kg_R}{Q_o}\right|`$, which would determine the branching ratio if the strong coupling limit was valid . This popular way of representing the data has the advantage that the ratio becomes constant when approaching the strong coupling limit. One may convert the ratios of $`\frac{B(M1)}{B(E2)}`$ calculated by means of TAC into effective ratios $`\left|\frac{g_Kg_R}{Q_o}\right|`$. The pertinent relation
$$\left|\frac{g_Kg_R}{Q_o}\right|=\left(\frac{5(J^2K^2)B(M1)}{16J^2B(E2)}\right)^{\frac{1}{2}}$$
(86)
is obtained from the expressions in section II F by making the assumption of strong coupling, $`J_3=K`$, $`\mu _3=(g_Kg_R)K`$ and $`\mu _1=0`$. The square root of the branching ratio becomes the product of $`\left|\frac{g_Kg_R}{Q_o}\right|`$ and the inverse of the geometric factor on the right hand side of (86). In order to avoid any misunderstanding it is noted that (86) is just a way to present the results of the exact TAC calculations, which do not make any strong coupling approximation.
## III Multi-quasi particle configurations near <sup>174</sup>Hf
This section will explain how to construct multi quasi particle configurations in the TAC scheme and how to interpret them as rotational bands. The nuclides <sup>174,175</sup>Hf, <sup>175</sup>Ta and <sup>174</sup>Lu serve as examples. The SCTAC scheme is used for the calculations. In order to simplify the discussion, the same deformation parameters $`\epsilon =0.258`$, $`\epsilon _4=0.034`$ and $`\gamma =0^o`$ are assumed for all the nuclides considered. The set represents an average of equilibrium shapes calculated for several configurations and frequencies $`\omega `$. The actual values scatter within the interval $`0.25<\epsilon <0.27`$, but the differences in deformation do not change the discussed quantities in a substantial way. We consider both the cases of no pairing and a constant pair field. For the case of finite pairing we use the prescription of the CSM , which has turned out to give very reasonable description of multi-quasi particle bands in the traditional PAC scheme. Accordingly, $`\mathrm{\Delta }_n=0.69MeV`$ and $`\mathrm{\Delta }_p=0.75MeV`$, which is 80% of the experimental even -odd mass difference. The chemical potentials $`\lambda _n`$ and $`\lambda _p`$ are fixed to the values that give the correct particle numbers for the ground state (configuration at $`\omega =0`$). This scenario provides a good description of configurations up to two excited quasi particles and a frequency of about $`0.35MeV`$. It will be discussed in the subsections III D \- III H. For higher frequency and more excited quasi particles we take zero neutron pairing into consideration. Self-consistency for $`\mathrm{\Delta }`$ is invoked along the lines described in sect. II J in order to decide where the transition to zero pairing is located. Subsection III J discusses this regime.
### A Construction of multi quasi particle configurations
For zero pairing the configurations are generated by filling up the lowest $`Z`$ and $`N`$ single particle levels and then making particle-hole excitations.
In the case of pairing the configurations are constructed from the quasi particle Routhians. The quasi particle spectrum is symmetric with respect to $`e^{}=0`$ and the double dimensional occupation scheme, discussed for the PAC solutions in ref , is applied: If a quasi particle state $`i`$ is occupied, its conjugate partner $`i^+`$ must be free. In contrast to the PAC case, the conjugate states in general do not have opposite signature, which is only for $`\vartheta =90^o`$ a good quantum number.
Diabatic tracing turns out very practical for identifying the conjugate states. They always cross sharply because they are orthogonal. As in the PAC scheme, one must be careful in choosing the right particle number parity when quasi particle trajectories cross the zero line. The most simple way is to start at sufficiently low $`\omega `$, where there is still a gap between the positive and negative solutions. There it is clear how to excite an odd or even number of quasi particles . Keeping the occupation by diabatic tracing, the particle number parity of the configuration is conserved.
In order to efficiently label the configurations a compact notation which indicates the quasi particle composition is desirable. We follow the well-tried practice of the CSM assigning letters to the quasi particle trajectories and quoting the excited quasi particles in parenthesis. The letters A, B, C, D denote positive parity quasi neutrons , E, F, G, H, … negative parity quasi neutrons , a, b, c, d positive parity quasi protons and e, f, g, h ,… negative parity quasi protons .
The letter code becomes to some extend ambiguous when the structure of the quasi particles strongly changes with the frequency $`\omega `$ and the tilt angle $`\vartheta `$. The positive parity $`i_{13/2}`$ orbitals are most susceptible to the inertial forces. Fig. 3 shows the complex pattern of quasi particle trajectories, which strongly interact with each other and interchange their character as functions of $`\omega `$ and $`\vartheta `$. An example are the orbitals B and C in fig. 3. As discussed already in sect. II L, they quasi-cross each other near $`\vartheta =60^o`$. For $`\omega =0.2MeV`$, (not shown) the crossing is still rather sharp and it would be natural to follow each trajectory diabatically, i. e. for $`\vartheta <60^o`$ to call the upper trajectory B and the lower one C. For $`\omega =0.4MeV`$ (see fig. 11) they interchange their character very gradually. Now it is more natural to call the lower trajectory B and the higher C throughout the mixing region. Fig. 10 shows the quasi neutron trajectories as functions of $`\omega `$. For $`\vartheta =45^o`$ (lower panel) the crossings are rather sharp for most trajectories. Here it is natural to keep the labels in a diabatic way, as indicated. For $`\vartheta =90^o`$ the crossings between the $`i_{13/2}`$ trajectories are much softer and the question arises of how to label them after the first quasi crossing. The suggested labeling tries to follow the quasi particle trajectories both in the $`\omega `$ and the $`\vartheta `$ direction such that the structural change is as gradual as possible. It connects the two diagrams $`e_i^{}(\omega ,\vartheta =45^o)`$ and $`e_i^{}(\omega ,\vartheta =90^o)`$ the most natural way. via the $`\vartheta `$ degree of freedom at high frequency (cf. figs. 3 and 11). This implies that the smooth crossing between A and B<sup>+</sup> as function of $`\omega `$ at $`\vartheta =90^o`$ must be treated diabatically.
It should be pointed out that the suggested labeling is a compromise. As discussed above, for $`\omega =0.2MeV`$ the quasi neutrons B and C cross sharply as functions of $`\vartheta `$. In the adopted labeling the lower of the two levels is B and the higher C. It would be more natural to follow the structures in $`\vartheta `$ direction diabatically through the crossing. But a relabeling that accounts for this leads to problems a high $`\omega `$, where the adopted labeling is most natural. The difficulty to label the strongly interacting quasi particle trajectories in a simple way has a topological origin, which can be best understood if one follows in one of the quasi particle diagram 2, 3, 5, 7 a trajectory on the $`(\omega ,\vartheta )`$-path: $`(0,90^o)(0.3MeV,90^o)(0.3MeV,0^o)(0,0^o)`$. For weakly interacting trajectories, as most with normal parity, one returns to the same quasi particle . For the intruder trajectories, as $`i_{13/2}`$ and $`h_{11/2}`$, this is not always the case.
Trying to keep the notation as simple as possible we shall assign the low-$`\omega `$ composition to a configuration and shall not change it when a crossing is encountered. The structural change can be figured out from the quasi particle diagrams.
### B Elimination of spurious states
Each configuration with an equilibrium angle $`\vartheta _o90^o`$ is associated with a $`\mathrm{\Delta }I=1`$ rotational band (TAC solution), whereas each configuration with an equilibrium angle $`\vartheta _o=90^o`$ is associated with a $`\mathrm{\Delta }I=2`$ rotational band (PAC solution). Of course, the number of quantal states cannot abruptly double when the equilibrium angle moves away from $`90^o`$. Hence, one has to be careful in avoiding spurious states. This problem was studied in ref. for the model system of one and two quasi particles coupled to a rotor. An elimination scheme has been suggested which is based on the following principle: The number of TAC configurations must be the same as the number of PAC configurations at $`\vartheta =90^o`$, where they emanate from. This means, for each TAC minimum, which is interpreted as a $`\mathrm{\Delta }I=1`$ band composed of two $`\mathrm{\Delta }I=2`$ sequences, one has to discard one configuration. One finds this spurious state most easily by taking into account that the function $`E^{}(\vartheta )`$ is symmetric with respect to $`90^o`$. If one configuration ($`J_3K>0`$ ) has a minimum at $`\vartheta _o`$ its mirror image ($`K<0`$ ) has the minimum at $`90^o\vartheta _o`$. Tracing the function $`E^{}(\vartheta )`$ of the mirror image diabatically through $`\vartheta =90^o`$, one arrives at the spurious configuration that must be discarded.
The tilt angle $`\vartheta _o`$ increases with $`\omega `$. When it approaches $`90^o`$ one has to switch from the TAC to the PAC interpretation. This results in a discontinuity of the function $`E^{}(\omega )`$ for the unfavored (upper) signature branch: In the TAC scheme it is is degenerate with the favored branch whereas in the PAC scheme it is the discarded configuration , which is now taken into account. The angle for switching from TAC to PAC is to some extend arbitrary. We have found it reasonable to use the PAC interpretation when the equilibrium angle $`\vartheta _o>80^o`$. If several quasi particles combine into high-$`K`$ and low-$`K`$ configurations it is important to switch from TAC to PAC for both the high- and the low-$`K`$ configurations at the same $`\omega `$. As discussed in and in sect. III G for a concrete example, changing to PAC only for a part of the configurations results in highly nonorthogonal states.
As an example, let us consider the most simple case of the one quasi neutron configurations denoted by \[E\] and \[F\] in fig. 12, which represent, respectively, the two branches $`j_35/2`$ and -5/2 of the DAL orbital 5/2. For $`\omega =0.2`$ and $`0.3MeV`$, \[E\] has a minimum below $`80^0`$, which is interpreted as a $`\mathrm{\Delta }I=1`$ band. The diabatic continuation of \[F\] becomes the mirror image of \[E\] for $`\vartheta >90^o`$. Hence, \[F\] is spurious and must be discarded. The kink - like minimum of the upper branch of 5/2 must also be disregarded.
The elimination rules are somewhat differently formulated in ref. . The reader might find this complementary formulation instructive. The proposed scheme has been tested for the model system of one and two quasi particles coupled to a rotor . No spurious states have been found in the low lying spectrum after applying the elimination rules.
### C Band heads
Generally, a band is a quasi particle configuration whose angular momentum increases with the rotational frequency $`\omega `$. Its structure changes gradually with $`\omega `$, such that it remains similar for adjacent quantal states of the band. This is a natural definition which permits calculating both the start and the termination of a band. Since we restrict ourselves to well deformed nuclei, we shall discuss only the start in this paper.
Fig. 13 illustrates how the configuration \[E\] in <sup>175</sup>Hf starts. The function $`E^{}(\omega ,\vartheta )`$ has a minimum at $`\vartheta =0`$ (and 180<sup>o</sup>) for $`\omega `$ below the band head frequency of $`\omega _h=0.08MeV`$. In this range of $`\omega `$ the band has not yet started, because angular momentum does not depend on $`\omega `$, being $`J=K`$. The band actually starts at $`\omega _h`$ when the equilibrium value $`\vartheta _o`$ becomes finite, i. e. when the minimum of $`E^{}(\vartheta )`$ at $`\vartheta =0^o`$ turns over into a maximum and there appears a minimum at $`\vartheta _o>0^o`$. That is, the frequency $`\omega _h`$ where
$$\frac{^2E^{}(\omega _h,\vartheta )}{\vartheta ^2}_{\vartheta =0}=0$$
(87)
has the physical meaning of the rotational frequency of the band head. Fig. 14 shows the equilibrium angles $`\vartheta _o`$ for several one quasi neutron bands. The bands heads lie where $`\vartheta _o`$ bifurcates from the zero line.
The experimental band head frequency is the energy of the first transition $`I=K+1I=K`$. It should be compared with the frequency where TAC gives $`J(\omega )=K+1`$, which is somewhat larger than $`\omega _h`$. In some of the figures (e. g. 19) this frequency is indicated by a fat dot. In most of the figures the calculated curves start with the first grid point for which $`\vartheta _o>0`$, i. e. $`\omega _h`$ is only determined with the accuracy of $`\mathrm{\Delta }\omega =0.05MeV`$, which is the step used in the calculations.
One may distinguish between strong coupling behavior and more complex response near the band head. Strong coupling behavior corresponds to $`J_3=K=const.`$ and $`J_1=\omega _1𝒥`$, where $`𝒥`$ is the moment of inertia of the collective rotation. In this case one has
$$J=\omega 𝒥,\vartheta _o=\mathrm{arccos}(\frac{K}{\omega 𝒥})for\omega >\omega _h=\frac{K}{𝒥}.$$
(88)
Axial nuclei are close to the strong coupling limit near the band head if only DAL quasi particles are excited. The configuration \[E\] illustrated in figs. 12 \- 14 is of his type. At a first glance, one might expect that the TAC approximation, which treats the orientation angle $`\vartheta `$ in a static way, becomes a bad approximation near the band head, because the $`E^{}(\vartheta )`$ is very flat there. The model studies demonstrated that this is not the case. In fact, the wave function becomes narrow at the band head, because $`J_3`$ approaches the good quantum number $`K`$. One may interpret this as follows. The mass parameter associated with the zero point motion in $`\vartheta `$ increases faster than the curvature near the minimum, $`^2E^{}(\omega _h,\vartheta )/\vartheta ^2_{\vartheta =\vartheta _o}`$.
A more complex situation is encountered when one or more quasi particles easily align with the rotational axis. Configuration \[A\] in fig. 14 is an example. The band starts significantly earlier than expected for a strongly coupled 7/2 band with a jump of $`\vartheta _o`$. For cases like this ref. found that TAC approximates the quantal particle rotor calculation less well near the band head, but becomes again a very good approximation for higher frequency.
### D The zero quasi particle configuration
In <sup>174</sup>Hf, the lowest configuration at $`\vartheta =90^o`$ is the vacuum with all negative levels occupied, which is the ground (g-) configuration at low $`\omega `$. Around $`\omega =0.30MeV`$, the neutron system gradually changes into s-configuration, which is seen in figs. 3 and 10 (upper panel) as the quasi crossings of trajectories A with B<sup>+</sup> and B with A<sup>+</sup> (AB crossing). Near these crossings the $`\vartheta `$ dependence of the trajectories is complicated. We shall return to the interpretation of this region in sect. III H.
First, let us discuss the proton system, which does not have such a crossing in the considered frequency interval. At $`\vartheta =90^o`$, the configuration has the character of the g-band. It keeps this character when $`\vartheta `$ decreases, provided the occupation is followed diabatically, i. e. the crossings between the $`\pi =+`$ trajectories at $`\vartheta =22^o`$ and the $`\pi =`$ trajectories at $`\vartheta =8^o`$ are ignored. It becomes the ground state for $`\vartheta =0`$, because the wave function does not depend on $`\omega `$ for this orientation. The ground state is not the lowest configuration at $`\vartheta =0`$, because a number of quasi particles have crossed the zero line and crossed each other. This example demonstrates the advantage of diabatic tracing, which automatically finds when starting from either small $`\omega `$ and $`\vartheta =0^o`$ or $`\vartheta `$ close to $`90^o`$, where is the lowest configuration . The neutron system has an analogous structure for $`\omega <0.25MeV`$, where where has the character of the g-band. Fig. 3 shows the quasi neutron trajectories at $`\omega =0.3MeV`$. It is seen that the configuration , which has a mixed g- and s- character $`\vartheta =90^o`$, becomes the ground state at $`\vartheta =0^o`$, where it is no longer the lowest configuration .
Fig. 15 shows the total Routhian $`E^{}(\omega =0.2MeV,\vartheta )`$ of the combined proton and neutron configurations . Its $`\vartheta `$ dependence reflects the g-band character: The angular momentum is collective, i. e.
$$E^{}(\omega ,\vartheta )\frac{\omega _1^2}{2}𝒥=\frac{(\omega \mathrm{sin}\vartheta )^2}{2}𝒥.$$
(89)
The minimum lies a $`\vartheta =90^o`$, the signature is $`\alpha =0`$, corresponding to the even spin g -band. For $`\omega =0.3MeV`$ the level repulsion near between A and B<sup>+</sup> modifies the slightly the $`\vartheta `$ dependence of the total Routhian.
In order to make a first qualitative estimate of the tilt angle for multi quasi particle configurations one has to add this zero quasi particle Routhian to the sum of the Routhians $`e_i^{}(\vartheta )`$ of the excited quasi particles , which can be taken from the quasi particle diagrams.
For $`\omega >0.35MeV`$ the quasi particle vacuum has the character of the s-configuration at $`\vartheta =90^o`$. With decreasing $`\vartheta `$, it changes into the t-configuration which becomes the $`K^\pi =8^+`$ configuration $`[7/2^+,9/2^]`$ at $`\vartheta =0^o`$. We shall discuss these changes and their consequences in sect. III H, together with the two quasi neutron excitations of positive parity.
### E One quasi neutron configurations
They are generated by adding one quasi neutron to the configuration . Fig. 12 shows their total Routhians $`E^{}(\omega ,\vartheta )`$.
The configurations \[G\] and \[H\] have $`\vartheta _o=90^o`$ for all $`\omega `$. They are interpreted as $`\mathrm{\Delta }I=2`$ bands. They represent the two signatures $`\alpha =\pm 1/2`$ of the pseudo spin singlet 1/2. Since the pseudo spin is decoupled (cf. sect. II K), $`E^{}(\vartheta )=E_{[0]}^{}(\vartheta )+const\pm \omega /2`$ and $`\vartheta _o=90^o`$.
For $`\omega =0.2`$ and $`0.3MeV`$ the configurations \[A\] and \[E\] have minima at $`90^o>\vartheta _o>0^o`$. They are interpreted as $`\mathrm{\Delta }I=1`$ bands ($`K^\pi =7/2^+`$ and $`5/2^{}`$). The configurations \[B\] and \[F\] are the continuations of \[A\] and \[E\] reflected through $`\vartheta =90^o`$. Accordingly they are discarded as spurious states together with the kink at $`90^o`$. Then the condition is satisfied, that the number of states is the same as for the PAC interpretation at $`\vartheta =90^o`$. For $`\omega =0.4MeV`$ the minima of \[A\] and \[E\] have moved above $`80^o`$. Now we change to the PAC interpretation and refer to the calculations at $`\vartheta =90^o`$. Both \[A\] and \[B\] are interpreted as $`\mathrm{\Delta }I=2`$ bands. They form the signature pair $`(\pi ,\alpha )=(+,\pm 1/2)`$. The configurations \[E\] and \[F\] represent the two $`\mathrm{\Delta }I=2`$ bands combining to the signature pair $`(,\pm 1/2)`$.
Fig. 16 shows the calculated total Routhians. The change from the TAC interpretation to the PAC is seen as the sudden onset of the signature splitting. The width of the transition region is determined by the size of the calculation grid in $`\omega `$, which is $`0.05MeV`$ and does not bear any physical relevance. As discussed above, the TAC approach is not able to describe the smooth transition from broken to conserved signature symmetry. In order to describe the gradual onset of the signature splitting one has to go beyond the pure mean field theory .
Since the orbital E obeys the DAL coupling, the configuration \[E\] is expected to be close to the strong coupling limit. Fig. 14 shows that the $`5/2^{}`$ band starts at $`\omega _h=0.09MeV`$ near the strong coupling estimate $`K/𝒥`$. Also for higher $`\omega `$ the tilt angle $`\vartheta _o`$ remains close to the strong coupling value. The $`7/2^+`$ band starts at $`\omega _h=0.02MeV`$, below the $`5/2^{}`$ band \[E\] and much below the strong coupling estimate for the $`K=7/2`$ band. This indicates a substantial deviation from strong coupling. In fact, fig. 14 shows that $`\vartheta _o`$ jumps to a finite value at a low frequency, corresponding to the rapid transition from the DAL to FAL coupling at the band head.
Fig. 16 also shows the experimental Routhians . The relative position of the Routhians as well as their slopes (i. e. the angular momentum ) are reasonably well reproduced by the calculation. The frequency of the first transition is well described too. In particular the low value of $`\omega _h`$ for the $`7/2^+`$ band indicates that the FAL coupling is seen in the experiment.
In the TAC calculation the configuration \[C\] starts at $`\omega _h=0.19MeV`$ with $`J=6.1`$ and $`\vartheta _o=60^o`$ (cf. fig. 12). It represents the $`K=9/2^+`$ orbital. The configuration \[D\] is discarded because it is the continuation of the mirror image of \[C\]. The minimum rapidly moves towards $`90^o`$, due to the strong admixture of $`i_{13/2}`$ components with low $`K`$. After $`\omega =0.25MeV`$ one has to change to the PAC interpretation, where both \[C\] and \[D\] are interpreted as the signature pair $`(+,\pm 1/2)`$. Experimentally, only the transition $`I=9/211/2`$ is seen at $`\omega =0.153MeV`$, which is lower than the calculated value of $`\omega _h`$.
The panel $`\omega =0.2MeV`$ in fig. 12 shows that the minimum of \[C\] is very shallow. For $`\omega =0.19MeV`$ it becomes a shoulder. In contrast to the experiment, there is no $`\vartheta >0`$ solution for lower frequency. Here, a limitation of the TAC approximation is encountered. The tilt angle $`\vartheta _o`$ is found in a static way by searching for the minimum of the Routhian $`E^{}(\vartheta )`$. This is an approximation to studying the dynamics of the $`\vartheta `$ degree of freedom. The static TAC treatment is expected to give good results as long as there exists a certain convex region around $`\vartheta _o`$. Then $`\vartheta `$ will execute symmetric oscillations and averaging over them will result in values close to the ones for $`\vartheta _o`$. The model studies in ref. have demonstrated this for the lowest configurations. It is clear for a curve like \[C\] that averaging over the the collective wave function in $`\vartheta `$ needs not to give values close to the ones obtained for $`\vartheta _o`$, in particular, when the minimum has become a shoulder. Then the dynamics of $`\vartheta `$ must be explicitly calculated. Refs. have addressed this problem in the frame work of the Generator Coordinate Method.
Situations like the discussed one become more likely if one considers excited configurations . As seen in fig. 12, the flat behavior of \[C\] may be thought as the consequence of interaction (repulsion) with the configurations below and above. This problem is not special to the orientation degree of freedom. Analog restrictions of the static HFB approximation are encountered when it is used to calculate the shape of excited configurations.
Since the $`J_3<4.5`$ for all the one quasi neutron configurations the tilt angle $`\vartheta _o`$ rapidly increases with the frequency. As seen in fig. 14, $`\vartheta _o>60^o`$ for $`\omega >0.2MeV`$. Accordingly, the quasi particle trajectories become similar to the ones at $`\vartheta =90^o`$. One recognizes the familiar CSM pattern of band crossings. The $`\pi =`$ bands show the AB crossing and the $`\pi =+`$ bands the delayed BC crossing, because AB is blocked (cf. e. g. ).
### F One quasi proton configurations
The lowest proton configurations are generated by occupying the orbitals e, a and c in fig. 2. They are all of DAL type and $`\vartheta _o<80^o`$. Accordingly, the configurations \[e\], \[a\] and \[c\] are interpreted as the $`\mathrm{\Delta }I=1`$ bands with $`K^\pi =7/2^{},7/2^+`$ and $`5/2^+`$, respectively. The configurations \[f\], \[b\] and \[d\] are discarded. The configuration \[g\] has always $`\vartheta _o=90^o`$, as can be expected from fig. 2. It is interpreted as the $`\mathrm{\Delta }I=2`$ band $`(,1/2)`$, i. e. the favored signature sequence of the $`h_{9/2}`$ orbital. Fig. 17 shows the calculated and the experimental Routhians in $`{}_{73}{}^{}{}_{}{}^{175}`$Ta<sub>102</sub>. All bands show the neutron AB crossing at $`\omega =0.3MeV`$. The TAC calculation for the configuration \[g\] gives too high energy and shows too early the neutron AB crossing. This is a well known problem of the $`h_{9/2}`$ band which has been discussed in the literature. The discrepancies can partially be attributed to a larger deformation. Since these questions have been addressed before and are not at the focus of this paper we have not tried to improve the agreement by optimizing the deformation.
### G One quasi proton one quasi neutron configurations
The Routhians $`E^{}(\vartheta `$) for the four combinations of the quasi protons a and b emanating from Nilsson states 7/2 with the quasi neutrons E and F emanating from 5/2 are shown in fig. 18. They are nearly degenerate at $`\vartheta =90^o`$. The configuration \[aE\] has its minimum at $`\vartheta _o=35^o`$ and represents the $`\mathrm{\Delta }I=1`$ band $`K^\pi =6^{}`$. The configuration \[aF\] has its minimum at $`\vartheta _o=78^o`$ and represents the $`\mathrm{\Delta }I=1`$ band $`K^\pi =1^{}`$. The other two configurations \[bF\] and \[bE\] continue the mirror images of \[aE\] and \[aF\]. They are discarded as spurious states. Both $`\mathrm{\Delta }I=1`$ bands are seen in <sup>174</sup>Lu. Fig. 19 shows that separation and the slope of the bands is well reproduced by the TAC calculation. Another bundle of Routhians are four combinations of a and b with A and B, emanating from the neutron states 7/2. The configuration \[aA\] represents the $`\mathrm{\Delta }I=1`$ band $`K^\pi =7^+`$. The configuration \[bB\] has no minimum, only the kink at $`\vartheta =90^o`$. It continues the mirror image of \[aA\] and is discarded. The configurations \[aB\] and \[bA\] have both their minimum at $`\vartheta _o=90^o`$. One is interpreted as a the $`\mathrm{\Delta }I=1`$ band $`K^\pi =0^+`$ and the other one is discarded. It does not matter which is chosen, because the two configurations differ only by their orientation ($`|aB>=e^{i\pi \widehat{J}_2}|bA>`$). To be definite we choose \[aB\].
When two quasi particles of the DAL type are combined into a low-$`K`$ and a high-$`K`$ configuration , one must switch from the TAC to the PAC interpretation for both configurations simultaneously. As discussed in detail in ref. , interpreting one configuration as PAC and the other one as TAC makes them highly nonorthogonal: The TAC solution is of the type $`\psi _{K_1}\psi _{K_2}`$ whereas the PAC solution is of the type $`(\psi _{K_1}\pm \psi _{K_1})(\psi _{K_2}\pm \psi _{K_2})/2`$. This has the following consequences for our example:
i) One cannot take \[aB\] at its minimum at $`\vartheta =90^o`$, because there it is of PAC type with good signature and thus nonorthogonal to the $`K=7`$ configuration \[aA\]. It suffices to take the configuration at a somewhat smaller tilt angle, for example at $`\vartheta =85^o`$, where \[aB\] has changed to the $`K=0`$ configuration . Energywise this makes barely a difference. However it is important for the calculation of the $`B(M1)`$ values, which according to (45, 46) are zero for a PAC configuration .
ii) Since the $`K^\pi =1^{}`$ configuration \[aF\] has a larger $`\vartheta _o`$ than the $`6^{}`$ combination \[aE\], its minimum approaches $`90^o`$ at a lower value of $`\omega `$. However, one must keep the TAC interpretation as long as \[aE\] is tilted. As discussed for \[aB\], one may take $`\vartheta =85^o`$. In fact, substantial $`M1`$ transitions are seen in this band, which are illustrated by fig. 20.
Fig. 19 compares the experimental Routhians in $`{}_{71}{}^{}{}_{}{}^{174}`$Lu<sub>103</sub> with the TAC calculations. The relative positions and slopes are well reproduced. Within the frequency range, all configurations are of TAC type. For the experimental $`K^\pi =0^+`$ band \[aB\] the two signatures are separated (signature splitting). This is at variance with the calculations, which assign a $`\mathrm{\Delta }I=1`$ band (cf. preceding paragraph). The discrepancy is attributed to the residual interaction, which may lead to signature dependent correlations in the low-$`K`$ bands, as well as to the zero point motion in $`\vartheta `$.
It is noted that in $`{}_{}{}^{180}{}_{73}{}^{}`$Ta<sub>107</sub> the $`K^\pi =9^{}`$ and $`0^{}`$ bands composed of the quasi neutron 9/2 and the quasi proton 9/2 both do not show signature splitting and both have a substantial M1 transitions. This is consistent with the predictions by TAC.
The orbitals H and G emanate from the pseudo spin singlet 1/2 (cf. subsection II K). The $`K^\pi =4^{}`$ configuration \[aG\] and the $`K^\pi =3^{}`$ configuration \[aH\] correspond to the parallel and anti parallel orientation of the pseudo spin with respect to the rotational axis. Accordingly, the distance between the $`4^{}`$ and $`3^{}`$ bands is equal to $`\omega `$ in the TAC calculation. The experimental distance somewhat deviates from this value, what can be seen as evidence for a pseudo spin dependence of the proton neutron interaction.
### H Two quasi neutron excitations
Fig. 15 shows the Routhians $`E^{}(\omega ,\vartheta )`$ of the lowest positive parity configurations. Let us start with $`\omega =0.2MeV`$. For $`\vartheta >60^o`$, the first three configurations are , \[AB\] and \[AC\], which have the character of g-, s- and t- configurations, respectively. For both the s- and the t- configurations the $`j_1`$ components of the two quasi neutrons add up to a large value of $`J_1`$. In the case of the s-configuration the two $`j_3`$ components are opposite in sign, resulting in $`J_30`$, whereas for the t-configuration the two $`j_3`$ components add up to a large $`J_3`$. The structure of the t- and s-bands was first discussed in , where illustrations can be found.
For $`\omega =0.2MeV`$, the configurations \[AB\] and \[AC\] change order at $`\vartheta =60^o`$. For $`\omega =0.3MeV`$, they mix strongly around $`\vartheta =65^o`$, interchanging their character. This reflects the quasi crossing between the orbital B and C, which can be seen in figs. 3 and 9. For $`\omega =0.4MeV`$, the crossing feature has disappeared. The orbitals B and C are now substantially different from what they were at low frequency. The labeling of the quasi neutron trajectories suggested in figs. 10, always assigns \[AB\] to the lower and \[AC\] to higher configuration . For $`\omega =0.2MeV`$ this results in an abrupt exchange of the structure when the configurations cross at $`\vartheta =60^o`$. \[AB\] takes the character of a t-band for $`\vartheta <60^o`$ (which changes into the $`K^\pi =8+`$ state at $`\vartheta =0`$). Their sudden exchange becomes a smooth transition at high frequency.
The change of configuration \[AB\] from the s- to the t- structure reflects the strong response of the $`i_{13/2}`$ quasi neutrons to the inertial forces, which depend on the orientation of the rotational axis. In fig. 9 it is seen for the quasi neutron B as the rapid change of $`j_3`$ from negative to positive values near $`60^o`$. This is an example of the complex rotational response which cannot be guessed within the traditional PAC scheme.
As seen in fig. 15, \[AB\] has its minimum at $`\vartheta =90^o`$. It has the signature $`\alpha =0`$ and is interpreted as the even spin s - band. It becomes yrast after the AB crossing at $`\omega =0.3MeV`$. The configuration \[AC\] has also its minimum at $`\vartheta =90^o`$. Since its signature is $`\alpha =1`$ it represents an odd spin band. The t-character of \[AB\] is only explored when additional quasi particles of DAL type change $`\vartheta `$ to smaller values. This will be discussed in subsection III I. For neutron numbers $`N106`$ the t-configuration is more favored, becoming a stable minimum. The pure two quasi neutron t-band is seen for example in $`{}_{74}{}^{}{}_{}{}^{180}`$W<sub>106</sub> and $`{}_{76}{}^{}{}_{}{}^{182}`$Os<sub>106</sub> .
Now we consider the combinations of the $`i_{13/2}`$ orbitals A, B, C, D with E and F, emanating from the Nilsson state 5/2. Fig. 21 shows the Routhians. In the lower bundle the quasi neutrons E and F are combined with A and B, emanating from 7/2. The configuration \[AE\] is the $`K^\pi =6^{}`$ band and \[AF\] the $`1^{}`$ band, both being $`\mathrm{\Delta }I=1`$ sequences. The configurations \[BE\] and \[BF\] are discarded. At $`\omega =0.4MeV`$ the minimum of \[AE\] has moved to $`80^o`$ and we switch to the PAC interpretation. Now, \[AE\] and \[BF\] have $`(\pi ,\alpha )=(,1)`$, i. e. they represent two odd spin bands, and \[AF\] and \[BE\] have $`(\pi ,\alpha )=(,0)`$, i.e. they represent two even spin bands. Fig. 22 compares the calculation with the experimental bands in $`{}_{72}{}^{}{}_{}{}^{174}`$Hf<sub>102</sub>. Whereas the $`6^{}`$ band is seen as a $`\mathrm{\Delta }I=1`$ sequence, as expected, the experimental $`1^{}`$ band shows a substantial signature splitting. The discrepancy is attributed to the residual interaction. For the low-$`K`$ negative parity configurations the octupole correlations are important. Usually they are stronger for the $`(\pi ,\alpha )=(,1)`$ bands than for the $`(,0)`$ bands . This can explain why the experimental $`(,1)`$ sequence has a low energy relative to the experimental $`(,0)`$ but also relative to the TAC calculation.
The interpretation of the bundle formed combining the quasi neutrons E and F with C and D is less straightforward. The configuration \[DE\] has a TAC minimum for $`\omega <0.25MeV`$. The component $`J_32`$, i. e. it represents the $`2^{}`$ band (cf. fig 22). Some small signature splitting is seen in experiment.
For $`\omega 0.3MeV`$ we find $`\vartheta 80^o`$ for \[DE\] and \[CE\]. Hence, the PAC interpretation is applied to the whole bundle. Accordingly, the TAC configuration \[DE\] splits into the PAC configurations \[DE\] and \[DF\], which have $`(\pi ,\alpha )=(,0)`$ and $`(,1)`$. The TAC configuration \[CE\] splits into the PAC configurations \[CE\] and \[CF\], which have $`(,1)`$ and $`(,0)`$, giving rise to two $`\mathrm{\Delta }I=2`$ sequences with odd and even spin, respectively. As seen in fig. 22, there is a $`(,0)`$ band observed, which can be interpreted as \[CF\]. The odd spin band \[CE\] should be nearby. It is not given in the experimental level scheme, but ref. reports two unplaced $`\mathrm{\Delta }I=2`$ sequences.
The interpretation of the low frequency part of \[CE\] and \[CF\] has the same problems as discussed for the configuration \[C\] in sect. III E. As seen in fig. 21 for $`\omega =0.2MeV`$, \[CE\] is rather flat. It has a minimum at $`\vartheta =0^o`$, which is due to the $`K=9/2`$ component in C. In addition it has a shoulder around $`\vartheta =60^o`$ due the $`K=5/2`$ component, which is hardly visible. However the condition of uniform rotation (35) has solutions $`\vartheta =58^o`$ and $`66^o`$ for $`\omega =0.2MeV`$ and $`0.25MeV`$, respectively. These points are included in fig. 22. The experimental $`(,0)`$ band, which we assign to \[CF\], is seen to substantial lower frequency than $`\omega =0.2`$, below which no TAC solution is found. This is another example for the limitations of the TAC approach, which treats the orientation angle $`\vartheta `$ as a static variable. The static approach is expected to work best if the function $`E(\vartheta )`$ is relatively symmetric around the minimum. Then, the fluctuations of $`\vartheta `$ are also expected to be symmetric and the contribution linear in $`\vartheta \vartheta _o`$ will average out.
The combinations of the $`i_{13/2}`$ quasi neutrons with the 7/2 orbitals will not be discussed, because the results of the calculations are similar to the ones obtained for the $`5/2^{}`$ orbitals and there is no experimental information about them.
The combinations of the quasi neutron E with G and H, which emanate from the pseudo spin singlet 1/2, are shown in fig. 23. The situation is analogous to $`{}_{71}{}^{}{}_{}{}^{174}`$Lu<sub>103</sub>, where the same orbitals combine with the quasi protons a and b. The decoupled pseudo spin just adds or subtracts one half unit of angular momentum to the total angular momentum . This means,
$`E_{[EG]}^{}=E_{[E]}^{}\omega /2+const,`$ (90)
$`E_{[EH]}^{}=E_{[E]}^{}+\omega /2+const.`$ (91)
The Routhians of \[AG\] and \[AH\] are related to the Routhian of \[A\] in a similar way (cf. figs. 16 and 23). The configurations \[AG\] and \[EG\] are observed as the $`K^\pi =4^{}`$ and $`3^+`$ bands. It would be interesting to observe the configurations \[AH\]$`3^{}`$ and \[EH\]$`2^+`$ with the pseudo spin being anti parallely oriented with respect to the rotational axis. Then one could investigate to what extend the the residual interaction depends on the pseudo spin orientation.
### I Quasi neutron configurations at a large tilt
The strong coupling estimate for the tilt angle $`\vartheta =\mathrm{arccos}(K/I)`$ shows that the number of steps in $`I`$ needed to change $`\vartheta `$ from $`0^o`$( band head) to $`90^o`$ increases with $`K`$. Hence, the high-$`K`$ bands offer the possibility to explore the quasi particle spectrum at an angle substantially different from 0<sup>o</sup> and 90<sup>o</sup>. The bands with low or moderate $`K`$, discussed in the preceding sections, only permit a sketchy view, because the rotational axis rapidly reorients from the 3- to the 1- axis. For <sup>174,175</sup>Hf there is the family of bands built on the $`K^\pi =8^{}`$ proton configuration \[ae\]. They permit us studying the quasi neutron configurations at a large tilt. These high-$`K`$ bands will be discussed in terms of the quasi neutron spectrum for $`\vartheta =45^o`$. Although the tilt angle varies along the bands, it stays below $`60^o`$ within the considered frequency range. Like the quasi particle diagrams for $`\vartheta =90^o`$ (upper panel of fig. 10) give a first orientation of the structure of the low-$`K`$ bands, the diagram for $`\vartheta =45^o`$ (lower panel of fig. 10) permits qualitative interpretation of the quasi neutron configurations in the \[ae\] family.
Comparing the two panels of fig. 10 one notices a substantial reordering of the trajectories. As a function of $`\omega `$, they cross much more sharply for $`\vartheta =45^o`$ than for $`90^o`$. This reflects the tilt of the rotational axis towards the symmetry axis, where all levels cross sharply. The sharp crossings have the consequence that the zero quasi neutron configuration \[ae\]$`8^{}`$ remains unperturbed up to the relatively high frequency of $`\omega 0.4MeV`$, where the quasi neutron A<sup>+</sup> interacts with the quasi neutron C’. This is confirmed by the full TAC calculation and the experiment shown in fig. 24. The trajectory A<sup>+</sup> crosses with a number of quasi neutron trajectories before it interacts with C’. This means that there are several two quasi neutron configurations below \[ae\]$`8^{}`$. The lowest of these, \[aeAE\], is seen in <sup>174</sup>Hf as the $`14^+`$ band. The lowest configuration of the opposite parity, \[aeAB\]$`16^{}`$, has not yet been observed. It is the quasi neutron t - configuration, which is shifted to low energy by the two DAL quasi protons a and e. Fig. 24 shows that with increasing $`\omega `$ it interchanges its character with \[aeAC\]. This is the manifestation of the quasi crossing between B and C seen in fig. 3 at $`\vartheta =65^o`$. The contrast to the quasi neutron spectrum at $`\vartheta 90^o`$ is noted. There, the configuration \[AB\] interchanges gradually character with the vacuum . For substantially smaller $`\vartheta `$, \[AB\] has a t - structure which is very different from ($`J_38`$ and 0, respectively) and, as a consequence, couples only very weakly to .
The $`18^+`$ band \[aeAEGI\] is predicted too high relative to \[aeAE\]$`14^+`$. As discussed in section III J, the discrepancy is probably due a substantial reduction of the pair field by blocking the two quasi neutrons G and I.
The quasi crossing of A<sup>+</sup> with C’, seen $`\omega 0.4MeV`$ in the lower panel of fig. 10, causes the down bend of the zero quasi neutron configuration \[ae\]$`8^{}`$ seen at $`\omega _c=0.38MeV`$ in fig. 24. It is the AB crossing at the tilt angle of $`\vartheta _o=60^o`$. It is delayed as compared to the yrast band where it appears at $`\omega _c=0.30MeV`$. Only part of the delay can be explained by geometry within the standard PAC scheme: The projection of the angular frequency on the 1 - axis is $`\omega _{c1}=\omega _c\mathrm{sin}\vartheta =\omega _c\mathrm{sin}60^o=0.33MeV`$, which is larger than $`\omega _1=0.30MeV`$, where the AB crossing appears in PAC.
Let us now discuss the \[ae\] family in the $`N=103`$ system by means of the $`\vartheta =45^o`$ quasi neutron diagram (lower panel of 10). The one quasi neutron configuration \[aeA\]$`23/2^{}`$ is lowest. It is observed at about the right energy in <sup>175</sup>Hf as the $`23/2^{}`$ band. Since E lies above A the configuration \[aeE\] is expected at higher energy. This confirmed by the full TAC calculation shown in fig. 25. In the $`\vartheta =45^o`$ diagram 10, the two quasi neutron excitations EB and EI have negative energy above $`\omega =0.4MeV`$. Fig. 25 demonstrates that after minimizing with respect to $`\vartheta `$ the configurations \[aeABE\]$`39/2^+`$ and \[aeAEI\]$`35/2^{}`$ lie below \[aeA\]$`23/2^{}`$. This order of the bands is seen in <sup>175</sup>Hf.
Hence, the experimental high $`K`$ spectra clearly reflect the modification of the quasi neutron spectrum with decreasing $`\vartheta `$.
### J Unpaired neutron configurations
In this subsection we discuss, how to interpret the high-$`K`$ bands of the \[ae\] family in terms of configurations of the unpaired neutron system. Fig. 26 shows the single neutron Routhians at $`\vartheta =45^o`$. The labels of the states are chosen such that they are as close as possible to the quasi neutron neutron labels in fig. 10. We consider $`N=102`$. In order to keep the notation similar to the case of finite pairing, let us denote by the yrast configuration for $`\omega <0.1MeV`$ ( below the crossing between the levels A and E ). At $`\vartheta =90^o`$ it has the character of the s - configuration and at $`\vartheta =0^o`$ it is the ground state. The particle hole excitations are constructed relative to this configuration \[ae0\]$`8^{}`$. Above the AE - crossing, the yrast configuration is \[aeA<sup>-1</sup>E\]$`14^+`$. The next higher configuration is \[aeA<sup>-1</sup>G<sup>-1</sup>EI\]$`18^+`$. They form the two lowest bands of the \[ae\] family. Both are shown in fig. 27 as results of the full TAC calculation. The relative position and slopes compare well with the experiment shown in the upper panel of fig. 24. Close by there is \[aeA<sup>-1</sup>G<sup>-1</sup>BE\]$`19^{}`$, which has not yet been observed. Above $`\omega =0.35MeV`$ one expects a number of configurations with higher $`K`$, generated by lifting a neutron from the levels $`9/2^{}`$ and C to I and B. One example is \[ae A<sup>-1</sup>C<sup>-1</sup>EI\]$`20^+`$.
The unpaired configurations \[aeA<sup>-1</sup>E\]$`14^+`$ and \[aeA<sup>-1</sup>G<sup>-1</sup>EI\]$`18^+`$ are also shown in fig. 24 for a comparison with the calculation at finite neutron pairing. There they are labeled in terms of the quasi particle notation as \[aeAE\]$`14^+`$ and \[aeAEGI\]$`18^+`$, respectively. The unpaired configuration \[aeA<sup>-1</sup>E\]$`18^+`$ lies significantly below the paired one. In this configuration four quasi neutrons are blocked and the description as unpaired neutron state is better (within the HFB scheme to which we restrict here). For \[aeAE\]$`14^+`$, the paired calculation is favored at low and the unpaired at high $`\omega `$. The zero pairing calculation compares better with the experiment. The zero pairing Routhians of both configurations lie too low as compared to the yrast band . This may be a consequence of the dynamical pair correlations which are not taken into account.
Fig. 28 shows the branching ratios for two bands \[aeAE\]$`14^+`$ and \[aeAEGI\]$`18^+`$. In case of the $`K^\pi =14^+`$ band, both the paired and the unpaired calculation give similar results. In case of the $`18^+`$ band, the unpaired calculation shows a similar increase at low frequency as the experiment, however it underestimates the experimental ratio. We have to underline here that the microscopic calculation of the magnetic transition probabilities by means of (45, 46) is expected to be less accurate than the popular strong coupling estimates (if applicable), which are based on quasi particle g - factors that are adjusted to the experiment.
The unpaired configuration \[ae0\]$`8^{}`$ lies slightly below the paired configuration \[ae\]$`8^{}`$ after the latter has bend down. It has the character of the neutron s - configuration . That is the pairfield becomes small when the pair of $`i_{13/2}`$ quasi neutrons is broken in the AB crossing.
The unpaired configuration \[aeA<sup>-1</sup>B\]$`8^{}`$ lies above the two paired bands \[aeAC\] and \[aeAB\], which interchange character. Thus the interpretation in terms of a quasi particle structure is favored. It will be interesting to see if the experimental bands shows the paired or unpaired pattern, which are markedly different.
As seen in fig. 26, the lowest configuration in the $`N=103`$ system is \[ae A<sup>-1</sup>EI\]$`35/2^{}`$, which is shown in fig. 25 as \[aeAEI\]$`35/2^{},\mathrm{\Delta }=0`$. It lies below the paired configuration \[aeAEI\]. This is expected because three quasi neutrons are blocked, destroying the static pair field. The relatively high band head frequency in experiment is consistent with the zero pairing calculation. In the calculation with finite pairing the band starts at a much lower frequency. The configuration \[ae A<sup>-1</sup>EF\]$`23/2^{}`$ is shown in fig. 25 as \[aeA\]$`23/2^{},\mathrm{\Delta }=0`$. It lies above the paired band \[aeA\]$`23/2^{}`$ for most of the frequency range. Finite pairing is favored for the one quasi neutron band. The zero pairing solution wins only at high frequency where a band crossing (down bend) is encountered. The band \[aeABI\]$`39/2^+`$ corresponds in the unpaired scheme to \[ae A<sup>-1</sup>BI\]$`39/2^+`$. This configuration , shown as \[aeABI\]$`39/2^+,\mathrm{\Delta }=0`$ in fig. 25, lies above the finite pairing calculation. This somewhat surprising result (static pairing for the three quasi neutron configuration ) is understood from fig. 26. In contrast to \[ae A<sup>-1</sup>EI\]$`35/2^{}`$, where the EF is blocked for pair scattering, for \[ae A<sup>-1</sup>BI\]$`39/2^+`$ the pair BB’ blocked. Since BB’ is much further away from the Fermi - surface than EF, blocking is much less effective.
For higher frequency it becomes favorable to generate angular momentum by exciting quasi protons . As seen in fig. 2, the next pair at large tilt is \[cg\]. In fact, the configuration \[aecgAEGI\]21<sup>-</sup> appears at the yrast line. In the TAC calculation the paired four quasi proton configuration \[aecg\] has a lower energy than the unpaired. It is not unexpected that the TAC calculation gives a too high energy, because already the one quasi particle configuration \[g\] is predicted too high by TAC (cf. sect. III F).
## IV Rules for TAC
Let us summarize the experience gained in applying the TAC approach to the analysis of rotational spectra in axial well deformed nuclei. It is assumed that for $`\vartheta =0`$ the rotational axis coincides with the symmetry axis. We shall use the form of rules.
1. In order to construct the configurations use the particle - hole scheme in the case of zero pairing. For finite pairing use the quasi particle occupation scheme. The quasi particle levels appear in pairs of conjugate levels, one of which is occupied and the other is empty. Only for $`\vartheta =90^o`$ the conjugate levels have opposite signature.
2. In order to check if a configuration corresponds to even or odd particle number trace it diabatically back to low $`\omega `$, where gap between the quasi particle levels of positive and negative energy exists.
3. In searching the equilibrium orientation $`\vartheta _o`$ try to stay within structurally the same configuration. Use diabatic tracing.
4. If there are avoided crossings between the levels, diabatic tracing may end up in an unwanted configuration. Usually this shows up as an irregularity of the calculated quantities as functions of $`\omega `$. In such cases one has to resort to the quasi particle diagrams and manually reassign the desired configuration .
5. Reduce the chance of unwanted configuration changes by choosing the start angle $`\vartheta _s`$ of the diabatic tracing close a to the expected equilibrium angles $`\vartheta _o`$. Draw the quasi particle diagram $`e_i^{}(\omega ,\vartheta _s)`$. Never start diabatic tracing at $`\vartheta _s=90^o`$, use $`85^o`$.
6. The crossings between quasi particle configurations represent crossings between real bands, but the description of the mixing region itself is incorrect.
7. Rotational bands correspond to a function $`J(\omega )`$ which increases with $`\omega `$. So long as $`\vartheta _o=0`$ the band has not started yet. The band head lies at the frequency where a minimum at $`\vartheta _00`$ appears.
8. Solutions with $`\vartheta _o<80^o`$ are of the TAC type. They represent $`\mathrm{\Delta }I=1`$ bands. The two signature partners $`(\pi ,\alpha )`$ and $`(\pi ,\alpha +1)`$ are degenerate.
9. Solutions with $`\vartheta _o>80^o`$ are of the PAC type. They represent $`\mathrm{\Delta }I=2`$ bands of given $`(\pi ,\alpha )`$, i. e. $`I=\alpha mod2`$. The signature $`\alpha `$ is given by its value at $`\vartheta =90^o`$.
10. When the tilt angle $`\vartheta _o(\omega )`$ becomes larger than $`80^o`$ the change from the TAC to the PAC interpretation results in unphysical jumps of the energy distance between signature partners and of other quantities. The experimental quantities show a gradual transition between the two cases, which cannot be calculated by TAC.
11. Since the number of $`\mathrm{\Delta }I=2`$ bands of given $`(\pi ,\alpha )`$ must be the same in the PAC and TAC interpretation, one half of the TAC configurations is spurious.
12. Each configuration in the region $`\vartheta <90^o`$ has its mirror image $`E^{}(\pi \vartheta )=E^{}(\vartheta )`$ in the region $`\vartheta >90^o`$. The diabatic continuation of the mirror image of an adopted configuration into the region $`\vartheta <90^o`$ is spurious and must be discarded.
13. The spurious configurations have minima or kinks only at $`\vartheta =90^o`$. If such a configuration has another minimum for $`\vartheta <80^o`$ this must be considered as physical. Then configuration has changed its character with $`\vartheta `$, being no longer spurious.
14. For a strong tilt ($`0^o\vartheta _o90^o`$), the spurious configuration are usually high in energy and do not interfere.
15. In the case of multi quasi particle configurations there are bundles of configurations emanating from $`\vartheta =90^0`$, each of which has its own $`\vartheta _o`$. Only when lowest $`\vartheta _o`$ has reached $`80^o`$ one must change from the TAC to the PAC interpretation simultaneously for the whole bundle.
16. The intra band matrix elements of the electromagnetic operators can be calculated. PAC solutions provide their signature dependence. TAC solutions give only the average over both signatures.
The rules can be applied to triaxial nuclei with few obvious modifications. Rule 5 must be complemented by: Do not use $`\vartheta _s=0^o`$, start at $`5^o`$. Only the first, general statement of rule 7 remains valid. The second which assumes that $`\vartheta =0`$ is a symmetry axis does not apply. The low frequency behavior of triaxial nuclei has not yet been studied, except the investigation of a model case in . PAC solutions are possible for all three principal axes and TAC solutions are possible in all three planes spanned by the principal axes. Accordingly one must extend the search. The simplest way is letting $`\gamma `$ vary from 120<sup>o</sup> to -60<sup>o</sup> (cf. tab. 1). If the rotational axis does not lie in one of the three planes spanned by the principal axes the rotation becomes chiral. The consequences of this 3D geometry are discussed in and .
## V Conclusions
The semiclassic concept of uniform rotation about an axis that is tilted with respect to the principal axes of the deformed density distribution leads to a mean field theory which describes energies and intra band electro-magnetic matrix elements of $`\mathrm{\Delta }I=1`$ bands in a quantitative way. The orientation of the rotational axis turns out to be as good a collective degree of freedom as the familiar shape degrees of freedom are. The tendency of high spin particles to align with the rotational axis, which in general does not have the direction of one of the principal axes of the deformed mean field, is a concept that permits to explain many features of high-$`K`$ bands from a new perspective.
The tilted solutions do not have the familiar $`C_2`$ symmetry, which appears when rotational axis coincides with one of the principal axes. The lower symmetry results in the loss of the signature quantum number, which manifests itself by the appearance of one $`\mathrm{\Delta }I=1`$ band instead two separate $`\mathrm{\Delta }I=2`$ sequences of opposite signature. The transverse magnetic dipole moment, which determines the rate of magnetic dipole transitions, plays the role of an order parameter. For tilted solutions it has a finite expectation value, which may become quite large, because it is the sum of contributions of several quasi particles . For rotation about a principal axis the expectation value of the transverse magnetic dipole moment is zero. The magnetic transition probability is given by a matrix element between two different quasi particle configurations , which is of single particle order.
Being a mean field approach, tilted axis cranking is not capable of describing the transition from a tilted to a principal axis solution in a correct way, because this involves the transition from a broken to a restored discrete symmetry. The signature dependence of the energy and other quantities appears in a sudden, unphysical way when switching from the broken symmetry to the conserved symmetry interpretation. Still, one can guess from the calculations at which rotational frequency the signature effects are expected and how strong they should be.
The breaking of the $`C_2`$ symmetry leads to the appearance of spurious states. An elimination method is suggested. After applying it the calculated sequence of the first bands above the yrast line which agrees with the observed one for the studied examples. No spurious states remain in the near yrast region.
It is the strength of the tilted axis cranking approach that treating many excited quasi particles is no more complicated than treating few. In order to demonstrate the application of the method we studied configurations with up to four excited quasi protons and four excited quasi neutrons in the nuclides with $`N=102,103`$ and $`Z=71,72,73`$. The calculated energies and branching ratios agree with the experimental values within an accuracy that is typical for microscopic mean field calculations. In particular, it is found that the order and structure of the high-$`K`$ bands can be qualitatively understood in terms configurations constructed from quasi particle levels, which are calculated as functions of the rotational frequency $`\omega `$ at a fixed tilt angle of about $`45^o`$.
The regime of quenched static pairing is encountered in the multi quasi particle bands of high seniority. Since the change of the pair field is not at the focus of this paper, it was treated in a rough way by comparing the paired with unpaired solutions and choosing the one with the lower energy. This schematic approach turned out quite practicable for a first analysis of the high-$`K`$ band structure. A refined description of pairing within the tilted cranking model, which includes dynamical pair correlations, will be given elsewhere .
The study of high-$`K`$ bands at the largest frequencies attainable is a interesting problem of nuclear physics. More systematic studies than the present are expected to reveal the response of the single particle levels to the tilt of the rotational axis. A particularly interesting question is how the $`K`$ quantum number is eroded with increasing rotational frequency. Tilted axis cranking is the proper mean field theory to address this question. It is also the appropriate starting point for theories that go beyond the mean field, like RPA.
## VI Acknowledgment
I should like to thank F. Dönau and Jing-ye Zhang for carefully reading the manuscript. The work was partially carried out under the Grant DE-FG02-95ER40934.
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# Time-reparametrization-invariant dynamics of a relativistic string
## 1. Introduction
The group of diffeomorphisms of the Hamiltonian description of relativistic systems (particles, string, branes, general relativity) - contains the Abelian subgroup of the reparametrization of the coordinate time . All known descriptions of a relativistic string are based on the reduction of the extended phase space by the fixation of gauge which breaks reparametrization - invariance from very beginning. The questions arise: Can one describe the time - reparametrization - invariant dynamics of a relativistic string dynamics directly in the terms of reparametrization - invariant variables, and what is a difference of this description from the gauge-fixing method?
To answer these questions, in the present paper, we apply a method of a reparametrization - invariant Hamiltonian description developed for gravitation .
The method of a reparametrization - invariant description is based on the reduction of an action by the explicit resolving of the first class constraints. An important element of the invariant reduction is the Levi-Civita - Shanmugadhasan canonical transformation that linearizes the energy constraint as the generator of reparametrizations of the coordinate time.
The content of the paper is the following. In Section 2 we formulate the method of the invariant Hamiltonian reduction using the simplest examples of classical mechanics and relativistic particle. Section 3 is devoted to the generalized Hamiltonian approach to a relativistic string and the statement of the problem. In Section 4, local excitations are separated from the ”center of mass” coordinates of the string. In Section 5, the Levi-Civita transformations and the invariant Hamiltonian reduction are performed to resolve the global constraint and to convert the time-like variable of the global motion into the proper time. In Section 6, the classical and quantum dynamics of local excitations are described in terms of the proper time. Section 7 is devoted to the generating functional for the Green functions.
## 2. Invariant Hamiltonian Reduction
### 2.1. Mechanics
To illustrate the time-reparametrization-invariant Hamiltonian reduction and its difference from the gauge-fixing method, let us consider an extended form of a classical-mechanical system
$$W=\underset{\tau ^1}{\overset{\tau _2}{}}𝑑\tau \left(p\dot{q}\mathrm{\Pi }_0\dot{Q}_0\lambda [\mathrm{\Pi }_0+H(p,q)]\right),$$
(1)
that is invariant under reparametrizations of the coordinate evolution parameter $`\tau `$ and ”lapse” function $`\lambda `$
$$\tau \tau ^{}=\tau ^{}(\tau ),\lambda \lambda ^{}=\lambda \frac{d\tau }{d\tau ^{}}.$$
(2)
The problem of the classical description is to obtain the evolution of the physical variables of the world space $`q,Q_0`$ in terms of the geometric time $`T`$ defined as
$$dT:=\lambda d\tau ,T=\underset{0}{\overset{\tau }{}}𝑑\tau ^{}\lambda (\tau ^{}),$$
(3)
that is also invariant under reparametrizations (2).
The second problem (connected with quantization) is to present the effective action of the equivalent unconstrained theory directly in terms of $`T`$, the equations of which reproduce this evolution. The solution of the second problem will be called the invariant Hamiltonian reduction.
The resolving of the first problem for the considered system is trivial, as the equations of motion of this system
$$\dot{q}=\lambda _pH,\dot{p}=\lambda _qH,\dot{Q}_0=\lambda ,\dot{\mathrm{\Pi }}_0=0$$
(4)
in terms of the geomeric time (3)
$$\frac{dq}{dT}=_pH,\frac{dp}{dT}=_qH,\frac{dQ_0}{dT}=1,\frac{d\mathrm{\Pi }_0}{dT}=0$$
(5)
are completely equivalent to the equations of the conventional unconstrained mechanics in the reduced phase space $`(p,q)`$
$$W_{reduced}=\underset{T(\tau _1)=T_1}{\overset{T(\tau _2)=T_2}{}}𝑑T\left(p\frac{dq}{dT}H(p,q)\right).$$
(6)
The problem is how to derive this system from the extended one (1) to apply the simplest Hamiltonian quantization with a clear physical interpretation of the invariant quantities.
The solution of the problem of the invariant Hamiltonian reduction considered in the present review is the explicit resolving of three equations of the extended system (1):
i) for the variable $`\lambda `$ (treated as constraint)
$$\frac{\delta W}{\delta \lambda }=\mathrm{\Pi }_0+H(p,q)=0,$$
(7)
ii) for the momentum $`\mathrm{\Pi }_0`$ with a negative contribution to the constraint (7)
$$\frac{\delta W}{\delta \mathrm{\Pi }_0}=0\frac{dQ_0}{d\tau }=\lambda ,$$
(8)
and iii) for its conjugate variable $`Q_0`$
$$\frac{\delta W}{\delta Q_0}=\frac{d\mathrm{\Pi }_0}{d\tau }=0.$$
(9)
(We call these three equations (7) - (9) the geometric sector.)
The resolving of the constraint (7) expresses the ”ignorable” momentum $`\mathrm{\Pi }_0`$ through $`H(p,q)`$ with a positive value $`\mathrm{\Pi }_0=H(p,q)>0`$. The second equation (8) identifies the dynamic evolution parameter $`Q_0`$ with the proper time (3) $`Q_0=T`$. It is not the gauge but the invariant solution of the equation of motion (8). The third equation (9) is the conservation law.
As a result of the invariant Hamiltonian reduction (i.e., a result of the substitution of $`\mathrm{\Pi }_0=H`$ and $`Q_0=T`$ into the initial action (1) ) this action is reduced to the one of the conventional mechanics (6) in terms of the proper time $`T`$ where the role of the nonzero Hamiltonian of evolution in the proper time $`T`$ is played by the constraint-shell value of the ”ignorable” momentum $`\mathrm{\Pi }_0=H(p,q)`$. In other words, this constraint-shell action $`W(\text{constraint})=W^M`$ determines the nonzero Hamiltonian $`H(p,q)`$ in the proper time $`T`$, instead of the zero generalized Hamiltonian in the coordinate time $`\tau `$ in (1) $`\lambda (\mathrm{\Pi }_0+H)`$.
Thus, the equivalent unconstrained system was constructed without any additional constraint of the type:
$$\lambda =1,\tau =T$$
(10)
which confuse quantities of the measurable sector with noninvariant ones. This confusion is contradictable. The ”gauge-fixing” identification of the coordinate evolution parameter $`\tau `$ and the geometric time $`dT=\lambda d\tau `$ in the form of the gauges (10) contradicts to the difference of their Hamiltonians $`\lambda (\mathrm{\Pi }_0+H)H(p,q)`$.
The second difference of the ”gauge-fixing” from the invariant Hamiltonian reduction is more essential, namely, the formulation of the theory in terms of the invariant geometric time (3) is achieved by the explicit resolving of the constraint (7) and equation of motion (8), as a result of which ”ignorable” variables $`\mathrm{\Pi }_0,Q_0`$ are excluded from the phase space.
### 2.2. Special Relativity
Let us apply the invariant Hamiltonian reduction to relativistic particle.
To answer the question: Why is the reparametrization-invariant reduction needed?, let us consider relativistic mechanics in the Hamiltonian form
$$W[P,X|N|\tau _1,\tau _2]=\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau [P_\mu \dot{X}^\mu \frac{N}{2m}(P_\mu ^2+m^2)],$$
(11)
which is classically equivalent to the conventional square root form
$$W[X|\tau _1,\tau _2]=m\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \sqrt{\dot{X}^\mu \dot{X}_\mu }$$
(12)
Both these action is invariant with respect to reparametrizations of the coordinate evolution parameter
$$\tau \tau ^{}=\tau ^{}(\tau ),N^{}d\tau ^{}=Nd\tau $$
(13)
given in the one-dimensional space with the invariant interval
$$dT:=Nd\tau ,T=\underset{0}{\overset{\tau }{}}𝑑\overline{\tau }N(\overline{\tau })$$
(14)
We called this invariant interval the geometric time whereas the dynamic variable $`X_0`$ (with a negative contribution in the constraint) we called dynamic evolution parameter.
In terms of the geometric time (14) the classical equations of the generalized Hamiltonian system (11) takes the form
$$\frac{dX_\mu }{dT}=\frac{P_\mu }{m},\frac{dP_\mu }{dT}=0,P_\mu ^2m^2=0.$$
(15)
The classical problem is to find the evolution of the world space variables with respect to the geometric time $`T`$.
The quantum problem is to obtain the equivalent unconstrained theories directly in terms of the invariant times $`X_0`$ or $`T`$ with the invariant Hamiltonians of evolution. The solution of the second problem is called the dynamic (for $`X_0`$), or geometric (for $`T`$) reparametrization-invariant Hamiltonian reductions.
The dynamic reduction of the action (11) means the substitution of the explicit resolving of the energy constraint $`(P_\mu ^2+m^2)=0`$ with respect to the momentum $`P_0`$ into this action
$$\frac{\delta W}{\delta N}=0P_0=\pm \sqrt{m^2+P_i^2}.$$
(16)
In accordance with two signs of the solution (16), after the substitution of (16) into (11), we have two branches of the dynamic unconstrained system
$$W(\text{constraint})_\pm =\underset{X_0(\tau _1)=X_0(1)}{\overset{X_0(\tau _2)=X_0(2)}{}}𝑑X_0\left[P_i\frac{dX_i}{dX_0}\sqrt{m^2+P_i^2}\right].$$
(17)
The role of the time of evolution, in this action, is played by the variable $`X_0`$ that abandons the Dirac sector of ”observables” $`P_i,X_i`$, but not the sector of ”measurable” quantities. At the same time, its conjugate momentum $`P_0`$ converts into the corresponding Hamiltonian of evolution, values of which are energies of a particle.
This invariant reduction of the action gives an ”equivalent” unconstrained system together with definition of the dynamic evolution parameter $`X_0`$ corresponding to a nonzero Hamiltonian $`P_0`$.
Thus, we need the reparametrization-invariant Hamiltonian reduction to determine the dynamic evolution parameter and its invariant Hamiltonian for a reparametrization-invariant system and to apply the symplest canonical quantization to it.
In quantum relativistic theory, we get two Schrödinger equations
$$i\frac{d}{dX_0}\mathrm{\Psi }_{(\pm )}(X|P)=\pm \sqrt{m^2+P_i^2}\mathrm{\Psi }_{(\pm )}(X|P),$$
(18)
with positive and negative values of $`P_0`$ and normalized wave functions
$$\mathrm{\Psi }_\pm (X|P)=\frac{A_P^\pm \theta (\pm P_0)}{(2\pi )^{3/2}\sqrt{2P_0}}\mathrm{exp}(iP_\mu X^\mu ),\left([A_P^{},A_P^{}^+]=\delta ^3(P_iP_i^{})\right).$$
(19)
The coefficient $`A_P^+`$, in the secondary quantization, is treated as the operator of creation of a particle with positive energy; and the coefficient $`A_P^{}`$, as the operator of annihilation of a particle also with positive energy. The physical states are formed by action of these operators on the vacuum $`<0|,|0>`$ in the form of out-state ( $`|P>=A_P^+|0>`$ ) with positive frequencies and in-state ( $`<P|=<0|A_P^{}`$ ) with negative frequencies. This treatment means that positive frequencies propagate forward ($`X_{0}^{}{}_{2}{}^{}>X_{0}^{}{}_{1}{}^{}`$); and negative frequencies, backward ($`X_{0}^{}{}_{1}{}^{}>X_{0}^{}{}_{2}{}^{}`$), so that the negative values of energy are excluded from the spectrum to provide the stability of the quantum system in QFT . For this causal convention the geometric time (14) is always positive in accordance with the equations of motion (15)
$$\left(\frac{dT}{dX_0}\right)_\pm =\pm \frac{m}{\sqrt{P_i^2+m^2}}T(X_{0}^{}{}_{2}{}^{},X_{0}^{}{}_{1}{}^{})=\pm \frac{m}{\sqrt{P_i^2+m^2}}(X_{0}^{}{}_{2}{}^{}X_{0}^{}{}_{1}{}^{})0$$
(20)
In other words, instead of changing the sign of energy, we change that of the dynamic evolution parameter, which leads to the arrow of the geometric time (20) and to the causal Green function
$$G^c(X)=G_+(X)\theta (X_0)+G_{}(X)\theta (X_0)=i\frac{d^4P}{(2\pi )^4}\mathrm{exp}(iPX)\frac{1}{P^2m^2iϵ},$$
(21)
where $`G_+(X)=G_{}(X)`$ is the ”commutative” Green function
$$G_+(X)=\frac{d^4P}{(2\pi )^3}\mathrm{exp}(iPX)\delta (P^2m^2)\theta (P_0)=$$
(22)
$$\frac{1}{2\pi }d^3Pd^3P^{}<0|\mathrm{\Psi }_{}(X|P)\mathrm{\Psi }_+(0|P^{})|0>.$$
The question appears: How to construct the path integral without gauges?
To obtain the reparametrization-invariant form of the functional integral adequate to the considered gauge-less reduction (17) and the causal Green function (21), we use the version of composition law for the commutative Green function with the integration over the whole measurable sector $`X_{1\mu }`$
$$G_+(XX_0)=G_+(XX_1)\overline{G}_+(X_1X_0)𝑑X_1,\overline{G}_+=\frac{G_+}{2\pi \delta (0)},$$
(23)
where $`\delta (0)=𝑑N`$ is the infinite volume of the group of reparametrizations of the coordinate $`\tau `$. Using the composition law $`n`$-times, we got the multiple integral
$$G_+(XX_0)=G_+(XX_1)\underset{k=1}{\overset{n}{}}\overline{G}_+(X_kX_{k+1})dX_k,(X_{n+1}=X_0).$$
(24)
The continual limit of the multiple integral with the integral representation for $`\delta `$-function
$$\delta (P^2m^2)=\frac{1}{2\pi }𝑑N\mathrm{exp}[iN(P^2m^2)]$$
can be defined as the path integral in the form of the average over the group of reparametrizations
$$G_+(X)=\underset{X(\tau _1)=0}{\overset{X(\tau _2)=X}{}}\frac{dN(\tau _2)d^4P(\tau _2)}{(2\pi )^3}\underset{\tau _1\tau <\tau _2}{}\left\{d\overline{N}(\tau )\underset{\mu }{}\left(\frac{dP_\mu (\tau )dX_\mu (\tau )}{2\pi }\right)\right\}$$
(25)
$$\mathrm{exp}(iW[P,X|N|\tau _1,\tau _2]),$$
where $`\overline{N}=N/2\pi \delta (0)`$, and $`W`$ is the initial extended action (11).
### 2.3. Geometric unconstrained system for a relativistic particle
The Hamiltonian of the unconstrained system in terms of the geometric time $`T`$ can be obtained by the canonical Levi-Civita - type transformation
$$(P_\mu ,X_\mu )(\mathrm{\Pi }_\mu ,Q_\mu )$$
(26)
to the variables ($`\mathrm{\Pi }_\mu ,Q_\mu `$) for which one of equations identifies $`Q_0`$ with the geometric time $`T`$. This transformation converts the constraint into a new momentum
$$\mathrm{\Pi }_0=\frac{1}{2m}[P_0^2P_i^2],\mathrm{\Pi }_i=P_i,Q_0=X_0\frac{m}{P_0},Q_i=X_iX_0\frac{P_i}{P_0}$$
(27)
and has the inverted form
$$P_0=\pm \sqrt{2m\mathrm{\Pi }_0+\mathrm{\Pi }_i^2},P_i=\mathrm{\Pi }_i,X_0=\pm Q_0\frac{\sqrt{2m\mathrm{\Pi }_0+\mathrm{\Pi }_i^2}}{m},X_i=Q_i+Q_0\frac{\mathrm{\Pi }_i}{m}.$$
(28)
After transformation (27) the action (11) takes the form
$$W=\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \left[\mathrm{\Pi }_\mu \dot{Q}^\mu N(\mathrm{\Pi }_0+\frac{m}{2})\frac{d}{d\tau }S^{lc}\right],S^{lc}=(Q_0\mathrm{\Pi }_0).$$
(29)
The invariant reduction is the resolving of the constraint $`\mathrm{\Pi }_0=m/2`$ which determines a new Hamiltonian of evolution with respect to the new dynamic evolution parameter $`Q_0`$, whereas the equation of motion for this momentum $`\mathrm{\Pi }_0`$ identifies the dynamic evolution parameter $`Q_0`$ with the geometric time $`T`$ ($`dQ_0=dT`$). The substitution of these solutions into the action (29) leads to the reduced action of a geometric unconstrained system
$$W(\text{constraint})=\underset{T_1}{\overset{T_2}{}}𝑑T\left(\mathrm{\Pi }_i\frac{dQ_i}{dT}\frac{m}{2}\frac{d}{dT}(S^{lc})\right)(S^{lc}=Q_0\frac{m}{2}),$$
(30)
where variables $`\mathrm{\Pi }_i,Q_i`$ are cyclic ones and have the meaning of initial conditions in the comoving frame
$$\frac{\delta W}{\delta \mathrm{\Pi }_i}=\frac{dQ_i}{d\tau }=0Q_i=Q_i^{(0)},\frac{\delta W}{\delta Q_i}=\frac{d\mathrm{\Pi }_i}{d\tau }=0\mathrm{\Pi }_i=\mathrm{\Pi }_i^{(0)}.$$
(31)
The substitution of all geometric solutions
$$Q_0=T,\mathrm{\Pi }_0=\frac{m}{2},\mathrm{\Pi }_i=\mathrm{\Pi }_i^{(0)}=P_i,Q_i=Q_i^{(0)}$$
(32)
into the inverted Levi-Civita transformation (28) leads to the conventional relativistic solution for the dynamical system
$$P_0=\pm \sqrt{m^2+P_i^2},P_i=\mathrm{\Pi }_i^{(0)},X_0(T)=T\frac{P_0}{m},X_i(T)=X_i^{(0)}+T\frac{P_i}{m}.$$
(33)
The Schrödinger equation for the wave function
$$\frac{d}{idT}\mathrm{\Psi }(T,Q_i|\mathrm{\Pi }_i)=\frac{m}{2}\mathrm{\Psi }(T,Q_i|\mathrm{\Pi }_i),$$
(34)
$$\mathrm{\Psi }(T,Q_i|\mathrm{\Pi }_i)=\mathrm{exp}(iT\frac{m}{2})\mathrm{exp}(i\mathrm{\Pi }_i^{(0)}Q_i)$$
contains only one eigenvalue $`m/2`$ degenerated with respect to the cyclic momentum $`\mathrm{\Pi }_i`$. We see that there are differences between the dynamic and geometric descriptions. The dynamic evolution parameter is given in the whole region $`\mathrm{}<X_0<+\mathrm{}`$, whereas the geometric one is only positive $`0<T<+\mathrm{}`$, as it follows from the properties of the causal Green function (21) after the Levi-Civita transformation (27)
$$G^c(Q_\mu )=\underset{\mathrm{}}{\overset{+\mathrm{}}{}}d^4\mathrm{\Pi }_\mu \frac{\mathrm{exp}(iQ^\mu \mathrm{\Pi }_\mu )}{2m(\mathrm{\Pi }_0m/2iϵ/2m)}=\frac{\delta ^3(Q)}{2m}\theta (T),T=Q_0.$$
Two solutions of the constraint (a particle and antiparticle) in the dynamic system correspond to a single solution in the geometric system.
Thus, the reparametrization-invariant content of the equations of motion of a relativistic particle in terms of the geometric time is covered by two ”equivalent” unconstrained systems: the dynamic and geometric. In both the systems, the invariant times are not the coordinate evolution parameter, but variables with the negative contribution into the energy constraint. The Hamiltonian description of a relativistic particle in terms of the geometric time can be achieved by the Levi-Civita-type canonical transformation, so that the energy constraint converts into a new momentum. Whereas, the dynamic unconstrained system is suit for the secondary quantization and the derivation of the causal Green function that determine the arrow of the geometric time.
## 3. Relativistic String
### 3.1. The generalized Hamiltonian formulation
We begin with the action for a relativistic string in the geometrical form
$$W=\frac{\gamma }{2}d^2u\sqrt{g}g^{\alpha \beta }_\alpha x^\mu _\beta x_\mu ,u_\alpha =(u_0,u_1)$$
(35)
where the variables $`x_\mu `$ are string coordinates given in a space-time with a dimension $`D`$ and the metric $`(x_\mu x^\mu :=x_0^2x_i^2)`$; $`g_{\alpha \beta }`$ is a second-rank metric tensor given in the two-dimensional Riemannian space $`u_\alpha =(u_0,u_1)`$.
To formulate the Hamiltonian approach, one needs to separate the two-dimensional Riemannian space $`u_\alpha =(u_0,u_1)`$ on the set of space-like lines $`\tau =\mathrm{constant}`$ in the form of the Dirac-Arnovitt-Deser-Misner parametrization of the two-dimensional metric
$$g_{\alpha ,\beta }=\mathrm{\Omega }^2\left(\begin{array}{cc}\lambda _1^2\lambda _2^2& \lambda _2\\ \lambda _2& 1\end{array}\right),\sqrt{g}=\mathrm{\Omega }^2\lambda _1$$
(36)
with the invariant interval
$$ds^2=g_{\alpha \beta }du^\alpha du^\beta =\mathrm{\Omega }^2[\lambda _1^2d\tau ^2(d\sigma +\lambda _2d\tau )^2],u_\alpha =(u_0=\tau ,u_1=\sigma )$$
(37)
where $`\lambda _1`$ and $`\lambda _2`$ are known in general relativity (GR) as the lapse function and shift ” vector”, respectively . The action (35) after the substitution (37) does not depend on the conformal factor $`\mathrm{\Omega }`$ and takes the form
$$W=\frac{\gamma }{2}\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \left[\frac{(D_\tau x)^2}{\lambda _1}\lambda _1x^2\right]$$
(38)
where
$$D_\tau x_\mu =\dot{x}_\mu \lambda _2x_\mu ^{}(\dot{x}=_\tau x,x^{}=_\sigma x)$$
(39)
is the covariant derivative with respect to the two-dimensional metric (37). The metric (37), the action (38), and the covariant derivative (39) are invariant under the transformations (see Appendix A)
$$\tau \stackrel{~}{\tau }=f_1(\tau ),\sigma \stackrel{~}{\sigma }=f_2(\tau ,\sigma ).$$
(40)
A similar group of transformations in GR is well-known as the ”kinemetric” group of diffeomorphisms of the Hamiltonian description .
The variation of action (38) with respect to $`\lambda _1`$ and $`\lambda _2`$ leads to the equations
$$\frac{\delta W}{\delta \lambda _2}=\frac{x^{}D_\tau x}{\lambda _1}=0\lambda _2=\frac{\dot{x}x^{}}{x^2};$$
(41)
$$\frac{\delta W}{\delta \lambda _1}=\frac{(D_\tau x)^2}{\lambda _1^2}+x^2=0\lambda _1^2=\frac{(\dot{x}x^{})^2\dot{x}^2x^2}{(x^2)^2}$$
The solutions of these equations convert the action (38) into the standard Nambu-Gotto action of a relativistic string
$$W=\gamma \underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \sqrt{(\dot{x}x^{})^2\dot{x}^2x^2}.$$
The generalized Hamiltonian form is obtained by the Legendre transformation of the action (38)
$$W=\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \left(p_\mu D_\tau x^\mu +\lambda _1\varphi _1\right)=\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \left(p_\mu \dot{x}^\mu +\lambda _1\varphi _1+\lambda _2\varphi _2\right),$$
(42)
where
$$\varphi _1=\frac{1}{2\gamma }[p_\mu ^2+(\gamma x_\mu ^{})^2],\varphi _2=x^\mu p_\mu ,$$
(43)
and the generalized Hamiltonian
$$=\lambda _1\varphi _1+\lambda _2\varphi _2$$
(44)
is treated as the generator of evolution with respect to the coordinate time $`\tau `$, and $`\lambda _1,\lambda _2`$ play the role of variables with the zero momenta
$$P_{\lambda _1}=0,P_{\lambda _2}=0$$
(45)
considered as the first class primary constraints . The equations for $`\lambda _1,\lambda _2`$
$$\frac{\delta W}{\delta \lambda _1}=\varphi _1=0;\frac{\delta W}{\delta \lambda _2}=\varphi _2=0$$
(46)
are known as the first class secondary constraints . The Hamiltonian equations of motion take the form
$$\frac{\delta W}{\delta x^\mu }=\dot{p}_\mu _\sigma [\gamma \lambda _1x_\mu ^{}+\lambda _2p_\mu ]=0,\frac{\delta W}{\delta p^\mu }=p_\mu +\gamma \frac{D_\tau x_\mu }{\lambda _1}=0$$
(47)
The problem is to find solutions of the Hamiltonian equations of motion (47) and constraints (46) which are invariant with respect to the kinemetric transformations (40).
There is the problem of the solution of the linearized ”gauge-fixing” equation in terms of the evolution parameter $`\tau `$ (as the object reparametrizations in the initial theory) being adequate to the initial kinemetric invariant and relativistic invariant system. In particular, the constraints mix the global motion of the ”center of mass” coordinates with local excitations of a string $`\xi _\mu `$, which contradicts to the relativistic invariance of internal degrees of freedom of a string. In this context, it is worth to clear up a set of questions: Is it possible to introduce the reparametrization-invariant evolution parameter for the string dynamics, instead of the non-invariant coordinate time $`(\tau )`$ used as the evolution parameter in the gauge-fixing method? Is it possible to construct the observable nonzero Hamiltonian of evolution of the ”center of mass” coordinates? What is relation of the ”center of mass” evolution to the unitary representations of the Poincare group?
## 4. The separation of the ”center of mass” coordinates
To apply the reparametrization-invariant Hamiltonian reduction discussed before to a relativistic string, one should define the proper time in the form of the reparametrization-invariant functional of the lapse function (of type (14)), and to point out, among the variables, a dynamic evolution parameter, the equation of which identifies it with the proper time of type (8). As any extended object admits to define the coordinates of its center of mass, we identify this dynamic evolution parameter with the time-like coordinate of the center of mass of a string
$$X_\mu (\tau )=\frac{1}{l(\tau )}\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma x_\mu (\tau ,\sigma ),l(\tau )=\sigma _2(\tau )\sigma _1(\tau ).$$
(48)
We see that the invariant reduction requires to separate the ”center of mass” variables before variation of the action. This separation is fulfilled by the substitution of
$$x_\mu (\tau ,\sigma )=X_\mu (\tau )+\xi _\mu (\tau ,\sigma )$$
(49)
into the action (38), which takes the form
$$W=\frac{\gamma }{2}\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \left\{\frac{\dot{X}^2l(\tau )}{N_0(\tau )}+2\dot{X}_\mu \underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \frac{D_\tau \xi ^\mu }{\lambda _1}+\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \left(\frac{(D_\tau \xi )^2}{\lambda _1}\lambda _1\xi ^2\right)\right\},$$
(50)
where the global lapse function $`N_0(\tau )`$ is defined as the functional of $`\lambda _1(\tau ,\sigma )`$
$$\frac{1}{N_0[\lambda _1]}=\frac{1}{l(\tau )}\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \frac{1}{\lambda _1(\tau ,\sigma )}.$$
(51)
¿From definition (48) and equality (53) it follows that the local variables $`\xi _\mu `$ are given in the class of functions (with the nonzero Fourier harmonics) which satisfy the conditions
$$\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \xi _\mu (\tau ,\sigma )=0.$$
(52)
The formulation of the Hamiltonian approach (consistent with (48)) supposes the similar separation of the conjugate momenta $`p_\mu `$ defined by equation (47). If we substitute the definition (53) in these equations, we get
$$p_\mu (\tau ,\sigma )=\gamma \left(\frac{\dot{X}_\mu (\tau )}{\lambda _1}+\frac{D_\tau \xi _\mu (\tau ,\sigma )}{\lambda _1}\right).$$
(53)
Defining the total momentum of a string $`P_\mu `$
$$P_\mu =\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma p_\mu (\tau ,\sigma )=\gamma \underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \left(\frac{\dot{X}_\mu (\tau )}{\lambda _1}+\frac{D_\tau \xi _\mu (\tau ,\sigma )}{\lambda _1}\right),$$
(54)
and taking into account (51) we obtain the following expresion
$$P_\mu =\gamma \frac{\dot{X}_\mu l}{N_0(\tau )}\gamma \underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \frac{D_\tau \xi _\mu (\tau ,\sigma )}{\lambda _1},$$
(55)
therefore the equality
$$\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \frac{D_\tau \xi ^\mu }{\lambda _1}=\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \pi _\mu (\tau ,\sigma )=0.$$
(56)
should be valid. This separation conserves the group of diffeomorphisms of the Hamiltonian and leads to the Bergmann-Dirac generalized action
$$W=\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \left[\left(\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma [\pi _\mu D_\tau \xi ^\mu \lambda _1]\right)P_\mu \dot{X}^\mu +N_0\frac{P_\mu ^2}{2\overline{\gamma }}\right],(\overline{\gamma }=\gamma l(\tau ))$$
(57)
where $``$ is the Hamiltonian of local excitations
$$=\frac{1}{2\gamma }[\pi _\mu ^2+(\gamma \xi _\mu ^{})^2].$$
(58)
The variation of the action (57) with respect to $`\lambda _1`$ results in the equation
$$\frac{\delta W}{\delta \lambda _1}=\left(\frac{1}{l\overline{\lambda }_1}\right)^2\frac{P^2}{2\gamma }=0,$$
(59)
where
$$\overline{\lambda }_1(\tau ,\sigma )=\frac{\lambda _1(\tau ,\sigma )}{N_0(\tau )}$$
(60)
is the reparametrization-invariant component of the local lapse function. Here we have used the variation of the functional $`N_0[\lambda _1]`$ (51)
$$\frac{\delta N_0[\lambda _1]}{\delta \lambda _1}=\frac{1}{l(\tau )\overline{\lambda }_1^2}.$$
In accordance with our separation of dynamic variables onto the global and local sectors, the first class constraint (59) has two projections onto the global sector (zero Fourier harmonic) and the local one. The global part of the constraint (59) can be obtained by variation of the action (57) with respect to $`N_0`$ (after the substitution of (60) into (57))
$$\frac{\delta W}{\delta N_0}=\frac{P^2}{2\overline{\gamma }}H=0,H=\underset{\sigma _1}{\overset{\sigma _2}{}}𝑑\sigma \overline{\lambda }_1,$$
(61)
or, in another way, by the integration over $`\sigma `$ of (59) multiplied by $`\lambda _1`$. Then, the local part of the constraint (59) can be obtained by the substitution of (61) into (59)
$$\overline{\lambda }_1\frac{1}{l\overline{\lambda }_1}\underset{\sigma _1}{\overset{\sigma _2}{}}𝑑\sigma \overline{\lambda }_1=0.$$
(62)
The integration of the local part over $`\sigma `$ is equal to zero if we take into account the normalization of the local lapse function
$$\frac{1}{l(\tau )}\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \frac{1}{\overline{\lambda }_1}=1.$$
(63)
This follows from the definition of the global lapse function (51).
Finally, we can represent the action (57) in the equivalent form
$$W=\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \left[\left(\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma [\pi _\mu D_\tau \xi ^\mu ]\right)P_\mu \dot{X}^\mu N_0(\frac{P_\mu ^2}{2\overline{\gamma }}+H)\right],$$
(64)
where the global lapse function $`N_0`$ and the local one $`\overline{\lambda }_1`$ are treated as independent variables, with taking the normalization (63) into account after the variation.
According to (40) and (51) the invariant proper time $`T`$ measured by the watch of an observer in the ”center of mass” frame of a string is given by the expression
$$\sqrt{\gamma }dT:=N_0d\tau ,\sqrt{\gamma }T=\underset{0}{\overset{\tau }{}}𝑑\tau ^{}\left[\frac{1}{l(\tau ^{})}\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \frac{1}{\lambda _1(\tau ^{},\sigma )}\right]^1.$$
(65)
We include the constant $`\sqrt{\gamma }`$ to provide the dimension of the time measured by the watch of an observer.
Now we can see from (64) that the dynamics of the local degrees of freedom $`\pi ,\xi `$, in the class of functions of nonzero harmonics (52), is described by the same kinemetric invariant and relativistic covariant equations (47) where $`x,p`$ are changed by $`\xi ,\pi `$, with the set of the first class (primary and secondary) constraints
$$P_{\lambda _1}=0,P_{\lambda _2}=0,\pi _\mu \xi ^\mu =0,\overline{\lambda }_1\frac{1}{l\overline{\lambda }_1}\underset{\sigma _1}{\overset{\sigma _2}{}}𝑑\sigma \overline{\lambda }_1=0.$$
(66)
The separation of the ”center of mass” (CM) variables on the level of the action removes the interference terms which mix the CM variables with the local degrees of freedom; as a result, the new local constraints (66) do not depend on the total momentum $`P_\mu `$, in contrast to the standard ones. In other words, there is the problem: when can one separate the CM coordinates of a relativistic string; before the variation of the action or after the variation of the action? The relativistic invariance dictates the first one, because an observer in the CM frame (which is the preferred frame for a string) cannot measure the total momentum of the string.
The first class local constraints (66) can be supplemented by the second class constraints
$$\overline{\lambda }_11=0,\lambda _2=0,n^\mu \xi _\mu =0,n^\mu \pi _\mu =0,$$
(67)
where $`n_\mu `$ is an arbitrary time-like vector. In particular, for $`(n_\mu =(1,0,0,0))`$ the equations of the local constraint-shell action
$$W(\text{loc.constrs.})=\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \left[\left(\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma \pi _i\dot{\xi }_i\right)P_\mu \dot{X}^\mu N_0(\frac{P_\mu ^2}{2\overline{\gamma }}+H)\right]$$
(68)
coincide with the complete set of equations and the same constraints (66), (67) of the extended action, i.e., the operations of constraining and variation commute. The substitution of the global constraint (61) with $`\overline{\lambda }_1=1`$ into the action (68) leads to the constraint-shell action
$$W_\pm ^D=\underset{X_0(\tau _1)}{\overset{X_0(\tau _2)}{}}𝑑X_0\left[\left(\underset{\sigma _1(X_0)}{\overset{\sigma _2(X_0)}{}}𝑑\sigma \pi _i\frac{d\xi _i}{dX_0}\right)+P_i\frac{dX_i}{dX_0}\sqrt{P_i^2+2\overline{\gamma }H}\right].$$
(69)
This action describes the dynamics of a relativistic string with respect to the time measured by an observer in the rest frame with the physical nonzero Hamiltonian of evolution. However, in this system, equations become nonlinear. To overcome this difficulty, we pass to the ”center of mass” frame.
## 5. Levi-Civita geometrical reduction of a string
To express the dynamics of a relativistic string in terms of the proper time (65) measured by an observer in the comoving (i.e. ”center of mass”) frame, we use the Levi-Civita-type canonical transformations (as in Section 2.3)
$$(P_\mu ,X_\mu )(\mathrm{\Pi }_\mu ,Q_\mu );$$
they convert the global part of the constraint (61) into a new momentum $`\mathrm{\Pi }_0`$
$$\mathrm{\Pi }_0=\frac{1}{2\overline{\gamma }}[P_0^2P_i^2],\mathrm{\Pi }_i=P_i,Q_0=X_0\frac{\overline{\gamma }}{P_0},Q_i=X_iX_0\frac{P_i}{P_0}.$$
(70)
The inverted form of these transformations is
$$P_0=\pm \sqrt{2\overline{\gamma }\mathrm{\Pi }_0+\mathrm{\Pi }_i^2},P_i=\mathrm{\Pi }_i,X_0=\pm Q_0\frac{\sqrt{2\overline{\gamma }\mathrm{\Pi }_0+\mathrm{\Pi }_i^2}}{\overline{\gamma }},X_i=Q_i+Q_0\frac{\mathrm{\Pi }_i}{\overline{\gamma }}.$$
(71)
As a result of transformations (70), the extended action (64) in terms of the Levi-Civita geometrical variables takes the form (compare with (1))
$$W=\underset{\tau _1}{\overset{\tau _2}{}}𝑑\tau \left[\left(\underset{\sigma _1(\tau )}{\overset{\sigma _2(\tau )}{}}𝑑\sigma [\pi _\mu D_\tau \xi ^\mu ]\right)\mathrm{\Pi }_\mu \dot{Q}^\mu N_0(\mathrm{\Pi }_0+H)\frac{d}{d\tau }(Q_0\mathrm{\Pi }_0)\right].$$
(72)
The Hamiltonian reduction means to resolve constraint (61) with respect to the momentum $`\mathrm{\Pi }_0`$
$$\frac{\delta W}{\delta N_0}=0\mathrm{\Pi }_0=H.$$
(73)
The equation of motion for the momentum $`\mathrm{\Pi }_0`$
$$\frac{\delta W}{\delta \mathrm{\Pi }_0}=0\frac{dQ_0}{d\tau }=N_0(i.e.,dQ_0=N_0d\tau :=\sqrt{\gamma }dT)$$
(74)
identifies (according to our definition (65)) the new variable $`Q_0`$ with the proper time $`T`$, whereas the equation for $`Q_0`$
$$\frac{\delta W}{\delta Q_0}=0\frac{d\mathrm{\Pi }_0}{d\tau }=0,i.e.,\frac{dH}{dT}=0,$$
(75)
in view of (73), gives us the conservation law.
Thus, resolving the global energy constraint $`\mathrm{\Pi }_0=H`$, we obtain, from (72), the reduced action for a relativistic string in terms of the proper time $`T`$
$$W^G=\underset{T_1}{\overset{T_2}{}}𝑑T\left[\left(\underset{\sigma _1}{\overset{\sigma _2}{}}𝑑\sigma [\pi _\mu D_T\xi ^\mu ]\right)+\mathrm{\Pi }_i\frac{dQ_i}{dT}H\frac{d}{dT}(TH)\right],$$
(76)
where in analogy with (60) we introduced the factorized ”shift-vector” $`\lambda _2=N_0\overline{\lambda }_2/\sqrt{\gamma }`$; in this case the covariant derivative (39) takes the form
$$D_T\xi _\mu =_T\xi _\mu \overline{\lambda }_2\xi _\mu ^{}=\frac{D_\tau \xi _\mu }{N_0}\sqrt{\gamma }.$$
(77)
The reduced system (76) has trivial solutions for the global variables $`\mathrm{\Pi }_i,Q_i`$
$$\frac{\delta W^R}{\delta \mathrm{\Pi }_i}=0\frac{dQ_i}{dT}=0;Q_i=\text{const};$$
(78)
$$\frac{\delta W^R}{\delta Q_i}=0\frac{d\mathrm{\Pi }_i}{dT}=0,\mathrm{\Pi }_i=\text{const}$$
which have the meaning of initial data.
If the solutions of equations (73), (74), and (78) for the system (76)
$$\mathrm{\Pi }_0=H:=\frac{M^2}{2\overline{\gamma }},\mathrm{\Pi }_i=P_i,Q_0=T\sqrt{\gamma },Q_i=X_i(0),$$
(79)
are substituted into the inverted Levi-Civita canonical transformations (71)
$$P_0=\pm \sqrt{M^2+P_i^2},X_0(T)=T\frac{P_0}{\sqrt{\gamma }l},X_i(T)=Q_i+T\frac{P_i}{\sqrt{\gamma }l},$$
(80)
the initial extended action (64) can be described in the rest frame of an observer who measures the energy $`P_0`$ and the time $`X_0`$ and sees the rest frame evolution of the ”center of mass” coordinates
$$X_i(X_0)=Q_i+X_0\frac{P_i}{P_0}.$$
(81)
The Lorentz scheme of describing a relativistic system in terms of the time and energy $`(X_0,P_0)`$ in the phase space $`P_i,X_i,\pi _\mu ,\xi _\mu `$ is equivalent to the above-considered the Levi-Civita scheme in terms of the proper time and the evolution Hamiltonian $`(T,H)`$ in the phase space $`\mathrm{\Pi }_i,Q_i,\pi _\mu ,\xi _\mu `$, where the variables $`\mathrm{\Pi }_i,Q_i`$ are cyclic.
## 6. Dynamics of the local variables
### 6.1. Reparametrization - invariant reduction for an open string
We restrict ourselves to an open string with the boundary conditions
$$\sigma _1(T)=0,\sigma _2(T)=\pi ,l(T)=\pi .$$
(82)
In the gauge-fixing method, by using the kinemetric transformation, we can put
$$\overline{\lambda }_1=1,\overline{\lambda }_2=0.$$
(83)
This requirement does not contradict the normalization of $`\overline{\lambda }_1`$ (63).
In view of (66), it means that the reduced Hamiltonian $`H`$ (61) coincides with its density (58)
$$\overline{\varphi }_1=\frac{1}{\pi }\underset{0}{\overset{\pi }{}}𝑑\sigma =0,\overline{\varphi }_2=\pi _\mu \xi ^\mu =0$$
(84)
In this case, the reparametrization-invariant equations for the local variables obtained by varying the action (76)
$$\frac{\delta W_s^R}{\delta \xi ^\mu }=0_T\pi _\mu _\sigma (\overline{\lambda }_2\pi _\mu )=\gamma _\sigma (\overline{\lambda }_1\xi _\mu ^{}),\frac{\delta W_s^R}{\delta \pi ^\mu }=0\gamma D_T\xi _\mu =\overline{\lambda }_1\pi _\mu $$
(85)
lead to the D’Alambert equations
$$_T^2\xi _\mu _\sigma ^2\xi _\mu =0.$$
(86)
The general solution of these equations of motion in the class of functions (52) with the boundary conditions (82) is given by the Fourier series
$$\xi _\mu (T,\sigma )=\frac{1}{2\sqrt{\pi \gamma }}[\psi _\mu (z_+)+\psi _\mu (z_{})],\psi _\mu (z)=i\underset{n0}{}e^{(inz)}\frac{\alpha _{n\mu }}{n},z_\pm =T\sqrt{\gamma }\pm \sigma .$$
(87)
$$\xi _\mu ^{}(T,\sigma )=\frac{1}{2\sqrt{\pi \gamma }}[\psi _\mu ^{}(z_+)\psi _\mu ^{}(z_{})],\pi _\mu (T,\sigma )=\frac{1}{2}\sqrt{\frac{\gamma }{\pi }}[\psi _\mu ^{}(z_+)+\psi _\mu ^{}(z_{})].$$
The total coordinates $`Q_\mu ^{(0)}`$ and momenta $`P_\mu `$ are determined by the reduced dynamics of the ”center of mass” (78), (79), (80), and the string mass $`M`$ obtained from (61)
$$P_\mu ^2=M^2=2\pi \gamma H=2\pi \gamma \underset{0}{\overset{\pi }{}}𝑑\sigma .$$
(88)
The substitution of $`\xi _\mu `$ and $`\pi _\mu `$ from (87) into (58) leads to the Hamiltonian density
$$=\frac{1}{4\pi }\left[\psi _\mu ^2(z_+)+\psi _\mu ^2(z_{})\right],$$
and from (88) we obtain, for the mass, the expression
$$M^2=2\pi \gamma \overline{L}_0=\frac{\gamma }{2}\underset{0}{\overset{\pi }{}}𝑑\sigma \left[(\psi _\mu ^{}(z_+))^2+(\psi _\mu ^{}(z_{}))^2\right].$$
(89)
The second constraint (84) in terms of the vector $`\psi _\mu ^{}`$ in (87) takes the form
$$\xi _\mu ^{}\pi ^\mu =\frac{1}{4\pi }\left[\psi _\mu ^2(z_+)\psi _\mu ^2(z_{})\right]=0\psi _\mu ^2(z_+)=\psi _\mu ^2(z_{})=\text{const.},$$
(90)
and the first constraint (84) $`\overline{\varphi }_1=0`$ is satisfied identically. After the substitution of the constant value (90) into (89) we obtain that $`\text{const.}=M^2/\pi \gamma `$; thus, finally the reparamerization-invariant constraint takes the form
$$P_\mu ^2+\pi \gamma \psi _\mu ^2(z_\pm )=0(P_\mu ^2=M^2).$$
(91)
Unlike this constraint, the gauge-fixing reparametrization-noninvariant constraint
$$\left(P_\mu +\sqrt{\pi \gamma }\psi _\mu ^{}\right)^2=0$$
(92)
contains the interference of the local and global degrees of freedom $`\psi _\mu ^{}P^\mu `$. The latter violates the relativistic invariance of the local excitations which form the mass and spin of a string.
The constraint (91) in terms of the Fourier components (87) takes the form
$$\psi _\mu ^2(z_\pm )=\underset{k,m0}{}\alpha _{k,\mu }\alpha _m^\mu e^{i(k+m)z_\pm }=2\underset{n=\mathrm{}}{\overset{\mathrm{}}{}}\overline{L}_ne^{inz_\pm }=\frac{M^2}{\pi \gamma },$$
(93)
where $`\overline{L}_n`$ are the contributions of the nonzero harmonics
$$\overline{L}_0=\frac{1}{2}\underset{k0}{}\alpha _{k\mu }\alpha _k^\mu ,\overline{L}_{n0}=\frac{1}{2}\underset{k0,n}{}\alpha _{k\mu }\alpha _{nk}^\mu .$$
(94)
From (93) one can see that the zero harmonic of this constraint determines the mass of a string
$$M^2=2\pi \gamma \overline{L}_0=\pi \gamma \underset{k0}{}\alpha _{k\mu }\alpha _{k\mu }$$
(95)
and coincides with the gauge-fixing value. However, the nonzero harmonics of constraint (93)
$$\overline{L}_{n0}=\frac{1}{2}\underset{k0,n}{}\alpha _{k\mu }\alpha _{nk\mu }=0,\overline{L}_n=\overline{L}_n^{}$$
(96)
(as we dicussed above) strongly differ from the nonzero harmonics of the gauge-fixing constraints (92). The latter (in the contrast to (91)) contains the mixing the global motion of the center of mass $`P_\mu `$ with the local excitations $`\psi _\mu `$. It is clear that this mixing the global and local motions violates the the Poincare invariance of the local degrees of freedom.
The algebra of the local constraints (96) of the reparametrization-invariant dynamics of a relativistic string is not closed, as it does not contains the zero Fourier harmonic of the energy constraint (which has been resolved to express the dynamic equations in terms of the proper time).
The reparametrization-invariant dynamics of a relativistic string in the form of the first and second class constraints (66), (67) coincides with the Röhrlich approach to the string theory . This approach is based on the choice of the gauge condition
$$p_\mu \xi ^\mu =0,p_\mu \pi ^\mu =0G_n=P_\mu \alpha _n^\mu =0,n0,$$
instead of (67). As consequence of this gauge the constraints (91), (92) became equivalent. In quantum theory, this condition is used for eliminating the states with negative norm in the ”center of mass” (CM) frame (in our scheme, the CM frame appears as a result of the geometric Levi-Civita reduction). This reference frame is the only preferred frame for quantizing such a composite relativistic object as the string, as only in this frame one can quantize the initial data. This is a strong version of the principle of correspondence with classical theory: the classical initial data become the quantum numbers of quantum theory.
### 6.2. Quantum theory
Thus, our classical Hamiltonian reparametrization - invariant formalism provides the quantization of the string as in Röhlich gauge.
The Röhrlich approach distinguishes two cases: $`M^2=0`$ and $`M^20`$.
The first case, in our scheme, the equality $`M^2=0`$ together with the local constraints (96) form the Virasoro algebra. The reparametrization-invariant version of the Virasoro algebra (with all its difficulties, including the $`D=26`$ \- problem and the negative norm states) appears only in the case of the massless string $`2\pi \gamma \overline{L}_0=M^2=0`$.
In the second case $`M^20`$, the Röhrlich gauge $`\alpha _{n,0}=0`$. allows us to exclude the time Fourier components $`\alpha _{n0}`$, and it is just these components that after quantization
$$[\alpha _{n,\mu },\alpha _{n,\nu }^+]=m\eta _{\mu ,\nu }\delta _{m,n};(n,m0,\eta _{00}=\eta _{ii}=1)$$
lead to the states with negative norm because of the system being unstable. This means that the state vectors in the CM frame are constructed only by the action on vacuum of the spatial components of the operators $`a_{ni}^+=\alpha _{ni}/\sqrt{n},n>0`$
$$|𝚽_\nu >_{CM}=\underset{n=1}{\overset{\mathrm{}}{}}\frac{(a_{nx}^+)^{\nu _{nx}}}{\sqrt{\nu _{nx}!}}\frac{(a_{ny}^+)^{\nu _{ny}}}{\sqrt{\nu _{ny}!}}\frac{(a_{nz}^+)^{\nu _{nz}}}{\sqrt{\nu _{nz}!}}|0>,$$
(97)
where the three-dimensional vectors $`\nu _n=(\nu _{nx},\nu _{ny},\nu _{nz})`$ have only nonnegative integers as components. These state vectors automatically satisfy the constraint
$$\alpha _{n0}|𝚽_\nu >_{CM}=0,n>0$$
(98)
The physical states (97) are subjected to further constraints (96) with $`n0`$
$$\overline{L}_n|𝚽_\nu >_{CM}=0,n>0,P^2=M_\nu ^2=\pi \gamma <𝚽_\nu \underset{m0}{}\alpha _{m,i}\alpha _{m,i}|𝚽_\nu >,$$
(99)
where $`\overline{L}_n`$ can be represented in the normal ordering form
$$\overline{L}_{n>0}=\underset{k=1}{\overset{\mathrm{}}{}}\alpha _{k,i}^+\alpha _{n+k,i}+\frac{1}{2}\underset{k>1}{\overset{n1}{}}\alpha _{k,i}\alpha _{nk,i}.$$
(100)
Constraints (98) and (99) are the first class constraints, in accordance with the Dirac classification because they form a closed algebra for $`n,m>0`$
$$[G_n,G_m]=0,[\overline{L}_n,\overline{L}_m]=(nm)\overline{L}_{n+m},[G_n,\overline{L}_m]=nG_{m+n}.$$
(101)
Therefore the conditions (98) eliminating the ghosts and the conditions (99) defining the physical vector states are consistent. Note that the commutator $`[\overline{L}_n,\overline{L}_m]`$ does not contain a c-number since suffices $`n0`$ and $`m0`$ in the Virasoso operators $`\overline{L}_n`$ do not lead to the central term.
On the operator level, equations determining the resolution of the constraints are fulfilled in a weak sense, as only the ” annihilation” part of the constraints is imposed on the state vectors.
In quantum theory, one can introduce a complete set of eigen functions satisfing equations
$$H[\pi _i,\xi _i]<\xi |\nu >=\frac{M_\nu ^2}{2\pi \gamma }<\xi |\nu >,$$
(102)
where
$$<\xi |\nu >=<\xi |𝚽_\nu >,\underset{\nu }{}<\xi _1|\nu ><\nu |\xi _2>=\underset{\sigma }{}\delta ^3(\xi _1\xi _2).$$
## 7. The causal Green functions
Now we can construct the causal Green function for a relativistic string as the analogy of the causal Green function for a relativistic particle (23) - (25) discussed in Section 2.
The Veneziano-type causal Green function is the spectral series with the Hermite polynomials $`<\xi |\nu >`$ over the physical state vectors $`|𝚽_\nu >=|\nu >`$
$$G_c(X|\xi _1,\xi _2)=G_+(X|\xi _1,\xi _2)\theta (X_0)+G_{}(X|\xi _1,\xi _2)\theta (X_0)=$$
(103)
$$i\frac{d^4P}{(2\pi )^4}\mathrm{exp}(iPX)\underset{\nu }{}\frac{<\xi _1|\nu ><\nu |\xi _2>}{P^2M_\nu ^2iϵ}.$$
The commutative Green function for a relativistic string $`G_+(X|\xi _1,\xi _2)`$ can be represented in the form of the Faddeev-Popov functional integral in the local gauge (67)
$$G_+(X|\xi _2,\xi _1)=\underset{X(\tau _1)=0}{\overset{X(\tau _2)=X}{}}\frac{dN_0(\tau _2)d^4P(\tau _2)}{(2\pi )^3}\underset{\tau _1\tau <\tau _2}{}\left\{d\overline{N}_0(\tau )\underset{\mu }{}\left(\frac{dP_\mu (\tau )dX_\mu (\tau )}{2\pi }\right)\right\}F_+(\xi _2,\xi _1),$$
(104)
using the representation of the spectral series
$$F_+(\xi _2,\xi _1)=\underset{\nu }{}<\xi _2|\nu >\mathrm{exp}\left\{iW[P,X,N_0,M_\nu ]\right\}<\nu |\xi _1>=$$
(105)
in the form of the functional integral
$$F_+(\xi _2,\xi _1)=\underset{\xi _1}{\overset{\xi _2}{}}D(\xi ,\pi )\mathrm{\Delta }_{fp}\mathrm{exp}\left\{iW_{fp}\right\},$$
$`W[P,X,N_0,M_\nu ]`$ is the action (11) with the mass $`M_\nu `$
$$W_{fp}=\underset{0}{\overset{\tau (X_0)}{}}𝑑\tau \left[\left(\underset{0}{\overset{\pi }{}}𝑑\sigma \pi _\mu \dot{\xi }^\mu \right)P_\mu \dot{X}^\mu N_0\left(\frac{P^2}{2\pi \gamma }+H\right)\right]$$
(106)
is the constraint-shell action (68),
$$D(\xi ,\pi )=\underset{\tau ,\sigma }{}\underset{\mu }{}\frac{d\xi _\mu d\pi _\mu }{2\pi },$$
(107)
and
$$\mathrm{\Delta }_{fp}=\underset{\tau ,\sigma }{}\delta (\varphi _1))\delta (\pi _0)\delta (\varphi _2))\delta (\xi _0)detB^1,detB=det\{\varphi _1,\varphi _2,\pi _0,\xi _0\}$$
(108)
is the FP determinant given in the monograph .
## 8. Conclusion
To describe the invariant dynamics of constrained relativistic string we used the universal method of the Hamiltonian reduction of their actions by resolving the energy constraint, so that one of variables of the extended phase space (with a negative contribution to the energy constraint) converts into the invariant evolution parameter, and its conjugate momentum becomes the invariant Hamiltonian of evolution.
This method allows us to find integrals of motion by the Levi-Civita canonical transformations which converts the energy constraint into a new momentum, and the time-like variable of the world space into the proper time interval. For a particle and a string the Levi-Civita transformations are the Hamiltonian form of the Lorentz transformations which describe pure relativistic effects of the transition from the rest frame of reference to the comoving one.
We have shown that a relativistic string can be described directly in terms of the reparametrization-invariant parameter of evolution with the nonzero Hamiltonians of evolution in agreement with the equations of motion of the initial system.
A crucial point in our approach is the separation of the ”center of mass” coordinates on the level of the action. The definition of the proper time with the nonzero Hamiltonian of evolution consistent with the group of diffeomorphisms of the Hamiltonian description requires to separate the ”center of mass” coordinates before varying the action, whereas in the standard gauge-fixing method, the ”center of mass” coordinates are separated after varying the action. The operations of separation of the ”center of mass” coordinates and variation of the action do not commute. The relativistic invariance dictates the reparametrization - invariant way, as an observer in the comoving frame cannot measure the componets of the total momentum of a string. Unique admissible gauge is the Röhrlich gauge that leads directly to the quantum theory of a string without a critical dimension.
Thus, we can formulate the novelty of this work: i) the separation of the ”center of mass” coordinates on the level of the action, ii) finding of the integrals of motion by the Levi-Civita transformation, iii) deriving of the nonzero Hamiltonian of evolution of a string with respect to the proper time with the new algebra of the Poisson brackets, that provides the Rörlich gauge, and iv) constructing of new reparametrization - invariant path integral representations of the causal Green functions for relativistic particle and string.
Acknowledgments
We are happy to acknowledge interesting and critical discussions with H. Kleinert, E. A. Kuraev, V. V. Nesterenko, A. I. Pashnev, and I. V. Tyitin.
Appendix A: Kinemetric transformations
The kinemetric transformations of the differentials
$$\stackrel{~}{\tau }=\dot{f}_1(\tau )d\tau ,d\stackrel{~}{\sigma }=\dot{f}_2(\tau ,\sigma )d\tau +f_2^{}(\tau ,\sigma )d\sigma $$
correspond to transformations of the string coordinates
$$x_\mu (\tau ,\sigma )=\stackrel{~}{x}_\mu (\stackrel{~}{\tau }),\stackrel{~}{\sigma }),x^{}_\mu (\tau ,\sigma )=\stackrel{~}{x}^{}_\mu (\stackrel{~}{\tau },\stackrel{~}{\sigma })f^{}_2(\tau ,\sigma ),$$
$$\dot{x}_\mu (\tau ,\sigma )=\dot{\stackrel{~}{x}}_\mu (\stackrel{~}{\tau },\stackrel{~}{\sigma })\dot{f}_1(\tau )+\stackrel{~}{x}_\mu ^{}(\stackrel{~}{\tau },\stackrel{~}{\sigma })\dot{f}_2(\tau ,\sigma ),$$
¿From these equations, we can derive the transformation law for $`\lambda _1,\lambda _2`$ taking into account (41)
$$\lambda _1(\tau ,\sigma )=\frac{\sqrt{(\dot{x}x^{})^2\dot{x}^2x^2}}{x^2(\tau ,\sigma )}=\frac{\sqrt{(\dot{\stackrel{~}{x}}\stackrel{~}{x}^{})^2\dot{\stackrel{~}{x}}^2\stackrel{~}{x}^2}}{\stackrel{~}{x}^2(\stackrel{~}{\tau },\stackrel{~}{\sigma })}\frac{\dot{f}_1}{f_2^{}}=\stackrel{~}{\lambda }_1\frac{\dot{f}_1(\tau )}{f_2^{}(\tau ,\sigma )}.$$
$$\lambda _2(\tau ,\sigma )=\frac{\dot{x}x^{}}{x^2}=\frac{(\dot{\stackrel{~}{x}}\stackrel{~}{x}^{})\dot{f}_1f_2^{}+\stackrel{~}{x}^2\dot{f}_2f_2^{}}{\stackrel{~}{x}^2f_2^2}=\stackrel{~}{\lambda }_2\frac{\dot{f}_1}{f_2^{}}+\frac{\dot{f}_2}{f_2^{}}.$$
The kinemetric-invariance of the interval (37) with respect to (40) follows from these transformation laws and the transformation of the conformal factor
$$\mathrm{\Omega }(\tau ,\sigma )=f_2^{}(\tau ,\sigma )\stackrel{~}{\mathrm{\Omega }}(\stackrel{~}{\tau },\stackrel{~}{\sigma })$$
The covariant derivative (39) is transformed under (40) as
$$D_\tau x_\mu =\dot{x}_\mu \lambda _2x_\mu ^{}=\dot{f}_1(\tau )\left[\dot{\stackrel{~}{x}}_\mu \stackrel{~}{\lambda }_2\stackrel{~}{x}_\mu ^{}\right]=\dot{f}_1(\tau )D_{\stackrel{~}{\tau }}\stackrel{~}{x}_\mu .$$
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# The newest release of the Ortocartan set of programs for algebraic calculations in relativity
## 1 The motivation
Several programs for algebraic calculations are available on the market today, and some of them include packages for relativity. Still, it seems that not all possibilities in applying such programs are known to the users and to the authors of those systems. The newest developments in Ortocartan include a few facilities that are not yet in general use. It is hoped that the present paper will demonstrate what is still possible beyond the standard applications.
## 2 Calculating the curvature tensors from a given metric
This has been the main application of the program Ortocartan ever since its first release in 1978. It was described in several publications, most recently in Ref. 1, so it will be recalled only very briefly. All details can be found in Ref. 1, in the papers cited there, and in the latest release of the user’s manual .
The main program Ortocartan takes an orthonormal tetrad of Cartan forms as its input data, and calculates all quantities that appear in the usual calculations in relativity: the determinant of this tetrad, the components of the inverse tetrad, the metric tensor and its inverse, the Ricci rotation coefficients, the Christoffel symbols, the Riemann, Ricci, Einstein and Weyl tensors, the scalar curvature. For the tensors calculated along the way it finds the (orthonormal) tetrad components and the coordinate components with each index raised or lowered as the user wishes. The user can make all kinds of substitutions, including those by pattern-matching (for an example of the latter see sec. 8). As an example, the input data for a very simple application of Ortocartan (calculating the curvature tensors for the spherically symmetric metric in the curvature coordinates) is shown below.
This example, and all the other examples shown in this paper, were chosen trivial for illustrative purposes, so that the readers can easily see what happens and how. In actual research, all the programs were successfully applied also to very complicated examples.
```
(ortocartan ’(
(SPHERICAL METRIC IN STANDARD FORM)
(functions mu(t r) nu(t r))
(coordinates t r theta phi)
(ematrix (exp nu) 0 0 0 0 (exp mu) 0 0 0 0 r
0 0 0 0 (r *(sin theta)))
(stop after ricci)
(rmargin 60)
))
```
(The command (rmargin 60) will fit the linelength of the output shown in the next section to the present text.)
As before, the Ortocartan package includes the program Calculate that can carry out simple algebraic operations defined by the user. This has been described in the earlier publications, too, so again it is just recalled by the example below, in which it calculates the derivative by $`x`$ of a complicated expression that includes constants and functions. The input data is:
```
(calculate ’(
(print example)
(constants a b c d)
(coordinates x)
(functions f (x))
(operation (deriv x (a ^ (b ^ ((der x f) ^ (c ^ d))))))
))
```
and the output is<sup>1</sup><sup>1</sup>1Each piece of output shown further is a part of a disk-output from the appropriate program that was inserted verbatim in the latex code of this paper.:
```
(print example)
(I UNDERSTAND YOU REQUEST THE FOLLOWING EXPRESSION TO BE SIMPLIFIED)
d
c
f,
x
b
> deriv (x,a )
```
```
THE RESULT IS
d
c d
f, c d
x
b f, d - 1 + c
x
> result = a b c log (a) log (b) f, f,
1 x x x
(I REALLY LIKED THIS ! CAN I HAVE MORE ? PLEASE ?!?)
END OF WORK
(RUN TIME = 50 msec)
```
This is a fictitious example, not taken from any actual application, and was intended to demonstrate the power of our differentiating and printing procedures.
## 3 Output in the form of Latex input and in the form of Ortocartan’s input
The simple call to Ortocartan shown above would result in printing all the formulae in the standard mathematical format. One item of the output is shown below (the tetrad $`R_{00}`$ component of the Ricci tensor):
```
-1
> ricci = 2 r exp (- 2 mu) nu, + exp (- 2 mu) nu,
0 0 r
```
```
2
> + exp (- 2 mu) nu, - exp (- 2 mu) mu, nu,
r r r r
```
```
2
> - exp (- 2 nu) mu, - exp (- 2 nu) mu, +
r t t t
```
```
> exp (- 2 nu) mu, nu,
t t
```
This can be either just viewed on the screen or stored on a disk and then printed. However, if a formula like the one above is to be inserted in a text for publication, then one may avoid retyping it. The same formula can be obtained in the form of Latex input code if one inserts the following item anywhere between the arguments of the function Ortocartan (but after the title (SPHERICAL METRIC…)):
```
(output for latex)
```
The same component of the Ricci tensor will then appear on output in the following form
```
\begin{equation}
ricci_{00} = 2r^{-1}\exp (- 2{\mu}){\nu},_{r} + \exp (- 2{\mu}){{\nu},_{r}}^{2}
$$
$$ + \exp (- 2{\mu}){\nu},_{r r} - \exp (- 2{\mu}){\mu},_{r}{\nu},_{r} - \exp
(- 2{\nu}){{\mu},_{t}}^{2} - \exp (- 2{\nu}) $$
$$ {\mu},_{t t} + \exp (- 2{\nu}){\mu},_{t}{\nu},_{t}
\end{equation}
```
Note that Ortocartan has recognized the Greek letters and replaced them with the appropriate Latex constructs. Then the user can add his/her favourite Latex preamble to this and run Latex on the result obtaining the following output:
$$ricci_{00}=2r^1\mathrm{exp}(2\mu )\nu ,_r+\mathrm{exp}(2\mu )\nu ,_{r}^{}{}_{}{}^{2}$$
$$+\mathrm{exp}(2\mu )\nu ,_{rr}\mathrm{exp}(2\mu )\mu ,_r\nu ,_r\mathrm{exp}(2\nu )\mu ,_{t}^{}{}_{}{}^{2}\mathrm{exp}(2\nu )$$
$$\mu ,_{tt}+\mathrm{exp}(2\nu )\mu ,_t\nu ,_t$$
(1)
The `exp(f)` can be very easily replaced by $`\mathrm{e}^f`$ using the substitutions. The exponential function is represented as above in order to give the user a greater freedom: the symbol ”e” can be used for anything the user wishes and is not reserved to mean only the base of natural logarithms.
If the output or some part of it is to be used as input for one of the programs of the Ortocartan set, then the user can avoid rewriting again. It is then enough to insert the following instead of the ”(output for latex)” line:
```
(output for input)
```
The whole output will then be written in the Ortocartan input notation and any part of it can be inserted in the actual input just by cutting and pasting. The same component of the Ricci tensor that was shown above would then appear in the following form:
```
ricci (0 0) = ((2 * (r ^ -1) * (exp (-2 * mu)) * (der r nu)) + ((exp (-2 * mu
)) * ((der r nu) ^ 2)) + ((exp (-2 * mu)) * (der r r nu)) + (-1 * (exp (-2 *
mu)) * (der r mu) * (der r nu)) + (-1 * (exp (-2 * nu)) * ((der t mu) ^ 2)) +
(-1 * (exp (-2 * nu)) * (der t t mu)) + ((exp (-2 * nu)) * (der t mu) * (der t
nu)))
```
The automatically generated output in the input format has the tendency to use more parentheses than are necessary. The algebraically equivalent input written by the user would not contain any of the unnecessary parentheses because Ortocartan’s input notation recognizes the usual hierarchy of priorities among the algebraic operations – see manual . The additional parentheses do not change this hierarchy.
The ”(output for latex)” and ”(output for input)” options are available for all the other programs described below.
## 4 The program Ellisevol.
This program calculates all the quantities appearing in the evolution equations of the kinematical tensors of fluid flow, as defined by Ellis . Since all these equations are consequences of the Ricci identity $`u_{\alpha ;\beta \gamma }u_{\alpha ;\gamma \beta }=u_\mu R_{}^{\mu }{}_{\alpha \beta \gamma }{}^{}`$, they will be fulfilled identically in most cases. However, they may fail to be identically fulfilled when one makes assumptions about separate parts of the flow, e.g. if one assumes that the shear is zero. As is well known, such assumptions have consequences for the other characteristics of the flow, and the Ellis equations will show what the consequences are. Along the way, the program calculates all those quantities that are calculated by the program Ortocartan, and in addition the expansion, the acceleration, the shear tensor and scalar, the rotation tensor and scalar, and the electric and magnetic parts of the Weyl tensor with respect to the velocity vector.
Since the signature assumed in the calculation is $`(+)`$, the formulae may differ from those given in textbooks, and so some of them are quoted below for reference. The equations that the program verifies are the following.
The rotation constraint equations:
$$\omega _{[\alpha \beta ;\gamma ]}+\dot{u}_{[\alpha ;\gamma }u_{\beta ]}+\dot{u}_{[\alpha }\omega _{\beta \gamma ]}=0$$
(square brackets on indices denote antisymmetrization, round brackets on indices denote symmetrization).
The shear constraint equations:
$$h_{}^{\alpha }{}_{\beta }{}^{}(\omega _{}^{\beta \gamma }{}_{;\gamma }{}^{}\sigma _{}^{\beta \gamma }{}_{;\gamma }{}^{}+\frac{2}{3}\theta ^{;\beta })(\omega _{}^{\alpha }{}_{\beta }{}^{}+\sigma _{}^{\alpha }{}_{\beta }{}^{})\dot{u}^\beta =0,$$
where $`h_{}^{\alpha }{}_{\beta }{}^{}=\delta _{}^{\alpha }{}_{\beta }{}^{}u^\alpha u_\beta `$ is the projection tensor.
The rotation evolution equations:
$$h_{\alpha }^{}{}_{}{}^{\gamma }h_{\beta }^{}{}_{}{}^{\delta }\dot{\omega }_{\gamma \delta }h_{\alpha }^{}{}_{}{}^{\gamma }h_{\beta }^{}{}_{}{}^{\delta }\dot{u}_{[\gamma ;\delta ]}+2\sigma _{\delta [\alpha }\omega _{}^{\delta }{}_{\beta ]}{}^{}+\frac{2}{3}\theta \omega _{\alpha \beta }=0.$$
The Raychaudhuri equation:
$$\dot{\theta }+\frac{1}{3}\theta ^2\dot{u}_{}^{\alpha }{}_{;\alpha }{}^{}+\sigma ^{\alpha \beta }\sigma _{\alpha \beta }\omega ^{\alpha \beta }\omega _{\alpha \beta }+R_{\alpha \beta }u^\alpha u^\beta =0.$$
The (coordinate) electric components of the Weyl tensor:
$$E_{\alpha \beta }=C_{\alpha \rho \beta \sigma }u^\rho u^\sigma E_{\beta \alpha }.$$
The shear evolution equations:
$$h_{\alpha }^{}{}_{}{}^{\gamma }h_{\beta }^{}{}_{}{}^{\delta }\dot{\sigma }_{\gamma \delta }h_{\alpha }^{}{}_{}{}^{\gamma }h_{\beta }^{}{}_{}{}^{\delta }\dot{u}_{(\gamma ;\delta )}+\dot{u}_\alpha \dot{u}_\beta +\omega _{\alpha \gamma }\omega _{}^{\gamma }{}_{\beta }{}^{}+\sigma _{\alpha \gamma }\sigma _{}^{\gamma }{}_{\beta }{}^{}+\frac{2}{3}\theta \sigma _{\alpha \beta }$$
$$+\frac{1}{3}h_{\alpha \beta }[2(\omega ^2\sigma ^2)+\dot{u}_{}^{\gamma }{}_{;\gamma }{}^{}]+E_{\alpha \beta }=0.$$
The magnetic components of the Weyl tensor:
$$H_{\alpha \beta }=\frac{1}{2}\sqrt{g}\epsilon _{\alpha \gamma \mu \nu }C_{}^{\mu \nu }{}_{\beta \delta }{}^{}u^\gamma u^\delta H_{\beta \alpha },$$
where $`\epsilon _{\alpha \gamma \mu \nu }`$ is the Levi-Civita symbol.
The ”magnetic constraint” equations:
$$2\dot{u}_{(\alpha }w_{\beta )}\sqrt{g}h_{\alpha }^{}{}_{}{}^{\gamma }h_{\beta }^{}{}_{}{}^{\delta }(\omega _{(\gamma }^{}{}_{}{}^{\mu ;\nu }+\sigma _{(\gamma }^{}{}_{}{}^{\mu ;\nu })\epsilon _{\delta )\rho \mu \nu }u^\rho =H_{\alpha \beta },$$
where $`w_\beta `$ is the rotation vector field defined by:
$$w^\alpha =\frac{1}{2\sqrt{g}}\epsilon ^{\alpha \beta \gamma \delta }u_\beta \omega _{\gamma \delta }.$$
As an example, let us consider the application of this program to the metric of Lanczos .
$$ds^2=(dt+Crd\phi )^2\psi d\phi ^2\frac{1}{4}\mathrm{e}^rdr^2/\psi \mathrm{e}^rdz^2,$$
where $`C=`$ const and
$$\psi =(C^2r+\mathrm{\Lambda }\mathrm{\Lambda }\mathrm{e}^r).$$
This is a stationary cylindrically symmetric solution of Einstein’s equations with a rotating dust source and with a nonvanishing cosmological constant $`\mathrm{\Lambda }`$. The coordinates used in the metric shown above are comoving and the velocity vector field of the dust is one of the orthonormal tetrad vectors, hence the tetrad components of the velocity field are `(1 0 0 0)`. Since this is a solution of Einstein’s equations, this vector field is uniquely determined by the metric, and so, as expected, all the constraint and evolution equations will be identities. However, the acceleration (= 0), rotation, expansion (= 0), and shear (= 0) are all calculated, along with the electric and magnetic parts of the Weyl tensor.
The input data is here:
```
(setq !*lower nil)
(ellisevol’(
(LANCZOS METRIC)
(coordinates t phi r z)
(velocity 1 0 0 0)
(constants C Lambda)
(symbols psi = (C ^ 2 * r + Lambda - Lambda * (exp (- r))) )
(ematrix 1 (C * r) 0 0
0 ((C ^ 2 * r + Lambda - Lambda * (exp (- r)))^ (1 2 ))
0 0 0 0 ((1 2) * (exp ((-1 2) * r)) * (C ^ 2 * r + Lambda
- Lambda * (exp (- r))) ^ (-1 2)) 0
0 0 0 (exp (- (1 2) * r)))
(substitutions (C ^ 2 * r + Lambda - Lambda * (exp (- r))) = psi )
(dont print messages)
(tensors einstein)
))
(setq !*lower t)
```
The first line of the input tells the system to treat the lower-case letters and their corresponding capitals as different symbols, the last line reverses this command. The command ”(dont print messages)” tells the system not to print the messages about unsuccessful attempts to carry out the substitution specified in the previous line. Without this command, the program would write a message every time when there was no opportunity for this substitution in a given formula.
The most important parts of the output are given below. Each tensor is printed with its unique name to facilitate substitutions. Thus the covariant rotation tensor $`\omega _{\alpha \beta }`$ has the name ”rotdd”, while the mixed rotation tensor $`\omega _{\alpha }^{}{}_{}{}^{\beta }`$ has the name ”rotdu”, and similarly for shear.
```
ACCELERATION = 0
```
```
> rotdd = (1/2) C
1 2
```
```
(rotdd calculated)
(TIME = 1060 msec)
(rotdd completed)
(TIME = 1070 msec)
```
```
2
> rotdu = - 2 C exp (r) psi
1
```
```
0 2 -1
> rotdu = - (1/2) C r psi
2
```
```
1 -1
> rotdu = (1/2) C psi
2
```
```
(rotdu completed)
(TIME = 1100 msec)
```
```
2
> ROTATION SQUARED = C exp (r)
```
```
(ROTATION SCALAR calculated)
(TIME = 1100 msec)
> EXPANSION SCALAR = 0
(EXPANSION SCALAR calculated)
(TIME = 1150 msec)
(sheardd calculated)
(TIME = 1160 msec)
SHEAR = 0
(ALL THE ROTATION CONSTRAINTS ARE FULFILLED IDENTICALLY)
(ROTATION CONSTRAINTS calculated)
(TIME = 1160 msec)
(ALL THE SHEAR CONSTRAINTS ARE FULFILLED IDENTICALLY)
(SHEAR CONSTRAINTS calculated)
(TIME = 2500 msec)
(ALL THE ROTATION EVOLUTION EQUATIONS ARE FULFILLED IDENTICALLY)
(ROTATION EVOLUTION EQUATIONS calculated)
(TIME = 2800 msec)
> RAYCHAUDHURI EQUATION = 0
(RAYCHAUDHURI EQUATION calculated)
(TIME = 3150 msec)
```
```
2
> elweyl = - (1/3) C exp (r) psi
1 1
```
```
2 -1
> elweyl = - (1/12) C psi
2 2
```
```
2
> elweyl = (2/3) C
3 3
```
```
(elweyl calculated)
(TIME = 3310 msec)
(ALL THE SHEAR EVOLUTION EQUATIONS ARE FULFILLED IDENTICALLY)
(SHEAR EVOLUTION EQUATIONS calculated)
(TIME = 3820 msec)
```
```
> magweyl = - (1/2) C
2 3
```
```
(magweyl calculated)
(TIME = 3890 msec)
(ALL THE MAGNETIC CONSTRAINTS ARE FULFILLED IDENTICALLY)
(magcons calculated)
(TIME = 5170 msec)
(I REALLY LIKED THIS! CAN I HAVE MORE ? PLEASE ?!?)
END OF WORK
(RUN TIME = 5170 msec)
```
## 5 The program Curvature.
This program calculates the curvature tensor from given connection coefficients in any number of dimensions. The connection coefficients are assumed symmetric (i.e. torsion-free), but need not be metrical. The program was written for one special application, and hence the assumption of zero torsion; it can be removed without any difficulty.
## 6 The program Landlagr.
This program calculates the reduced lagrangian for the Einstein equations by the Landau–Lifshitz method . This is the Hilbert lagrangian with the divergences containing second derivatives of the metric already removed. In short, this (noncovariant) lagrangian is obtained by deleting from the scalar curvature those terms in which the Christoffel symbols are differentiated, and taking the remaining part with the reverse sign.
As is well-known, the Euler-Lagrange equations derived from a variational principle in which the generality of the metric was limited by various assumptions (e.g. about symmetry) are not always equivalent to the Einstein equations. It may happen that they will have nothing to do with the Einstein equations; this is known, for example, to occur for certain Bianchi-type models . Therefore, the user must make sure, when using the program ”landlagr”, that the situation is appropriate for applying the lagrangian methods.
## 7 The program Eulagr.
This program calculates the Euler-Lagrange equations starting from a lagrangian specified by the user. It is assumed that the resulting E-L equations will be ordinary differential equations (i.e. that there is only one independent variable in the lagrangian action integral).
As an example, the program will derive the Newtonian equations of motion for a point particle of mass $`m`$ in the cartesian coordinates $`\{x,y,z\}`$ from the lagrangian
$$L=\frac{1}{2}m(\dot{x}^2+\dot{y}^2+\dot{z}^2)V(x,y,z),$$
where $`V`$ is a potential and $`x(t),y(t),z(t)`$ are the equations of a trajectory of the particle. The input data is:
```
(setq !*lower nil)
(eulagr ’(
(The lagrangian for the Newtonian equations of motion
in 3 dimensions)
(constants m)
(parameter t)
(functions x(t) y(t) z(t) V(x y z) )
(variables x y z)
(lagrangian ((1 2) * m * ((der t x) ^ 2 + (der t y) ^ 2
+ (der t z) ^ 2) - V))
))
(setq !*lower t)
```
and the results are:
```
(The lagrangian for the Newtonian equations of motion in 3 dimensions)
```
```
2 2 2
> lagrangian = - V + (1/2) m x, + (1/2) m y, + (1/2) m z,
t t t
```
```
(THIS IS THE VARIATIONAL DERIVATIVE BY x)
> eulagr = m x, + V,
0 t t x
```
```
(THIS IS THE VARIATIONAL DERIVATIVE BY y)
> eulagr = m y, + V,
1 t t y
```
```
(THIS IS THE VARIATIONAL DERIVATIVE BY z)
> eulagr = m z, + V,
2 t t z
```
```
(I REALLY LIKED THIS! CAN I HAVE MORE ? PLEASE ?!?)
END OF WORK
(RUN TIME = 100 msec)
```
All the other programs described in this paper are, from the programmer’s point of view, rather similar to each other. The various programs are simply different sets of the same basic algebraic operations. The program Eulagr is somewhat exceptional, since, when calculating variational derivatives, the functional expressions $`\dot{x}^i,i=1,2,\mathrm{},n`$ have to be treated as independent variables at a certain stage. The program generates names for these variables, adds them to the list of independent variables for differentiation, replaces the $`\dot{x}^i`$ by the new variables, then calculates the derivatives $`L/\dot{x}^i`$ and $`L/x^i`$, then replaces the new variables back by $`\dot{x}^i`$ in the resulting expressions, and finally uses these results to calculate $`\frac{\mathrm{d}}{\mathrm{d}t}(L/\dot{x}^i)`$. All this happens automatically, and the user need not worry about any details of it. Ortocartan’s differentiating and substituting subprograms are flexible enough to handle this.
## 8 The program Squint.
This program verifies whether a given expression is a first integral of a given set of ordinary differential equations. The program was written for a specific application, therefore it is rather limited in its abilities. It is assumed that the (hypothetical or actual) first integral is a polynomial of first or second degree in the first derivatives of the functions that should obey the set of equations. It is also assumed that the equations in the set are all of second or first order and that they have the form:
$$f^i,_{tt}=\mathrm{\Psi }^i(f^1,\mathrm{},f^n,f^1,_t,\mathrm{},f^n,_t)$$
i.e. that they are algebraically solved with respect to the highest derivatives. ”Squint” is an abbreviation for ”square integral”.
In order to make the result easy to verify, the program ”squint” will be used to find a first integral of the equations found in the previous example. We shall at first pretend that we do not know what the integral $`I`$ should be and will assume that it is a general polynomial of second degree in the first derivatives by $`t`$ of the functions $`x(t)`$, $`y(t)`$ and $`z(t)`$, thus $`I=Q_{ij}\frac{\mathrm{d}x^i}{\mathrm{d}t}\frac{\mathrm{d}x^j}{\mathrm{d}t}+L_i\frac{\mathrm{d}x^i}{\mathrm{d}t}+E`$, where $`Q_{ij}`$, $`L_i`$ and $`E`$ are unknown functions of the $`x^i`$. The argument ”(markers M)” defines $`M`$ as a symbolic variable that can be used to represent any expression. Thanks to this facility the single equation (der t t M) = (- (der M V) / m) represents the 3 equations $`d^2x^i/dt^2=V/x^i`$ for $`i=1,2,3`$ simultaneously. The ”maineq” directs these substitutions to the main equation, which is the explicit formula for the total derivative $`\mathrm{d}I/\mathrm{d}t`$.
The input data is:
```
(setq !*lower nil)
(squint’(
(a first integral of the Newtonian equations of motion)
(constants m)
(parameter t)
(functions x(t) y(t) z(t) V(x y z) Q11(x y z) Q12(x y z) Q13(x y z)
Q22(x y z) Q23(x y z) Q33(x y z) L1(x y z) L2(x y z) L3(x y z)
E(x y z) )
(variables x y z)
(integral (Q11 * (der t x) ^ 2 + 2 * Q12 * (der t x) * (der t y)
+ 2 * Q13 * (der t x) * (der t z) + Q22 * (der t y) ^ 2
+ 2 * Q23 * (der t y) * (der t z) + Q33 * (der t z) ^ 2
+ L1 * (der t x) + L2 * (der t y) + L3 * (der t z) + E) )
(markers M)
(substitutions
maineq
(der t t M) = (- (der M V) / m)
)
(dont print maineq)
))
(setq !*lower t)
```
The output includes the integral $`I`$ printed in the mathematical format, the ”main equation” $`\mathrm{d}I/\mathrm{d}t`$ (the printing of $`\mathrm{d}I/\mathrm{d}t`$ was suppressed by the command ”(dont print maineq)”), the coefficients of all the expressions $`\frac{\mathrm{d}x^i}{\mathrm{d}t}\frac{\mathrm{d}x^j}{\mathrm{d}t}\frac{\mathrm{d}x^k}{\mathrm{d}t}`$, $`\frac{\mathrm{d}x^i}{\mathrm{d}t}\frac{\mathrm{d}x^j}{\mathrm{d}t}`$, $`\frac{\mathrm{d}x^i}{\mathrm{d}t}`$, and the remaining part of $`\mathrm{d}I/\mathrm{d}t`$ that does not contain any $`\frac{\mathrm{d}x^i}{\mathrm{d}t}`$. This will be done correctly only if the second derivatives $`\frac{\mathrm{d}^2x^i}{\mathrm{d}t^2}`$ appearing in $`\mathrm{d}I/\mathrm{d}t`$ had been replaced using the (substitutions …), in this way use is made of the set of equations for which $`I`$ is expected to be the first integral. If the set consists of $`n`$ equations, then the number of coefficients printed will be $`\frac{1}{6}(n+1)(n+2)(n+3)`$, i.e. 20 for $`n=3`$. When the explicit form of $`I`$ is unknown, as in this example, the coefficients will be partial differential equations that should determine the $`Q_{ij}`$, $`L_i`$ and $`E`$ (each of the ”equations” printed is a left-hand side of an equation that has zero on the right-hand side).
Only a few pieces of the output are reproduced below.
```
(a first integral of the Newtonian equations of motion)
```
```
2
> integral = E + Q11 x, + 2 Q12 x, y, + 2 Q13 x, z, + Q22
t t t t t
```
```
2 2
> y, + 2 Q23 y, z, + Q33 z, + L1 x, + L2 y, + L3 z,
t t t t t t
```
```
>
t
```
```
3
> THIS IS THE COEFFICIENT OF x,
t
```
```
> equation = Q11,
1 x
```
```
2
> THIS IS THE COEFFICIENT OF x, y,
t t
```
```
> equation = 2 Q12, + Q11,
2 x y
```
```
..............................
[part of the output deleted]
..............................
> THESE ARE THE TERMS THAT ARE FREE OF THE DERIVATIVES
```
```
-1 -1 -1
> equation = - m L1 V, - m L2 V, - m L3 V,
20 x y z
```
```
(I REALLY LIKED THIS! CAN I HAVE MORE ? PLEASE ?!?)
END OF WORK
(RUN TIME = 1030 msec)
```
Now let us substitute the well-known solution of these equations into the data and see what happens.
```
(setq !*lower nil)
(squint’(
(a first integral of the Newtonian equations of motion - the final result)
(constants m)
(parameter t)
(functions x(t) y(t) z(t) V(x y z))
(variables x y z)
(integral ((1 2) * m * ((der t x) ^ 2 + (der t y) ^ 2
+ (der t z) ^ 2) + V) )
(markers M)
(substitutions
maineq
(der t t M) = (- (der M V) / m)
)
(dont print maineq)
))
(setq !*lower t)
```
The result is:
```
(a first integral of the Newtonian equations of motion - the final result)
```
```
2 2 2
> integral = V + (1/2) m x, + (1/2) m y, + (1/2) m z,
t t t
```
```
(THE FIRST INTEGRAL IS ALREADY MAXIMALLY SIMPLIFIED
AND IS EXPLICITLY CONSTANT)
> maineq = 0
(I REALLY LIKED THIS! CAN I HAVE MORE ? PLEASE ?!?)
END OF WORK
(RUN TIME = 130 msec)
```
Similar programs for verifying first integrals can be written for any kind of dependence of the first integral on the first derivatives of the functions.
## Appendix A How to obtain Ortocartan.
The base language of this newest release of Ortocartan is Codemist Standard Lisp. The latter is implemented in Windows 98 and in Linux, and so should be usable for most computer users.
In order to use Ortocartan one must first buy the Codemist Standard Lisp. It can be bought from:
Professor John Fitch, Director
CODEMIST Limited
”Alta”, Horsecombe Vale
Combe Down
BATH, Avon, BA2 5QR
England
phone and fax (44-1225) 837 430
email jpffitch@maths.bath.ac.uk
Ortocartan is free of charge. If you wish to have it, just contact A. Krasiński, I will either email the programs to you or send you a diskette. How to install Ortocartan when Lisp is already working is described in sec. 7 of the manual .
REFERENCES
A. Krasiński, Gen. Rel. Grav. 25, 165 (1993).
A. Krasiński, M. Perkowski, The system Ortocartan - user’s manual. Fifth edition, Warszawa 2000. A Latex file distributed by email or on diskettes.
G. F. R. Ellis, in General relativity and cosmology. Proceedings of the International School of Physics ”Enrico Fermi”, Course 47: General Relativity and Cosmology. Edited by R. K. Sachs, Academic Press, New York 1971.
K. Lanczos, Zeitschrift für Physik 21, 73 (1924); Gen. Rel. Grav. 29, 363 (1997).
L. D. Landau and E. M. Lifshitz, Teoria polya, 6th Russian edition. Izdatelstvo ”Nauka”, Moskva 1973, sec, 93.
M. A. H. MacCallum, in General relativity. An Einstein centenary survey. Edited by S. W. Hawking and W. Israel. Cambridge University Press 1979, p. 552.
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# REDSHIFT DISTORTIONS AND CLUSTERING IN THE PSCZ SURVEY
## 1 Introduction
It has been understood for a long time that redshift surveys are systematically distorted by peculiar velocities $`^\mathrm{?}`$ – so called redshift distortions, which can cause the observed distribution of galaxies to become anisotropic. On large scales where linear theory applies, redshift distortions are characterised by the parameter <sup>a</sup><sup>a</sup>a$`\mathrm{\Omega }_\mathrm{m}`$ is the contribution to the density parameter from matter. The bias parameter $`b`$ depends on galaxy type, in this paper it refers to IRAS galaxies. $`\beta \mathrm{\Omega }_\mathrm{m}^{0.6}/b`$ – see Hamilton $`^\mathrm{?}`$ for a detailed review.
In a classic paper, Kaiser $`^\mathrm{?}`$ derived a simple formula for the effect of linear redshift distortions for a volume limited survey that subtends a small opening angle at the observer (the “distant observer” approximation – all lines of sight are treated as parallel). He showed that power is boosted by a factor $`(1+\beta \mu _k^2)^2`$, where $`\mu _k`$ is the cosine of the angle between wavevector and line of sight. Hamilton $`^\mathrm{?}`$ $`^\mathrm{?}`$ and Cole, Fisher & Weinberg $`^\mathrm{?}`$ $`^\mathrm{?}`$ analysed all-sky IRAS surveys using a method based on the Kaiser formalism, but were forced to break the survey up into sections because of the constraints of the distant observer approximation, losing information about the largest (and most reliably linear) scales.
Rather than fit a square peg into a round hole, Fisher et al $`^\mathrm{?}`$ and Heavens & Taylor $`^\mathrm{?}`$ (HT) dropped the plane parallel approximation and used a spherical harmonic decomposition to match the spherical nature of the IRAS redshift surveys. HT used spherical Bessel functions to decompose the density field radially; using eigenfunctions of the Laplacian retains all the advantages of Fourier analysis. HT analysed the IRAS 1.2Jy survey, fitting $`\beta `$ and the amplitude of the power spectrum; Ballinger, Heavens & Taylor $`^\mathrm{?}`$(BHT) extended this analysis to fit the shape of the power spectrum. Tadros et al $`^\mathrm{?}`$ (T99) applied these techniques to the IRAS PSCz survey – section 4 includes a review of those results.
## 2 Spherical Harmonic Formalism
We will briefly review the formalism for spherical harmonic analysis – see HT and T99 for details. The density field of the galaxy distribution $`\rho (𝐬)`$ is expanded in terms of spherical harmonics, $`Y_\mathrm{}m`$, and a discrete set of spherical Bessel functions, $`j_{\mathrm{}}`$,
$$\widehat{\rho }_{\mathrm{}mn}=c_\mathrm{}nd^3s\rho (𝐬)w(s)j_{\mathrm{}}\left(k_\mathrm{}ns\right)Y_\mathrm{}m^{}(\theta ,\varphi ),$$
(1)
where the $`c_\mathrm{}n`$ are normalization constants and $`k_{ln}`$ are discrete wavenumbers.
These observed coefficients can be related to those of the true underlying density field ($`\delta _{\mathrm{}mn}`$) by:
$$\widehat{\rho }_{\mathrm{}mn}=\left(\rho _0\right)_{\mathrm{}mn}+\underset{\mathrm{}^{}m^{}n^{}}{}\underset{n^{\prime \prime }}{}S^{nn^{\prime \prime }}W_{\mathrm{}\mathrm{}^{}}^{mm^{}}\left(\mathrm{\Phi }_{\mathrm{}\mathrm{}^{}}^{n^{\prime \prime }n^{}}+\beta V_{\mathrm{}\mathrm{}^{}}^{n^{\prime \prime }n^{}}\right)\delta _{\mathrm{}^{}m^{}n^{}}.$$
(2)
The transition matrices $`𝐖`$, $`𝚽`$, $`𝐕`$ and $`𝐒`$ describe the effects of the sky mask, the radial selection function, the linear redshift space distortion and the small scale distortion correction respectively; $`\left(\rho _0\right)_{\mathrm{}mn}`$ is a mean term, non-zero because of partial sky coverage. The transition matrices are derived and defined in HT and T99.
A likelihood approach is used to constrain $`\beta `$ and the real-space power spectrum:
$$2\mathrm{ln}[𝐃|\beta ,P(k)]=\mathrm{ln}(det𝐂)+\mathrm{𝐃𝐂}^1𝐃^T,$$
(3)
where $`𝐂=<\mathrm{𝐃𝐃}^T>`$ and elements of the data vector $`𝐃`$ are given by $`D_{\mathrm{}mn}[\widehat{\rho }_{\mathrm{}mn}(\rho _0)_{\mathrm{}mn}]/\overline{\rho }`$ where $`\overline{\rho }`$ is the mean number density – Gaussian statistics are assumed. Two different parameterisations of $`P(k)`$ are used: following HT a fixed shape is assumed and the amplitude is fitted, following BHT a stepwise maximum likelihood method is used, allowing the power spectrum to assume any shape (within bin discreteness).
## 3 Data – The IRAS PSCz Survey
The PSCz survey $`^\mathrm{?}`$ is a redshift catalogue complete down to 0.6Jy over $`83\%`$ of the sky with a total of $`15,000`$ redshifts – it is the largest all sky survey in existence. As our method is very precise, we minimise systematic errors by using a more conservative flux cut (0.75Jy) and sky mask reducing the number of redshifts by a factor of roughly two – see T99.
## 4 Previous Results
Both the fixed amplitude (HT) and stepwise $`P(k)`$ (BHT) methods were applied to the PSCz by T99 for modes with $`k0.13h`$ Mpc<sup>-1</sup>. The first method produced $`\beta =0.58\pm 0.26`$ and the amplitude of the real space power measured at wavenumber $`k=0.1h`$ Mpc<sup>-1</sup> of $`\mathrm{\Delta }_{0.1}=0.42\pm 0.03`$ – see Fig. 2 (dashed contours). Freeing the shape of the power spectrum we find the consistent results $`\beta =0.47\pm 0.16`$ (conditional error), and $`\mathrm{\Delta }_{0.1}=0.47\pm 0.03`$ – Fig. 1. T99 also carried out extensive tests on simulations, and the methods were found to be reliable. In addition, we carried out a suite of tests for systematics effects in the data and found the cut catalogue gave consistent results.
The method was restricted to a small range of wavenumber because the likelihood analysis involves the repeated inversion of a large $`n\times n`$ matrix; the time required for this process grows as $`n^3`$. More importantly, the matrix rapidly became numerically unstable. The results above do not strongly rule out either high ($`1`$) or low ($`0.5`$) values of $`\beta `$. It would be nice to overcome the matrix problem, extend the k-range and reduce the error bar.
## 5 Data Compression
It is possible to transform the (length $`n`$) data vector $`𝐃`$ to create a new, smaller dataset, while retaining most of the information about the parameters of interest $`^\mathrm{?}`$. A new dataset is constructed which is a linear combination of the original:
$$𝐃^{}=\mathrm{𝐁𝐃},$$
(4)
where $`𝐃^{}`$ is a new data vector of length $`n^{}<n`$, with a corresponding transformation for the covariance matrix. TTH show that the optimal matrix $`𝐁`$ which minimises the conditional error on a single parameter is made up from eigenvectors of <sup>b</sup><sup>b</sup>bThis, like the Wiener filter, is one of those marvellous data-handling methods which were invented by astronomers then sent back in time so that signal processing engineers could use them for the past fifty years.
$$𝐂,_i𝐛=\lambda \mathrm{𝐂𝐛}$$
(5)
where the comma refers to a derivative with respect to parameter $`i`$. The new covariance matrix is smaller and also close to diagonal – hence much more stable.
TTH suggested extending this to multi-parameter problems by constructing a separate matrix $`𝐁`$ for each parameter and then using singular value decomposition to combine the matrices efficiently. However, Ballinger $`^\mathrm{?}`$ (see also Taylor et al $`^\mathrm{?}`$) tested this and found that error ellipses/ellipsoids tended to grow along correlation axes – only conditional errors are constrained. Instead it was shown that equation (5) could be used to constrain the marginal errors by using one or more linear combinations of the original parameters which lie directly along the correlation axes.
Unfortunately solving equation (5) involves manipulating uncompressed $`n\times n`$ matrices. This usually isn’t a timescale problem – it need only be done once – but still suffers from numerical difficulties. To avoid this problem, we split the data into several sections and compressed them separately, while still retaining full information about correlations between modes in different sections. Strictly speaking this is slightly less optimal than compressing the whole dataset in one go, but it should make a negligible difference in practice and will not bias the result.
## 6 New Results and conclusions
Fig. 2 shows the old $`k_{\mathrm{max}}=0.13h`$ Mpc<sup>-1</sup> result together with the new $`k_{\mathrm{max}}=0.2h`$ Mpc<sup>-1</sup> obtained using data compression. The original 4644 modes were compressed down to 2278. The new, lower, value of $`\beta =0.4\pm 0.1`$ is consistent with the previous result but the error ellipse is considerably smaller; the best-fit amplitude of the power spectrum is essentially unchanged. $`\beta =1`$, corresponding to $`\mathrm{\Omega }_\mathrm{m}=1`$, $`b=1`$, is now strongly disfavoured. More details will be in Ballinger et al and Taylor et al (in preparation).
The low value of $`\beta `$ is consistent with currently popular cosmological models with a low value of $`\mathrm{\Omega }_{\mathrm{matter}}`$ if the IRAS bias parameter is close to unity. The value is somewhat lower than that from other analyses of the PSCz survey $`^\mathrm{?}`$ $`^\mathrm{?}`$ $`^\mathrm{?}`$ $`^\mathrm{?}`$, but not inconsistent. It is consistent with the recent velocity-velocity comparison results from peculiar velocity catalogues $`^\mathrm{?}`$, but somewhat lower than the corresponding density-density value $`^\mathrm{?}`$. See Willick $`^\mathrm{?}`$ for a discussion.
## References
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# MS-TPI-00-3hep-th/0005084 Distribution of Instanton Sizes in a Simplified Instanton Gas Model
## 1 Introduction
### 1.1 Instantons in gauge theories
In non-abelian gauge theories different topologically nontrivial configurations have been made responsible for non-perturbative features . An important class are instantons, which are solutions of the Euclidean field equations with non-vanishing topological charge . They give contributions to the saddle-point approximation of Euclidean functional integrals, which lead to non-perturbative effects . For a review see .
In the dilute gas approximation one considers superpositions of single instantons as quasi saddle points of the action. These configurations are characterized by the instanton positions $`\{𝐚_𝐣\}`$ in four-dimensional space-time, the instanton sizes $`\{\rho _j\}`$, and other internal parameters. For a single instanton the action is
$$S_1=\frac{8\pi ^2}{g_0^2},$$
(1)
where $`g_0`$ is the bare gauge coupling constant. Taking into account quadratic fluctuations around the one-instanton solution its contribution to the functional integral is
$$Z_1=d^4𝐚𝑑\rho C\rho ^5\mathrm{exp}\left\{\frac{8\pi ^2}{g^2(1/\rho )}\right\},$$
(2)
where $`g(\mu )`$ is the running coupling. In the case of supersymmetric Yang-Mills theory this formula has even been established at the two-loop-level, in ordinary Yang-Mills theory there are higher-order corrections to the integrand . In one-loop order the running coupling obeys
$$g^2(\mu )=\frac{8\pi ^2}{b\mathrm{log}(\mu /\mathrm{\Lambda })}$$
(3)
where
$$b=\frac{11N}{3},$$
(4)
and $`\mathrm{\Lambda }`$ is the renormalization-group invariant scale parameter, so that
$$Z_1=d^4𝐚𝑑\rho C\rho ^5(\rho \mathrm{\Lambda })^b.$$
(5)
Therefore the integrand is proportional to $`\rho ^{b5}`$ and increases with increasing instanton size.
The space-time integral over the instanton position gives the usual volume factor, which is needed in the large volume limit in order to get an extensive free energy. The integral over $`\rho `$, however, represents an infrared problem. Where the instanton density becomes important, for $`\rho O(1/\mathrm{\Lambda })`$, we leave the region of valididty of the semiclassical expansion because the running coupling becomes too large.
It should be noted that the apparent infrared divergence in (5) is an artifact of using the one-loop formula for $`g^2(1/\rho )`$ long after it has become invalid, i.e. for $`\rho 1/\mathrm{\Lambda }`$. Nevertheless we are confronted with the infrared problem in the size integration for the single instanton contribution $`Z_1`$. If the semiclassical approximation is meaningful at all, a solution of this problem in the context of the full instanton ensemble is required.
The quasi saddle points composed of any number of instantons and anti-instantons are treated as independent in the dilute gas approximation. Consequently their contributions exponentiate in the usual way. The problem with the integration over the sizes persists and has to be dealt with. The simplest way is to cut the integrations off at some ad-hoc value $`\rho _c`$. But since the integrand increases with increasing $`\rho _j`$ the dominant contribution comes from large $`\rho _j`$ near the cut-off where the assumption of diluteness fails. Moreover the introduction of an ad-hoc cut-off leads to inconsistencies with the renormalization group .
In order to solve the problem of the instanton size integrations is has been proposed that instanton sizes are cut off in a dynamical way . The dynamical cut-off should originate from configurations where instantons start to overlap. Configurations of overlapping instantons have an action which deviates from the sum of the single instanton actions. Therefore large instantons feel an interaction. Additionally, the fluctuations around the multi-instanton configurations contribute to the instanton interaction. The interaction between instantons is expected to suppress overlapping instantons and to result in a dynamical self-consistent cut-off.
Some consequences of this picture have been discussed in , based on certain assumptions about the repulsive instanton interactions. In a temporary infrared cut-off was introduced by means of a finite space-time volume $`V`$. The large $`V`$ limit was then considered with the help of renormalization group arguments. This led to some general results independent of the specific form of the repulsive instanton interactions. In particular, a finite renormalization factor
$$\frac{4}{b}=\frac{12}{11N}$$
(6)
appears in some quantities, e.g. in the trace anomaly , correcting inconsistencies with the renormalization group. The same factor is conjectured to multiply the instanton singularities in the Borel plane, which then coincide with the infrared renormalons. The conclusions were supported by considering a model of the instanton ensemble, where the repulsive interaction is approximated by a hard core.
The theory of instanton ensembles with a dynamical size cut-off has been developed further by Shuryak , see for a review. In his model of an “instanton liquid” various observables have been calculated and related to hadronic phenomenology.
In connection with the dynamical cut-off the distribution of instanton sizes is of central importance. The size-distribution reflects the way in which large instantons are suppressed and thus gives information about the instanton interactions. In recent years it has been studied by means of lattice Monte Carlo calculations by different groups . For small sizes the distribution is predicted to be
$$n(\rho )\rho ^{b5}$$
(7)
by the dilute gas approximation as well as by the “instanton liquid model”, in accordance with Eq. (5). For large sizes $`\rho `$, where the dynamical cut-off is in effect, not much is known about the distribution. There are arguments in favour of a suppression like
$$n(\rho )\mathrm{exp}(c\rho ^p)\text{with}p=2.$$
(8)
In this article we investigate the distribution of instanton sizes in a model where the instanton interactions are approximated by a repulsive hard core of variable size. Although this approximation appears to be crude, the general features of the instanton ensemble with a dynamical cut-off are present. In particular, using analytical and numerical methods we calculate the asymptotic behaviour for small and for large sizes $`\rho `$ and compare them with results from Monte Carlo simulations of lattice gauge theory. More details can be found in .
In it has been conjectured that the distribution $`n(\rho )`$ is affected by the finite renormalization factor $`4/b`$ in such a way that for small $`\rho `$ asymptotically
$$n(\rho )\rho ^1\rho ^{\frac{4}{b}(b4)}.$$
(9)
Using the simplified model, we shall show below that this conjecture is wrong and that instead the semiclassical result (7) holds.
### 1.2 Simplified model for ensembles of instantons
Consider an ensemble of instantons in $`d`$ space-time dimensions. In the spirit of we introduce a finite volume $`V`$ and study the approach to the thermodynamic limit $`V\mathrm{}`$.
In the sector with instanton number $`K`$ the partition function is written as
$$Z_K(V)=\frac{C^K}{K!}\underset{i=1}{\overset{K}{}}d𝐚_id\rho _i\underset{j=1}{\overset{K}{}}\rho _j^{bd1}\mathrm{e}^{U(\{𝐚_k\},\{\rho _k\})},$$
(10)
where the instanton positions and radii are denoted $`\{𝐚_i,\rho _i\}`$. $`C`$ is a constant, whose numerical value is unimportant here, $`b=11N/3`$ for SU($`N`$) Yang-Mills theory, and $`U`$ represents the interaction potential between instantons. Distances are measured in units of $`\mathrm{\Lambda }^1`$.
In our simplified model the repulsive potential is approximated by a hard core potential. The radius of an instanton core varies proportional to the size $`\rho `$ of the instanton. In a finite volume $`V`$ this means
$$\mathrm{e}^{U(\{𝐚_i\},\{\rho _i\})}=\mathrm{\Theta }(\{𝐚_i\},\{\rho _i\})$$
(11)
with
$$\begin{array}{cc}\mathrm{\Theta }(\{𝐚_i\},\{\rho _i\})\hfill & =1,\text{if}\{\begin{array}{c}1.𝐚_i𝐚_j>\left(\frac{\tau }{v_1}\right)^{\frac{1}{d}}(\rho _i+\rho _j)i,j,\text{and}\hfill \\ 2.|a_i^\mu |<\frac{1}{2}V^{\frac{1}{d}}\mu ,i,\text{and}\hfill \\ 3.0<\rho _i<\frac{1}{2}\left(\frac{v_1}{\tau }V\right)^{\frac{1}{d}}i\hfill \end{array}\hfill \\ \mathrm{\Theta }(\{𝐚_i\},\{\rho _i\})\hfill & =0,\text{else}.\hfill \end{array}$$
Here
$$v_1=\frac{\pi ^{\frac{d}{2}}}{\mathrm{\Gamma }\left(\frac{d}{2}+1\right)}$$
(12)
is the volume of the unit sphere in $`d`$ dimensions. The parameter $`\tau `$ specifies the effective volume $`\tau \rho _j^d`$ of an instanton and is of the order of $`v_1`$.
We introduce a reduced distribution by
$$Z_K^{red}(V,\rho _K)=\frac{C^K}{K!}\underset{i=1}{\overset{K}{}}d𝐚_i\underset{j=1}{\overset{K1}{}}d\rho _j\underset{k=1}{\overset{K}{}}\rho _k^{bd1}\mathrm{\Theta }(\{𝐚_l\},\{\rho _l\}),$$
(13)
such that
$$Z_K(V)=𝑑\rho Z_K^{red}(V,\rho ).$$
(14)
For the total system with variable instanton number the grand canonical partition function is
$$Z(V)=\underset{K=0}{\overset{\mathrm{}}{}}Z_K(V),\text{where}Z_0(V)=1.$$
(15)
We do not distinguish between instantons and anti-instantons in this model. In this way we neglect aspects of the interactions which differ between instantons and anti-instantons, but we do not expect that they play a significant role for our considerations.
The probability distribution of instanton numbers is given by
$$𝒫_K(V)=\frac{Z_K(V)}{Z(V)}.$$
(16)
In order to define the probability distribution of instanton sizes one has to specify how the sizes are sampled. The definition should be made in such a way that it is compatible with the Monte Carlo calculations to be discussed later. In the Monte Carlo runs configurations with a variable number of instantons are produced. A configuration of $`K`$ instantons contributes $`K`$ entries to the total histogram of instanton sizes. Therefore it has a relative weight proportional to $`K`$. Correspondingly the probability distribution in the $`K`$-instanton sector is normalized to $`K`$ :
$$\overline{n}_K(V,\rho )=K\frac{Z_K^{red}(V,\rho )}{Z_K(V)}.$$
(17)
In the total ensemble the sizes are then distributed according to
$`\overline{n}(V,\rho )`$ $`=`$ $`{\displaystyle \underset{K=1}{\overset{\mathrm{}}{}}}𝒫_K(V)\overline{n}_K(V,\rho )`$ (18)
$`=`$ $`{\displaystyle \frac{1}{Z(V)}}{\displaystyle \underset{K=1}{\overset{\mathrm{}}{}}}KZ_K^{red}(V,\rho )`$ (19)
with
$$_0^{\mathrm{}}\overline{n}(V,\rho )𝑑\rho =K_V.$$
(20)
As the expectation value of the instanton number grows linearly with the volume $`V`$ one is interested in the rescaled distribution
$$n(V,\rho )=\frac{\overline{n}(V,\rho )}{V}.$$
(21)
In the following sections we study the properties of $`n(V,\rho )`$, and its thermodynamic limit $`n(\rho )`$, respectively, utilizing analytical approaches as well as Monte Carlo-methods.
## 2 The one-dimensional instanton gas
In $`d=1`$ dimensions the model can be solved exactly in the thermodynamic limit. This case illustrates some general features and can serve as a testing ground for approximations used in higher dimensions. Therefore we shall discuss these results before we turn to the consideration of other dimensions.
In the one-dimensional case the canonical partition function for a system of spatial length $`L`$ can be written as
$$Z_K(L)=\frac{C^K}{K!}\underset{i=1}{\overset{K}{}}da_id\rho _i\underset{j=1}{\overset{K}{}}\rho _j^{b2}\mathrm{\Theta }(\{a_k\},\{\rho _k\})$$
(22)
with
$$\begin{array}{cc}\mathrm{\Theta }(\{a_i\},\{\rho _i\})\hfill & =1,\text{if}\{\begin{array}{c}1.)|a_ia_j|>(\frac{\tau }{2})(\rho _i+\rho _j)i,j,\text{and}\hfill \\ 2.)|a_i|<\frac{1}{2}Li,\text{and}\hfill \\ 3.)0<\rho _i<\frac{L}{\tau }i\hfill \end{array}\hfill \\ \mathrm{\Theta }(\{a_i\},\{\rho _i\})\hfill & =0,\text{else}.\hfill \end{array}$$
This represents a system of rods with variable lengths on a line. A pictorial representation is given in Fig. 1.
Integration over the instanton positions $`\{a_i\}`$ yields the effective free volume of $`K`$ indistinguishable particles on the line $`L`$:
$$\underset{i=1}{\overset{K}{}}da_i\mathrm{\Theta }(\{a_k\},\{\rho _k\})=\overline{\mathrm{\Theta }}(\{\rho _k\})\left(L\tau \underset{j=1}{\overset{K}{}}\rho _j\right)^K,$$
(23)
with
$$\overline{\mathrm{\Theta }}(\{\rho _k\})=\{\begin{array}{cc}1,\hfill & \text{if}\tau _{j=1}^K\rho _jL\hfill \\ 0,\hfill & \text{else},\hfill \end{array}$$
(24)
as can be shown by induction. With this result the partition function reads
$$Z_K(L)=\frac{C^K}{K!}_0^{\frac{L}{\tau }}𝑑\rho _K\mathrm{}_0^{\frac{L}{\tau }_{j=2}^K\rho _j}𝑑\rho _1\underset{j=1}{\overset{K}{}}\rho _j^{b2}\left(L\tau \underset{l=1}{\overset{K}{}}\rho _l\right)^K.$$
(25)
The integrations over the $`\rho _j`$ can be carried out successively employing
$$_0^1x^a(1x)^b𝑑x=\frac{\mathrm{\Gamma }(a+1)\mathrm{\Gamma }(b+1)}{\mathrm{\Gamma }(a+b+2)},$$
(26)
and one obtains for the size distribution
$$n(L,\rho )=\frac{_{K=1}^{\mathrm{}}KC\gamma ^{K1}(\mathrm{\Gamma }(b(K1)+2))^1\rho ^{b2}(L\tau \rho )^{b(K1)+1}}{_{K=0}^{\mathrm{}}\gamma ^K(\mathrm{\Gamma }(bK+1))^1L^{bK+1}}$$
(27)
with
$$\gamma =\frac{C\mathrm{\Gamma }(b1)}{\tau ^{b1}}.$$
(28)
The asymptotic behaviour for small instanton sizes is given by the power law
$$n(L,\rho )\text{const.}\rho ^{b2}\text{for}\rho 0$$
(29)
in the sense of
$$b2=\underset{\rho 0}{lim}\frac{\mathrm{ln}(n(L,\rho ))}{\mathrm{ln}(\rho )}.$$
(30)
In order to discuss the behaviour for large $`\rho `$ we take the infinite volume limit
$$n(L,\rho )\stackrel{L\mathrm{}}{}n(\rho ).$$
(31)
The grand canonical sums can be evaluated by replacing them by integrals over $`K`$ and performing a saddle point approximation, which becomes exact in the large-$`L`$ limit. The result is
$$n(\rho )=\frac{C}{b}\rho ^{b2}\mathrm{e}^{c\rho },$$
(32)
with
$$c=(C\mathrm{\Gamma }(b2)\tau )^{\frac{1}{b}}.$$
(33)
In addition to the power law with exponent $`b2=bd1`$ we recognize an exponential suppression of large instanton sizes. The exponent $`p`$ in the exponential, cp. Eq. (8), is equal to 1 in one dimension. Our next aim is to see how these results generalize to higher dimensions $`d`$.
## 3 The general instanton gas
In higher dimensions, $`d>1`$, it is not possible to derive closed expressions for the partition functions or the size distributions. The main difficulty is that the integrations over the instanton positions and the radii cannot be decoupled. In particular we have to use approximations for the effective free volume of sets of instantons in higher dimensions. Nevertheless one can obtain approximate expressions for $`n(V,\rho )`$ and derive its asymptotic behaviour for small $`\rho `$.
The canonical $`K`$-instanton partition function $`Z_K(V)`$ is written as
$$Z_K(V)=\frac{C^K}{K!}_0^{\overline{V}}𝑑v\underset{i=1}{\overset{K}{}}d𝐚_id\rho _i\underset{j=1}{\overset{K}{}}\rho _j^\alpha \mathrm{\Theta }(\{𝐚_k\},\{\rho _k\})\delta (v\tau \underset{l}{}\rho _l^d),$$
(34)
where
$$\alpha =bd1,$$
(35)
and
$$v=\tau \underset{l}{}\rho _l^d\overline{V}$$
(36)
is the total effective volume of instantons. The maximal accessible volume $`\overline{V}V`$ accounts for the fact that spheres in dimensions $`d>1`$ cannot completely occupy a given volume $`V`$.
An approximative decoupling of the integrations can be achieved through the observation that for given $`v=\tau _j\rho _j^d`$ the product
$$\underset{j=1}{\overset{K}{}}\rho _j^\alpha $$
develops a sharp maximum at
$$\rho _j=\left(\frac{v}{K\tau }\right)^{1/d}=\rho _0,j\{1,\mathrm{},K\}$$
(37)
in the thermodynamic limit. The main idea is then to perform a saddle point approximation for the $`\rho `$-integrations near this sharp maximum. The integrations over the positions $`𝐚_j`$ then correspond to a gas of hard spheres with equal radii $`\rho _0`$. With
$$\rho _j=\rho _0+\delta _j$$
(38)
and
$$\delta (v\tau \underset{j=1}{\overset{K}{}}\rho _j^d)\delta (\tau d\rho _0^{d1}\underset{j=1}{\overset{K}{}}\delta _j)=\frac{\rho _0^{1d}}{2\pi \tau d}𝑑q\mathrm{e}^{iq_{j=1}^K\delta _j}$$
(39)
the $`\rho `$-integrations can be solved straightforwardly.
In a similar way the reduced partition function $`Z_K^{red}(V,\rho )`$ can be evaluated. For a given radius $`\rho `$ of the $`K^{th}`$ instanton one assigns an effective volume $`\overline{V}\tau \rho ^d`$ to the remaining $`K1`$ instantons and performs the saddle point approximation for the integrations over $`\rho _1,\mathrm{},\rho _{K1}`$ in terms of Gaussian integrals.
The integral over the instanton positions
$$\frac{1}{K!}\underset{j=1}{\overset{K}{}}d𝐚_j\mathrm{\Theta }(\{𝐚_j\},\{\rho _0\})$$
(40)
can be estimated with the help of geometrical considerations . We use an approximation of the form
$$\frac{1}{K!}\underset{j=1}{\overset{K}{}}d𝐚_j\mathrm{\Theta }(\{𝐚_j\},\{\rho _0\})\frac{1}{K!}(Vv_{eff})^K,$$
(41)
where
$$v_{eff}=h(d)v,$$
(42)
and $`h(d)=V/\overline{V}`$ measures the inverse filling fraction of spheres in a volume $`V`$. In the one-dimensional case we get $`h(1)=1`$ because a given length can be completely filled with rods. For $`d>1`$ we consider $`h(d)`$ as a parameter. Lower and upper bounds are given by $`1h(d)2^{d1}`$. For more details on this point see .
Using these expressions we get by some lengthy but straightforward calculations for the canonical partition functions
$$Z_K(V)\sqrt{\frac{\alpha }{2\pi }}\frac{\sqrt{K}}{d}\left(\sqrt{\frac{2\pi }{\alpha }}\frac{C}{(K\tau h(d))^{\beta 1}}\right)^K\frac{\mathrm{\Gamma }((\beta 1)K)}{\mathrm{\Gamma }(\beta K+1)}V^{\beta K},$$
(43)
and for the reduced ones
$$Z_K^{red}(V,\rho )$$
$$\sqrt{\frac{\alpha }{2\pi }}\frac{\sqrt{K1}C\rho ^\alpha }{d}\left(\sqrt{\frac{2\pi }{\alpha }}\frac{C}{((K1)\tau h(d))^{\beta 1}}\right)^{K1}\frac{\mathrm{\Gamma }((\beta 1)(K1))}{\mathrm{\Gamma }(\beta (K1)+2)}\stackrel{~}{V}^{\beta (K1)+1},$$
(44)
where
$$\stackrel{~}{V}=Vh(d)\tau \rho ^d,\beta =\frac{b}{d}.$$
(45)
The next step is to perform the grand-canonical sums over the instanton numbers $`K`$. As in the one-dimensional case, the sums can be evaluated in the large volume limit by replacing them by integrals which are calculated by means of the saddle point method. The error of this approximation vanishes in the thermodynamic limit. For the size distribution we get in this way
$$n(\rho )=\frac{Cd}{b}\rho ^{bd1}\mathrm{exp}(c\rho ^d),$$
(46)
with
$$c^{\frac{b}{d}}=C\sqrt{\frac{2\pi }{bd1}}\left(\frac{b}{d}1\right)^{\frac{b}{d}1}\mathrm{e}^{(\frac{b}{d}1)}h(d)\tau .$$
(47)
For $`bd`$ this takes the form
$$c^{\frac{b}{d}}=C\mathrm{\Gamma }(\frac{b}{d}1)d^{\frac{1}{2}}h(d)\tau ,$$
(48)
which agrees with the one-dimensional result.
The expression for $`n(\rho )`$ is consistent with the general expectation mentioned in the introduction: for small $`\rho `$ it grows powerlike with an exponent $`\alpha =bd1`$, and for large $`\rho `$ this power-law is combined with an exponential decrease.
Although the canonical partition functions are dominated by configurations where the instantons are densely packed, the exponent $`\alpha `$ agrees with the one of the semiclassical dilute gas approximation. This result does not depend on the details of our approximations and follows from the general structure of the occuring terms in the grand-canonical sums. The conjecture, made in , that due to the denseness of instantons the small-$`\rho `$ behaviour of $`n(\rho )`$ gets modified, is therefore wrong.
On the other hand, the value of the exponent $`p=d`$ in the exponential decay at large $`\rho `$ should be considered with reservations, because it depends on the saddle point approximations which have been made. In $`d=1`$ dimensions it is correct, but we would not be surprised, if in higher dimensions the true value would differ from $`d`$. In order to get more insight into this question and to get an idea of the quality of the approximations being made so far, we have also studied the instanton gas by grand canonical Monte Carlo simulations.
## 4 Monte Carlo simulations
Usually Monte Carlo calculations are done in the canonical ensemble. In our case the particle number has to change and it is necessary to simulate a grand canonical ensemble. Simulations of grand canonical systems are not very common. They are rarely discussed in the literature and some important details remain unclear. Therefore it appears appropriate to describe the algorithm we have used in our calculations. For related work on this topic we refer to .
### 4.1 Grand canonical Monte Carlo algorithms
A stochastic process, which is realized in a Monte Carlo simulation, is specified by a transition matrix $`W(X,Y)`$, where $`X`$ and $`Y`$ denote states of the system. For the purpose of a simulation $`W`$ is usually decomposed as a product of two factors: $`\omega (X,Y)`$ represents a proposal probability for a transition from $`X`$ to $`Y`$, and $`a_{XY}`$ denotes the corresponding acceptance probability. In addition to normalization and ergodicity one has to require stationarity, which is often fulfilled by demanding the stronger Metropolis condition of detailed balance:
$$a_{XY}=\mathrm{min}(1,\frac{\omega (Y,X)P(Y)}{\omega (X,Y)P(X)}).$$
(49)
Here $`P`$ is the probability distribution, which we want to generate as the stationary distribution of the underlying stochastic process. In our context it is given by
$$P_K(V;𝐚_1,\mathrm{},𝐚_K,\rho _1,\mathrm{},\rho _K)=\frac{1}{Z(V)}\frac{C^K}{K!}\underset{j=1}{\overset{K}{}}\rho _j^\alpha \mathrm{\Theta }(\{𝐚_j\},\{\rho _j\}).$$
(50)
The states $`X`$ and $`Y`$ are characterized by the instanton number $`K`$ combined with the set of coordinates $`\{𝐚_j,\rho _j\}`$.
In canonical algorithms $`\omega (X,Y)`$ is usually chosen to be symmetric so that it is omitted in (49) without further comments. This is not possible in a grand canonical ensemble, where one has to consider transitions that change the instanton number. Independent of the choice of $`\omega (X,Y)`$ there will be additional volume factors in $`a_{XY}`$ for processes that do not conserve the instanton number. This results from the asymmetry in particle creation and destruction. If an instanton is created one has to specify a probability for the generation of its new coordinates, On the other hand, in the process of removing an instanton such a probability does not occur. In our case, for the space-coordinates as well as for the radii we choose a uniform distribution within the allowed volume.
In the algorithm three different kinds of steps occur with equal probability: creation, destruction and movement of an instanton. With the shortcut notation
$$=\{\begin{array}{c}1.𝐚_i𝐚_j>\left(\frac{\tau }{v_1}\right)^{\frac{1}{d}}(\rho _i+\rho _j)i,j,\text{and}\hfill \\ 2.|a_i^\mu |<\frac{1}{2}V^{\frac{1}{d}}\mu ,i,\text{and}\hfill \\ 3.0<\rho _i<\frac{1}{2}\left(\frac{v_1}{\tau }V\right)^{\frac{1}{d}}i\hfill \end{array}$$
we have chosen the following transition rules, where $`x`$ denotes a random number between 0 and 1.
* Creation:
The creation of a new instanton with number $`K+1`$ and coordinates $`(𝐚^{},\rho ^{})`$ is proposed, and
$$X\{\begin{array}{c}Y,\frac{CV(Vv_1/\tau )^{\frac{1}{d}}}{2(K+1)}\rho _{}^{}{}_{}{}^{\alpha }x\text{ and }\hfill \\ X,\frac{CV(Vv_1/\tau )^{\frac{1}{d}}}{2(K+1)}\rho _{}^{}{}_{}{}^{\alpha }<x\text{ or not }.\hfill \end{array}$$
* Destruction:
The destruction of an instanton with randomly chosen number $`j`$ is proposed, and
$$X\{\begin{array}{c}Y,\frac{2K}{CV(Vv_1/\tau )^{\frac{1}{d}}}\rho _j^\alpha x\hfill \\ X,\frac{2K}{CV(Vv_1/\tau )^{\frac{1}{d}}}\rho _j^\alpha <x.\hfill \end{array}$$
* Movement:
A movement of a randomly chosen instanton in a volume element $`[\delta _a,\delta _a]^d\times [\delta _\rho ,\delta _\rho ]`$ around the original coordinates is proposed, and
$$X\{\begin{array}{c}Y,\left(\frac{\rho _j^{}}{\rho _j}\right)^\alpha x\text{ and }\hfill \\ X,\left(\frac{\rho _j^{}}{\rho _j}\right)^\alpha <x\text{ or not }.\hfill \end{array}$$
The simulations were started with the empty configuration ($`K=0`$). Measuring was started after the instanton number reached saturation.
### 4.2 Simulation results
In the case of $`d=1`$ dimensions the available exact result (32) provides a useful check of the Monte Carlo calculations. In Fig. 2 Monte Carlo data for $`n(\rho )`$ in $`d=1`$ are compared with the exact formula. The size $`L`$ has been chosen large enough such that finite $`L`$ effects are negligible. Obviously the Monte Carlo data agree very well with the theoretical predictions, and the thermodynamic limit has been approached sufficiently.
With this check on the Monte Carlo algorithm we proceed to the more interesting case of four space-time dimensions ($`d=4`$). In order to compare the Monte Carlo data with the outcome of our analytical approximations, Eq. (46), we have to make assumptions concerning the parameter $`h(d)=h(4)`$ that describes the ability of instantons to fill a given volume. We consider three choices, namely the lower bound $`h_1=1`$, the upper bound $`h_2=2^{41}=8`$ and their geometric mean $`h_g=2\sqrt{2}`$. The parameter $`\tau `$ is taken to be equal to $`v_1`$. For the volume we chose $`V=15^4`$. This is based on simulations in different volumes, which showed that in this case the thermodynamic limit was approximately reached within the errors of the simulation. The parameter $`\alpha `$ is taken to be $`\alpha =7/3`$, which is the value for SU(2) gauge theory in 4 dimensions.
Fig. 4 shows the Monte Carlo data in comparison with the analytical approximation. Near the maximum of the distribution the approximation qualitatively reproduces the Monte Carlo results. Furthermore, the growth of the distribution for small instanton radii according to a power law with exponent $`\alpha `$ can be confirmed, as is shown in Fig. 4.
In order to study the behaviour of $`n(\rho )`$ for large $`\rho `$ we considered the ratio
$$F(\rho )=\frac{n(\rho )}{\rho ^\alpha }.$$
(51)
Inspired by the theoretical results we tried fits of the form
$$F_{fit}(\rho )=a\mathrm{exp}(c\rho ^p).$$
(52)
The parameter $`a`$ was obtained by extrapolating $`F(\rho )`$ to small $`\rho `$. The fit with parameters $`c`$ and $`p`$ was then obtained using the Marquardt–Levenberg-algorithm. We performed fits for various choices of the model parameters $`C`$ and $`\alpha `$. In agreement with the theoretical results they showed that $`c`$ depends on $`\alpha `$, while $`p`$ is nearly independent of it.
The main interest is in the exponent $`p`$. We present the results for the parameter set $`C=1`$, $`\alpha =7/3`$, $`V=15^4`$, because this value of $`\alpha `$ is relevant for gauge theory with gauge group SU(2). For $`\alpha =6`$, the SU(3) case, the results for $`p`$ are the same within the present errors. We find $`a0.89`$, and the fit leads to $`c=3.3\pm 0.2`$ and $`p=1.9\pm 0.2`$. In Fig. 5 the result of a fit in the interval $`[0,2.25]`$ is shown.
This Gaussian ($`p=2`$) decay of the probability distribution has already been predicted by some authors under various assumptions . A recent approach based on an idea of dual superconductivity also leads to the prediction $`p=2`$. Furthermore good agreement with the $`SU(3)`$ lattice gauge theory calculations of Hasenfratz et al. was found.
The exponent $`p=2`$ differs from the one predicted by our approximate analytical calculation, $`p=d`$. The saddle point approximation being made is, however, not beyond any doubt. In that case the exponent originates from the effective excluded volume being proportional to $`\rho ^d`$ at the considered saddle point. This would also coincide with the intuitive expectation based on the following picture. In the presence of a large instanton of size $`\rho `$ the remaining ones are excluded from a volume $`\rho ^d`$. If they behave like a dilute gas, one would expect that the excluded volume yields a suppression factor $`\mathrm{exp}(c\rho ^d)`$. The instanton ensemble in the effective remaining volume is, however, dominated by dense configurations, as the analytical calculation shows. Therefore the intuitive picture should be considered with reservations. Indeed, the calculations of take into account excluded volume effects in the framework of the theory of grand canonical pair distribution functions and, also employing certain approximations, arrive at $`p=2`$.
In recent years much effort has been devoted to lattice Monte Carlo calculations of properties of the instanton ensemble. There are still ambiguities due to smoothing procedures and only data with little statistics are yet available. Nevertheless some quantitative statements have been given. Concerning the size distribution for small $`\rho `$, lattice calculations appear to support the power law (7) rather than (9). A nice plot, using data of , can be found in . For the large-$`\rho `$ distribution, de Forcrand et al. predict an exponential decrease with $`p=3\pm 1`$ from their $`SU(2)`$ lattice data . In contrast to this, Smith and Teper conclude form their $`SU(3)`$ simulations a decay according to $`\rho ^\xi `$ with $`\xi 10\mathrm{}12`$ .
We have studied the distribution of instanton sizes $`\rho `$ in the framework of a model, where instanton interactions are approximated by a hard core potential with variable radius. This model incorporates the basic features of a dynamical cut-off on large instanton sizes. In the one-dimensional case an exact formula can be derived, which yields a power-like growth $`\rho ^\alpha `$ for small radii $`\rho `$ and an exponential decay for large $`\rho `$.
In four space-time dimensions we employed analytical approximations as well as Monte Carlo simulations. The theoretical calculations generalize the one-dimensional results and give a power-like behaviour for small $`\rho `$. For large radii $`\rho `$ they overestimate the decay which is found in the Monte Carlo data. Fits to the numerical Monte Carlo results suggest a behaviour like
$$n(\rho )\stackrel{\rho \mathrm{}}{}\mathrm{exp}(c\rho ^2),$$
(53)
in agreement with some other work on gauge theories.
The results indicate that our simplified model reproduces the main features of instanton ensembles with a dynamical infrared cut-off. Definite results about properties of instanton ensembles can of course only be expected from future Monte Carlo calculations of lattice gauge theories.
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# Frustrated trimer chain model and Cu3Cl6(H2O)₂⋅2H8C4SO2 in a magnetic field
## I Introduction
The trimerized $`S=1/2`$ Heisenberg chain in a strong external magnetic field has already received a substantial amount of theoretical attention, one reason being a plateau at one third of the saturation magnetization in the magnetization curve . Some frustrated variants of the trimer model have also been investigated since they can be shown to have dimer groundstates and thus a spin gap.
While many materials with trimer constituents exist (see e.g. ), the behavior in high magnetic fields has been investigated only in a few of them, for instance in 3CuCl$`{}_{2}{}^{}`$2dioxane . Also Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub> belongs to the known trimer materials , but its behavior in a strong magnetic field has been measured only recently and at the same time its magnetic susceptibility has been remeasured. Surprisingly, a spin gap of about 3.9Tesla (that is roughly 5.5K) is observed both in the magnetic susceptibility of Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub> as well as in the magnetization as a function of external magnetic field. This system probably exhibits also a plateau at one third of the saturation magnetization in addition to the spin gap.
Motivated by the crystal structure , the authors of have proposed the following model (see also Fig. 1) :
$`H`$ $`=`$ $`J_1{\displaystyle \underset{i=1}{\overset{L/3}{}}}\left\{𝐒_{3i}𝐒_{3i+1}+𝐒_{3i+1}𝐒_{3i+2}\right\}`$ (4)
$`+J_2{\displaystyle \underset{i=1}{\overset{L/3}{}}}𝐒_{3i+2}𝐒_{3i+3}`$
$`+J_3{\displaystyle \underset{i=1}{\overset{L/3}{}}}\left\{𝐒_{3i+1}𝐒_{3i+3}+𝐒_{3i+2}𝐒_{3i+4}\right\}`$
$`h{\displaystyle \underset{i=1}{\overset{L}{}}}S_i^z.`$
Since the spin is localized on Cu<sup>2+</sup> ions, the $`𝐒_i`$ are spin-1/2 operators at site $`i`$. In (4), the reduced field $`h`$ is related to the physical field $`H`$ by $`h=g\mu _BH`$ in units where $`k_B=1`$. The numerical prefactor is determined by $`\mu _B0.67171`$K/Tesla as well as the value of $`g`$ which for the present material is slightly above 2 (the precise numerical value depends on the direction of the external magnetic field relative to the crystal axes).
For a study of the phase diagram of the Hamiltonian (4), it is useful to observe that the Hamiltonian and therefore also the phase diagram are invariant under the exchange of $`J_1`$ and $`J_3`$ such that one can concentrate e.g. on $`|J_3||J_1|`$. In fact, the $`h=0`$ phase diagram with antiferromagnetic exchange constants ($`J_i0`$) has been explored in using bosonization and exact diagonalization (see also ) determining in particular a parameter region with a spin gap. Very recently, this was complemented by a computation of the magnetization curve at some values of the parameters using DMRG . The investigations of concentrated on the region with all coupling constants in (4) antiferromagnetic ($`J_i0`$) because suggested that this should be appropriate for Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub>. However, the parameters relevant to the experimental system have not really been determined so far. We believe that this is an important issue in particular in view of the fact that according to the crystallographic data , all angles of the Cu–Cl–Cu bonds lie in the region of $`91^{}`$ to $`96^{}`$ – a region where usually no safe inference on the coupling constants $`J_i`$ can be made, not even about their signs. We will therefore develop a high-temperature series for the magnetic susceptibility of the model (4) and use it to determine the coupling constants from the experimental data . It will turn out that some coupling constants are likely to be ferromagnetic, i.e. the experimentally relevant coupling constants lie presumably outside the region studied so far. We then proceed to study more general properties of the model and to address the question of a spin gap in the relevant parameter region. We use mainly perturbative arguments supplemented by numerical methods.
Some supplementary results on the trimer model are contained in the appendices or can be found in .
## II Magnetic susceptibility and specific heat
### A High-temperature series for zero field
First we discuss some high-temperature series in zero magnetic field. We have used an elementary approach to perform the computations. Denote the Hamiltonian of a length $`L`$ chain with $`h=0`$ by $`H_0`$. Then the fundamental ingredient for any higher-order expansion is that contributions of $`\mathrm{tr}\left(H_0^N\right)`$ to suitable physical quantities become independent of the system size $`L`$ if one uses a long enough chain with periodic boundary conditions. The concrete Hamiltonian $`H_0`$ given by (4) must be applied $`2L/3`$ times to wind once around the system and to feel that it is finite. On the other hand, contributions from $`\mathrm{tr}\left(H_0^N\right)`$ with $`N<2L/3`$ are independent of $`L`$. We have used this observation to determine the high-temperature series by simply computing the traces for the lowest powers $`N`$ on a chain with a fixed $`L`$ and periodic boundary conditions . Just two small refinements to this elementary approach have been made. The first one is that we computed the traces separately for all subspaces of the $`z`$-component of the total spin $`S_{\mathrm{tot}}^z`$. This is already sufficient to obtain series for the specific heat $`c_v`$ and the magnetic susceptibility $`\chi `$. The second one is to make also the order $`2L/3`$ usable: At this order, only the coefficient of $`J_1^{L/3}J_3^{L/3}`$ is affected by the finiteness of the chain and this coefficient can be corrected by hand using results for a Heisenberg ring of length $`2L/3`$.
For notational convenience, we introduce the partition function for $`L`$ sites by
$$Z_L=\mathrm{tr}\left(\mathrm{e}^{\beta H_0}\right)$$
(5)
with $`k_BT=1/\beta `$.
The lowest orders of a reduced magnetic susceptibility $`\chi `$ are found to be
$`\chi _{\mathrm{red}.}(\beta )`$ $`=`$ $`{\displaystyle \frac{1}{\beta LZ_L}}{\displaystyle \frac{^2}{h^2}}\mathrm{tr}\left(\mathrm{e}^{\beta \left(H_0hS_{\mathrm{tot}}^z\right)}\right)|_{h=0}={\displaystyle \frac{\beta }{L}}{\displaystyle \frac{\mathrm{tr}\left(\left(S_{\mathrm{tot}}^z\right)^2\mathrm{e}^{\beta H_0}\right)}{Z_L}}`$ (6)
$`=`$ $`{\displaystyle \frac{\beta }{4}}{\displaystyle \frac{\beta ^2}{24}}\left(2J_3+J_2+2J_1\right){\displaystyle \frac{\beta ^3}{96}}\left(J_1^2+J_2^2+J_3^26J_3J_12J_2J_12J_3J_2\right)`$ (11)
$`+{\displaystyle \frac{\beta ^4}{1152}}\left(8J_3^33J_3^2J_2+6J_2^2J_1+6J_3J_2^2+8J_1^33J_2J_1^2+J_2^3\right)`$
$`+{\displaystyle \frac{\beta ^5}{4608}}(28J_3^2J_2J_128J_3J_1^2J_2+36J_3^2J_1^2+8J_3^2J_2^2+8J_2^2J_1^22J_2^3J_134J_3J_1^334J_3^3J_1`$
$`10J_2J_1^32J_3J_2^310J_3^3J_228J_3J_1J_2^2+14J_1^4+14J_3^4+5J_2^4)`$
$`+𝒪\left(\beta ^6\right).`$
Similarly, we obtain the lowest orders of the high-temperature series for the specific heat
$`{\displaystyle \frac{c_v(\beta )}{k_B}}`$ $`=`$ $`{\displaystyle \frac{\beta ^2}{L}}{\displaystyle \frac{^2}{\beta ^2}}\mathrm{ln}\left(Z_L\right)`$ (12)
$`=`$ $`{\displaystyle \frac{\beta ^2}{16}}\left(2J_1^2+J_2^2+2J_3^2\right)+{\displaystyle \frac{\beta ^3}{32}}\left(2J_3^3+2J_1^3+J_2^36J_3J_1J_2\right)`$ (15)
$`{\displaystyle \frac{\beta ^4}{256}}\left(8J_3^2J_2J_1+8J_3J_1^2J_2+12J_3^2J_1^2+8J_3^2J_2^2+8J_2^2J_1^2+6J_1^4+6J_3^4+J_2^4+8J_3J_1J_2^2\right)`$
$`+𝒪\left(\beta ^5\right).`$
Complete 12th order versions of both series can be accessed via .
For a uniform Heisenberg chain ($`J_1=J_2`$ and $`J_3=0`$), the coefficients of the series for $`\chi `$ and $`\mathrm{ln}\left(Z_L\right)/L`$ (or $`c_v`$) agree with those given for instance in when they overlap.
### B Fit to the experimental susceptibility
Now we use our 12th order series for the susceptibility (11) to fit the experimental data and thus extract values for the coupling constants $`J_i`$. We used the data for the single crystal ($`Hb`$-axis) and the polycrystalline sample as well as some unpublished new measurements for all three axes of a single crystal . For the polycrystalline case the average $`g`$-factor is known to be $`g_{\mathrm{av}}2.1`$ from ESR while in the single crystal case we used the $`g`$-factor as a fitting parameter. The following prefactors are used to match the series (11) to the experimental data:
$$\chi _{\mathrm{exp}.}(T)=\frac{3N_Ag^2\mu _B^2}{k_B}\chi _{\mathrm{red}.}\left(\frac{1}{T}\right).$$
(16)
We performed fits in various intervals of temperature with a lower boundary ($`T_l`$) lying between 150K and 250K, while the upper boundary was kept fixed at 300K. Fits were performed with the raw 12th order series. For both experimental data sets of we obtained reasonable, though volatile fits around $`T_l=150\text{K}250\text{K}`$ yielding the following estimates: $`J_1=250\text{K}\pm 40\text{K}`$, $`J_2=250\text{K}\pm 40\text{K}`$, $`J_3=40\text{K}\pm 30\text{K}`$. For the single crystal sample we additionally determined $`g_b=1.95\pm 0.05`$.
We have further performed fits to unpublished single-crystal data sets where the $`g`$-factors are known from ESR . When a constant is added to (11), the data for all three crystal axes can be fitted consistently with $`J_1300`$K, $`J_2280`$K and $`J_360`$K in an interval of high temperatures ($`220\mathrm{K}T300`$K). This set of coupling constants is in agreement with our earlier fits and we will use the latter in the further discussion below.
Fig. 2 shows the measured susceptibility for the polycrystalline sample together with the series result. Since the parameters were obtained from a fit which was performed with a different data set, we have used $`g=2.03`$ (which differs slightly from the experimentally found $`g_{\mathrm{av}}2.1`$) in order to obtain agreement of the raw series with the experimental data for $`T240`$K. Clearly, the raw series should not be trusted down into the region of the maximum of $`\chi `$ where Padé approximants should be used instead. The region below the maximum cannot be expected to be described with a high-temperature series. The overall agreement is reasonable though the theoretical result reproduces the experimental one in the vicinity of the maximum only qualitatively. This discrepancy might be due to the frustration in the model which leads to cancellations in the coefficients. Note also that, due to the frustration, the maximum of $`\chi `$ is located at a lower temperature than would be expected for a non-frustrated model with coupling constants of the same order of magnitude. Consequently, higher orders are important in the entire temperature range covered by Fig. 2, precluding in particular the analysis of the high-temperature tail of $`\chi `$ in terms of a simple Curie-Weiss law.
The agreement for intermediate temperatures can be improved if the maximum is included in the fitting region and Padé approximants are used in the fit. The main change with respect to the fits discussed above is that $`J_3`$ tends to be closer to $`J_1`$. However, it will become clear from the discussion in later sections that the region with $`J_3`$ close to $`J_1`$ is not appropriate to describe the experimental observations of the low-temperature region.
Although we are not able to determine the coupling constants to high accuracy, all our fits lead to the conclusion that $`J_2`$ should be antiferromagnetic and $`J_1`$ and $`J_3`$ (or at least one of them) must be ferromagnetic if one wants to model the susceptibility measured at high temperatures with the frustrated trimer chain (4). In view of earlier theoretical investigations , this conclusion is somewhat surprising. Note that none of our fits converged to all $`J_i>0`$. Additional assumptions (including a constraint on the $`J_i`$) are necessary to determine from $`\chi (T)`$ what the optimal values of the $`J_i`$ would be in this antiferromagnetic region and thus allow for a comparison with Fig. 2. Such a fit and a comparison with the present one is discussed in appendix A. The upshot is that the experimentally observed $`\chi (T)`$ cannot be explained with only antiferromagnetic $`J_i`$.
The findings of this section necessitate a detailed re-analysis of the Hamiltonian (4) since earlier works did not look at the appropriate parameter region.
## III Lanczos results
In order to study the zero-temperature behavior of the frustrated trimer chain we have performed Lanczos diagonalizations of small clusters with periodic boundary conditions. Although computations were performed for various values of the parameters, we will present explicit results only for the final parameter set determined above. Further results in the region $`J_i>0`$ are in agreement with and are used in appendix A.
Fig. 3 presents the zero-temperature magnetization curve for the trimer chain model. Here and below the magnetization $`M`$ is normalized to saturation values $`\pm 1`$. First, it is reassuring that the system still has antiferromagnetic features despite two ferromagnetic coupling constants (note that we are now probing a region far from that used for determining the $`J_i`$). Since experiments found a spin gap , an important question clearly is if we also obtain a gap from the model with these parameters. We have therefore performed a finite-size analysis of the gap to $`S^z=1`$ excitations (corresponding to the first step in the finite-size magnetization curves of Fig. 3). All our approaches led to results compatible with a vanishing gap. However, it is difficult to reliably exclude a gap of a few K with system sizes $`L30`$. We will therefore return to this issue later and assume for the extrapolated thick line in Fig. 3 a vanishing spin gap. In general, this extrapolation was obtained by connecting the mid-points of the steps of the $`L=30`$ magnetization curve, except for $`M=1`$ and $`M=1/3`$ where the corners were used.
For the parameters of Fig. 3, $`M=1/3`$ is reached with a magnetic field $`H=2025`$Tesla. The order of magnitude agrees with the experimental finding even if the value found within the model is a factor of two to three below the experimental one. Above this field, Fig. 3 exhibits a clear $`M=1/3`$ plateau which is expected on general grounds .
We conclude this section by presenting in Fig. 4 the lowest three excitations for the $`S^z=1`$ sector as a function of momentum $`k`$, where $`k`$ is measured with respect to the groundstate, i.e. $`k=k_{S^z=1}k_{\mathrm{GS}}`$. This spectrum is very similar to that of an $`S=1/2`$ Heisenberg chain of length $`L/3`$ with coupling constant $`J_{\mathrm{eff}.}16`$K. In particular, one can recognize the two-spinon scattering continuum and a few higher excitations. This identification of the low-energy excitations of the frustrated trimer chain with an effective $`S=1/2`$ Heisenberg chain is one of the numerical indications for the absence of a spin gap.
## IV The line $`J_1=J_3`$
The Lanczos results of the previous section raise the question if the trimer chain model has a spin gap in the region with $`J_1`$, $`J_3<0`$: In this region, the model behaves like an antiferromagnet (see e.g. Fig. 3) which is frustrated since the number of antiferromagnetic coupling constants around a triangle is odd. Therefore, a spin gap appears possible in principle and we proceed with further arguments to decide whether it appears in the relevant parameter region.
Evidence for a spin gap in the parameter region $`J_1`$, $`J_3>0`$ was actually first obtained on the line $`J_1=J_3`$ . The reason is presumably that the line $`J_1=J_3`$ can be treated analytically at least to some extent because then the total spin is locally conserved on each bond coupled by $`J_2`$. In fact, one can easily discuss the entire magnetization process and not just the question of a spin gap and we refer the interested reader to for some comments on this aspect.
Recall that the model (4) with $`h=0`$ gives rise to three types of groundstates in different regions with $`J_1=J_3>0`$ (we will assume $`J_2>0`$ throughout this section) : That of the $`S=1/2`$-$`S=1`$ ferrimagnetic chain for $`J_2<0.90816J_1`$, a spontaneously dimerized state for $`0.90816J_1<J_2<2J_1`$ and, finally, singlets are formed on all bonds coupled by $`J_2`$ for $`J_2>2J_1`$ with effectively free spins in between. In the first and the third case, one finds ferrimagnetic behavior with a spontaneous magnetization $`M=1/3`$ and only the second region exhibits the requested gap.
For $`J_1=J_3<0`$, we found only two regions:
1. For
$$J_2>J_1$$
(17)
the groundstate is formed by singlets on the $`J_2`$-bonds and free $`S=1/2`$ spins in between. This again gives rise to ferrimagnetic behavior with a spontaneous magnetization $`M=1/3`$.
2. When
$$J_2<J_1$$
(18)
the entire system behaves like a ferromagnet. In this case the system is spontaneously completely polarized ($`M=1`$).
We conclude that –unlike for $`J_1=J_3>0`$– the groundstate is always gapless for $`J_1=J_3<0`$ which we have argued to be more appropriate for Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub>.
## V Effective Hamiltonians for the groundstate
The low-energy behavior of the model Hamiltonian (4) can be analyzed further using degenerate perturbation theory. Truncation at a certain order of the coupling constants leads to effective Hamiltonians which in some cases turn out to be well-known models.
We will use the abbreviations
$$\stackrel{~}{J}_i=\frac{J_i}{J_1},\overline{J}_i=\frac{J_i}{J_2}.$$
(19)
### A $`J_1`$ large and antiferromagnetic
To test the method, we first consider the case of antiferromagnetic $`J_1`$. For this purpose we extend the first-order effective Hamiltonian of for the case $`J_1J_2`$, $`J_30`$ to second order. For $`J_1`$ large and antiferromagnetic, the groundstate-space of a trimer is given by an $`S=1/2`$ representation. In this subspace of doublets, the effective Hamiltonian has the form of a $`J_1`$-$`J_2`$ chain when truncated after the second order:
$$H_{\mathrm{eff}.}=𝒥_1\underset{i}{}𝐒_i𝐒_{i+1}+𝒥_2\underset{i}{}𝐒_i𝐒_{i+2}.$$
(20)
Here, the $`𝐒_i`$ are effective spin-1/2 operators. The effective exchange constants are
$$\frac{𝒥_1}{J_1}=\frac{4}{9}\left(\stackrel{~}{J}_3+\stackrel{~}{J}_2\right)\frac{79}{405}\stackrel{~}{J}_3^2+\frac{8}{135}\stackrel{~}{J}_2\stackrel{~}{J}_3+\frac{211}{1620}\stackrel{~}{J}_2^2$$
(21)
and
$$\frac{𝒥_2}{J_1}=\frac{91}{486}\stackrel{~}{J}_3^2+\frac{22}{243}\stackrel{~}{J}_2\stackrel{~}{J}_3+\frac{10}{243}\stackrel{~}{J}_2^2.$$
(22)
In this approximation, the ferrimagnetic phase found in is given by an effective ferromagnetic Hamiltonian ($`𝒥_1<0`$) while the antiferromagnetic phase corresponds to $`𝒥_1>0`$. The transition line can thus be determined from $`𝒥_1=0`$ . We find
$`\stackrel{~}{J}_3`$ $`=`$ $`{\displaystyle \frac{12}{79}}\stackrel{~}{J}_2{\displaystyle \frac{90}{79}}+{\displaystyle \frac{1}{158}}\sqrt{17245\stackrel{~}{J}_2^2+48240\stackrel{~}{J}_2+32400}`$ (23)
$`=`$ $`\stackrel{~}{J}_2{\displaystyle \frac{1}{80}}\stackrel{~}{J}_2^2+𝒪\left(\stackrel{~}{J}_2^3\right),`$ (24)
which improves the agreement of the approximation $`\stackrel{~}{J}_3=\stackrel{~}{J}_2`$ with the numerical results of .
The dimer phase with a spin gap is characterized by $`𝒥_2/𝒥_1>0.241167(5)`$ (see and references therein). Using (21) and (22), it is found to open at $`\stackrel{~}{J}_23.60`$, $`\stackrel{~}{J}_31.361`$ with a square-root like behavior of $`\stackrel{~}{J}_3`$ as a function of $`\stackrel{~}{J}_2`$. Since this is not in the weak-coupling region, it is not surprising that the numbers differ substantially from those obtained numerically in . However, the topology of the groundstate phase diagram comes out correctly from our effective Hamiltonian: In particular, the dimerized spin gap phase is located inside the antiferromagnetic phase and arises because of a sufficiently large effective second neighbor frustration $`𝒥_2`$.
### B $`J_2`$ large and antiferromagnetic
The preceding argumentation is not applicable to the region $`J_2>0`$, $`J_1,J_3<0`$. However, a similar case has been discussed earlier and $`J_2|J_1|,|J_3|`$ has been found to be a useful limiting case. We will now analyse this region in the same manner as above.
For $`J_2|J_1|,|J_3|`$, the spins on all $`J_2`$-bonds couple to singlets and only the intermediate spins contribute to the low-energy excitations. In the space of these intermediate spins, we can again map the Hamiltonian (4) to the Hamiltonian (20) to the lowest orders in $`J_1`$, $`J_3`$. Up to fifth order, we find the effective coupling constants to be given by
$`{\displaystyle \frac{𝒥_1}{J_2}}`$ $`=`$ $`(\overline{J}_1\overline{J}_3)^2\{{\displaystyle \frac{1}{2}}+{\displaystyle \frac{3\left(\overline{J}_1+\overline{J}_3\right)}{4}}+3\overline{J}_1\overline{J}_3`$ (26)
$`{\displaystyle \frac{\left(\overline{J}_1+\overline{J}_3\right)\left(107\left(\overline{J}_{1}^{}{}_{}{}^{2}+\overline{J}_{3}^{}{}_{}{}^{2}\right)406\overline{J}_1\overline{J}_3\right)}{64}}\}`$
and
$$\frac{𝒥_2}{J_2}=\frac{\left(\overline{J}_1+\overline{J}_3\right)\left(\overline{J}_1\overline{J}_3\right)^4}{4}.$$
(27)
This mapping is now applicable regardless of the sign of $`J_1`$ and $`J_3`$ as long as $`J_2>0`$. First we consider the case of antiferromagnetic $`J_1,J_3>0`$. Then the effective coupling constants are essentially always antiferromagnetic, i.e. $`𝒥_1,𝒥_2>0`$ leading to a frustrated chain. If $`J_1`$ and $`J_3`$ are large enough, $`𝒥_2/𝒥_1`$ can exceed the critical value of about $`0.241`$ (see above) and a spin gap opens. These observations are again in qualitative agreement with the phase diagram of . As for the preceding limit, one should not expect good quantitative agreement since the required values of $`J_1`$ and $`J_3`$ are not small but of the same order as $`J_2`$.
Now we turn to the more interesting case $`J_1,J_3<0`$. Then the oupling constant (27) is always ferromagnetic: $`𝒥_2<0`$. If $`|J_1|`$ and $`|J_3|`$ are large enough, $`𝒥_1`$ also becomes ferromagnetic. This is compatible with the behavior found in section IV on the line $`J_1=J_3<0`$. If $`|J_1|`$ and $`|J_3|`$ are small, $`𝒥_1`$ remains antiferromagnetic. Since $`𝒥_2`$ is always ferromagnetic, no frustration arises in the effective model and a spin gap is not expected to open. This is true to the order which we have considered. Higher orders might actually yield frustrating contributions. In any case, frustration is substantially weaker for ferromagnetic $`J_1,J_3<0`$ than for antiferromagnetic $`J_1,J_3>0`$. It is therefore plausible that a spin gap is absent in the ferromagnetic region (unless $`|J_1|`$ and/or $`|J_3|`$ are very large and the present argument is not applicable).
It should be noted that (26) and (27) turn out to be small if $`\overline{J}_1\overline{J}_3`$ is small. In fact, one can argue that the results of this section remain qualitatively correct for $`\overline{J}_1\overline{J}_3`$ small even if $`\overline{J}_1`$ and $`\overline{J}_3`$ are not separately small: For $`J_1=J_3`$, the intermediate spins are effectively decoupled due to the presence of the singlets on the $`J_2`$-bonds (see section IV). A small detuning $`J_1J_3`$ generates an effective coupling of the intermediate spins via higher-order processes. However, the effective coupling will stay small as long as $`J_1J_3`$ is small. If one wants to model Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub>, $`|J_1J_3|`$ must therefore at least be on the same scale as e.g. the field $`h80`$K required to polarize the intermediate spins leading to $`M=1/3`$ . This observation rules out a $`J_1`$ very close to $`J_3`$.
Finally, we also calculated the effective Hamiltonian for a strong ferromagnetic intra-trimer interaction $`J_1`$. The problem then maps to a frustrated $`S=3/2`$ chain with four-spin interactions. Even if this is not a well-known Hamiltonian and the issue of a spin gap thus remains unclear in this case, we present it in appendix B in order to open the way for further investigation of this limit.
## VI Magnetization plateaux
We complete our theoretical analysis with a discussion of plateaux in the magnetization curves of the frustrated trimer chain model.
A plateau with $`M=1/3`$ is abundant in the magnetization curve (compare Fig. 3) and can be easily understood in the limits $`|J_2|,|J_3|J_1`$ or $`|J_1|,|J_3|J_2`$. This is readily done by adding the coupling $`J_3`$ to the series of . More details as well as the explicit series for the boundaries of the $`M=1/3`$ plateau are available under . Here we just mention that the main conclusions of regarding this plateau remain qualitatively unchanged in the presence of the additional coupling $`J_3`$.
Regarding plateaux with $`M1/3`$, observe first that, when a spin gap opens in the frustrated trimer model, the groundstate is dimerized, i.e. translational invariance is spontaneously broken by a period two. Spontaneous breaking of translational invariance by a period two also permits the appearance of a plateau with $`M=2/3`$ (see and references therein). We will now investigate this possibility further.
First we consider the case $`J_1>0`$ and start in the limit of strong trimerization ($`J_2=0,J_3=0`$). When one applies a magnetic field $`h_c=\frac{3}{2}J_1`$, the two states $`|`$ and $`\frac{1}{\sqrt{6}}(|2|+|)`$ are degenerate in energy. This degeneracy is then lifted by the couplings $`J_2,J_3`$. The effective Hamiltonian to first order is an $`XXZ`$ chain in a magnetic field . We obtain the following effective couplings for the $`XXZ`$ chain:
$`J_{xy}`$ $`=`$ $`{\displaystyle \frac{1}{6}}J_2{\displaystyle \frac{2}{3}}J_3`$ (28)
$`J_z`$ $`=`$ $`{\displaystyle \frac{1}{36}}(J_2+8J_3)`$ (29)
$`h_{\mathrm{eff}}`$ $`=`$ $`hh_c{\displaystyle \frac{1}{36}}(5J_2+22J_3)`$ (30)
and therefore the effective anisotropy $`\mathrm{\Delta }_{\mathrm{eff}}=J_z/|J_{xy}|`$ is
$$\mathrm{\Delta }_{\mathrm{eff}}=\frac{J_2+8J_3}{|6J_224J_3|}.$$
(31)
For $`5/32<J_3/J_2<7/16`$, we have $`\mathrm{\Delta }_{\mathrm{eff}}>1`$ and thus a gap, i.e. an $`M=2/3`$ plateau in the original model. A plateau with $`M=2/3`$ can be indeed observed numerically somewhere in this region (see e.g. ). The line $`J_3/J_2=1/4`$ describes the Ising limit $`\mathrm{\Delta }_{\mathrm{eff}}=\mathrm{}`$.
In order to address the region of ferromagnetic $`J_1`$, we now start from the limit $`J_1=J_3=0`$ and apply a magnetic field $`h_c=J_2`$. Then the two states $`|`$ and $`\frac{1}{\sqrt{2}}\left(\right||)`$ on the $`J_2`$-dimer become degenerate in energy while the intermediate spins are already polarized. This can be again treated by degenerate perturbation theory in $`1/J_2`$. Up to third order we find an $`XXZ`$ chain with
$`{\displaystyle \frac{J_{xy}}{J_2}}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(2+\overline{J}_1+\overline{J}_3\right)\left(\overline{J}_1\overline{J}_3\right)^2`$ (32)
$`{\displaystyle \frac{J_z}{J_2}}`$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(\overline{J}_1+\overline{J}_3\right)\left(\overline{J}_1\overline{J}_3\right)^2`$ (33)
$`{\displaystyle \frac{h_{\mathrm{eff}}}{J_2}}`$ $`=`$ $`{\displaystyle \frac{h}{J_2}}1{\displaystyle \frac{1}{2}}\left(\overline{J}_1+\overline{J}_3\right){\displaystyle \frac{1}{4}}\left(\overline{J}_1\overline{J}_3\right)^2`$ (34)
that is
$$\mathrm{\Delta }_{\mathrm{eff}}=\frac{\overline{J}_1+\overline{J}_3}{|2+\overline{J}_1+\overline{J}_3|}.$$
(35)
In the region where this treatment is valid, we always have a small $`\mathrm{\Delta }_{\mathrm{eff}}`$, i.e. no plateau at $`M=2/3`$. Indeed, one can see that the dimer excitations can hop at second order in $`1/J_2`$ while up to this order all diagonal terms involve only a single dimer site. Thus, up to second order the diagonal terms contribute only to $`h_{\mathrm{eff}}`$ and to this order one obtains an $`XY`$ chain in a magnetic field. A small anisotropy is restored at third order before terms that are not described by a simple $`XXZ`$ chain arise at fourth order.
## VII Conclusions
We have studied the frustrated trimer chain (4) (Fig. 1) using a variety of methods. First, we have computed 12th-order high-temperature series for the susceptibility $`\chi `$ and specific heat. Fits of the high-temperate tail of the susceptibility computed from the model to the one measured on Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub> lead to $`J_2=250\mathrm{K}\pm 40\mathrm{K}`$ and ferromagnetic $`J_1=260\mathrm{K}\pm 50\mathrm{K}`$, $`J_3=40\mathrm{K}\pm 30\mathrm{K}`$ (we showed in appendix A that $`\chi (T)`$ cannot be fitted with the antiferromagnetic parameters proposed in ). We assumed that these parameters remain valid down to low temperatures since we are not aware of any indication of a drastic change in the magnetic behavior of Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub> as temperature is lowered. In fact, features of other experimental observations at intermediate and low temperatures are roughly reproduced with the aforementioned parameters: We find a maximum in $`\chi (T)`$ in the region $`50\mathrm{K}T100\mathrm{K}`$ and a smooth increase of the low-temperature magnetization $`M`$ from $`0`$ to $`1/3`$ as the external magnetic field is increased from zero to several ten Tesla. From a quantitative point of view, the agreement may however not yet be entirely satisfactory: Deviations between the measured susceptibility from the one obtained within the model can be seen in the interval $`80\mathrm{K}T200\mathrm{K}`$ and the model predicts an $`M=1/3`$ magnetization for a magnetic field that is a factor two to three below the one actually required in the experiment.
Probably the most exciting experimental observation for Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub> is the existence of a spin gap of about 5.5K. We have therefore searched for a spin gap in the region of ferromagnetic $`J_1`$ and $`J_3`$ using several methods. Neither Lanczos diagonalization, discussion of the line $`J_1=J_3`$ nor an effective Hamiltonian for large $`J_2`$ provide any evidence in favor of a spin gap in this parameter region. A further careful analysis of this issue would certainly be desireable in particular in view of the small size of the actually observed gap. At present, however, it seems likely that the model does not reproduce a spin gap in the relevant parameter region.
It should be noted that the coupling constants which we have determined are about two orders of magnitude larger than the experimentally observed gap. Therefore, a small modification of the model is sufficient to produce a gap of this magnitude. The possibilities include dimerization of the coupling constants, exchange anisotropy as well as additional couplings. A modification of the model along these lines may also help to improve the quantitative agreement with the features observed in Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub> at energy scales of about 100K. Further measurements are however needed to discriminate between these possibilities. For example, it would be interesting to measure the specific heat and compare it with our series (15). It emerges also from our analysis that a temperature of 300K is still too small to allow for application of a simple Curie-Weiss law to the magnetic susceptibility $`\chi `$. It would therefore be useful to measure $`\chi `$ to higher temperatures in order to permit analysis via truncation of (11) after the order $`T^2`$ which would provide a more direct check that $`2(J_1+J_3)+J_2`$ is negative.
However, inelastic neutron scattering would presumably be most helpful: First, this should clearly decide if Cu<sub>3</sub>Cl<sub>6</sub>(H<sub>2</sub>O)$`{}_{2}{}^{}`$2H<sub>8</sub>C<sub>4</sub>SO<sub>2</sub> is really quasi-one-dimensional and secondly it would yield direct information on the excitation spectrum which could hopefully be interpreted in terms of coupling constants. Such a determination of the coupling constants would also circumvent the question whether model parameters change as a function of temperature since neutron scattering is carried out at low temperatures, i.e. the temperature scale of interest. We therefore hope that neutron scattering can indeed be performed and are curious if excitations will be observed that are similar to those computed in the trimer chain model (Fig. 4).
The frustrated trimer chain model is also interesting in its own right: It has a rich phase diagram which among others includes many aspects of the $`J_1`$-$`J_2`$ chain such as a frustration-induced spin gap in some parameter region . Also plateaux in the magnetization curve exist in this model: A plateau with $`M=1/3`$ is abundant both in the regions with antiferromagnetic and ferrimagnetic $`h=0`$ groundstates. Also a plateau with $`M=2/3`$ can be shown to exist in the region with $`J_1,J_3>0`$ (see and section VI). Like in the case of the spin gap, the opening of the latter plateau is accompanied by spontaneous breaking of translational invariance in the groundstate. Amusingly, however, the $`M=2/3`$ plateau opens already for $`J_2,J_3J_1`$ – a region where the spin gap is absent. In this context of magnetization plateaux, we hope that the magnetization measurements can be extended to slightly higher fields which should unveil the lower edge of the $`M=1/3`$ plateau.
## Acknowledgments
We are very grateful to M. Ishii and H. Tanaka for providing us with their partially unpublished data for the susceptibility and for discussions. In addition, useful discussions with D.C. Cabra, F. Mila, M. Troyer and T.M. Rice are gratefully acknowledged. A.H. is indebted to the Alexander von Humboldt-foundation for financial support during the initial stages of this work as well as to the ITP-ETHZ for hospitality.
## A Antiferromagnetic coupling constants
In this appendix we discuss a fit of the magnetic susceptibility $`\chi `$ with antiferromagnetic coupling constants $`J_i>0`$. A number of assumptions are necessary in order to obtain at all a convergent fit with parameters in the antiferromagentic spin gap region .
First we fix the ratio of the magnetic field $`h(M=1/3)`$ to the spin gap $`h_c(M=0)`$ approximately to the experimental value
$$\frac{h(M=1/3)}{h_c(M=0)}=14.1.$$
(A1)
To this end, we used numerical data for $`h(M=1/3)`$ and $`h_c(M=0)`$ on systems of size $`L=12`$, $`18`$ and $`24`$. This data was extrapolated to $`L=\mathrm{}`$ in the same manner as in , i.e. with a polynomial fit $`h_L(M)=h_{\mathrm{}}(M)+a/L+b/L^2`$. For the spin gap, this amounts to reproducing the computation of . The numerical solutions to eq. (A1) were then approximated by
$$\stackrel{~}{J}_3=0.3\left(\stackrel{~}{J}_20.85\right)^2+0.63$$
(A2)
where we used the notation (19). An analytic formula was needed in order to implement the constraint (A1) by inserting (A2) into (11) before performing a fit. Eq. (A2) is valid for $`0.6\stackrel{~}{J}_22`$.
The constraint (A2) is still not sufficient to ensure antiferromagnetic $`J_i>0`$ with $`0.6\stackrel{~}{J}_22`$. To achieve this goal, we had to make the following further adjustments when fitting our series (11) to the experimental data :
1. Keep $`g`$ as a fitting parameter,
2. add a constant to (11) and use this as another parameter in the fit,
3. start fitting at low temperatures $`T_l100`$K.
Note that both $`g`$ and the additive constant turn out to be quite large. For example, for the parameters used in Fig. 5, we found $`g2.9`$ and an additive constant of about $`0.1710^3\mathrm{K}^1/k_B`$. This means that the prefactor in (16) is off by a factor of about 2 from the value determined by ESR and that the absolute value of the additive constant is almost 40% of the susceptibility observed at $`T=300`$K !
On the basis of these unrealistic parameters, one could already discard this fit to $`\chi (T)`$. Nevertheless, we compare it to the one shown in Fig. 2: Fig. 5 shows the measured susceptibility for the polycrystalline sample together with the series evaluated at $`J_1=120`$K, $`J_2=141`$K and $`J_3=79`$K. This parameter set is close to parameters proposed in . This proposal was based on two assumptions: 1) The model should give rise to the experimentally observed spin gap of around 5K . 2) The maximum of $`\chi `$ is located at $`T0.7J_1`$. While we do indeed reproduce the spin gap rather accurately, the second assumption is falsified by our computation: The frustration pushes the maximum of $`\chi (T)`$ again to lower temperatures as compared to a non-frustrated system.
Fig. 5 should be compared to Fig. 2. The seemingly better agreement in the region $`100\mathrm{K}T200`$K is due to including this temperature interval in the fit for Fig. 5, but not in Fig. 2. Note that the $`|J_i|`$ are now smaller by a factor of about two than those used in Fig. 2. One would therefore expect better convergence in the vicinity of the maximum of $`\chi (T)`$, i.e. for $`50\mathrm{K}T100`$K. This expectation is confirmed by the fact that in Fig. 5, the , and Padé approximants are indistinguishable. However, while the series reproduces the maximum roughly in Fig. 2, this is definitely not the case in Fig. 5. The better agreement of the fit in Fig. 2 with the experimental data at $`T70`$K is particularly remarkable since this temperature range is far from the fitting region in this case, while closeby temperatures were used in Fig. 5. In combination with the unrealistic assumptions needed to obtain a convergent fit with all $`J_i>0`$ one can therefore conclude safely that only antiferromagnet coupling constants are not suitable for describing the experimental data for the susceptibility $`\chi (T)`$.
## B Effective Hamiltonian for $`J_1`$ large and ferromagnetic
For a strong ferromagnetic intra-trimer interaction $`J_1`$, the noninteracting groundstates are built from products of trimer $`S=3/2`$ states. Up to second order we find the following effective Hamiltonian in this subspace of low-lying trimer quartets:
$`H_{\mathrm{eff}.}`$ $`=`$ $`𝒥_a{\displaystyle \underset{i}{}}𝐒_i𝐒_{i+1}+𝒥_b{\displaystyle \underset{i}{}}𝐒_i𝐒_{i+2}`$ (B4)
$`+𝒥_c{\displaystyle \underset{i}{}}\left(𝐒_i𝐒_{i+1}\right)^2`$
$`+{\displaystyle \frac{𝒥_d}{2}}{\displaystyle \underset{i}{}}\{(𝐒_i𝐒_{i+1})(𝐒_{i+1}𝐒_{i+2})`$
$`+(𝐒_{i+2}𝐒_{i+1})(𝐒_{i+1}𝐒_i)\},`$
where the $`𝐒_i`$ are now effective spin-3/2 operators.
The coupling constants are found to be:
$`𝒥_a`$ $`=`$ $`{\displaystyle \frac{1}{9}}(J_2+2J_3)+{\displaystyle \frac{197J_2^2+212J_2J_3+212J_3^2}{2592|J_1|}},`$ (B5)
$`𝒥_b`$ $`=`$ $`{\displaystyle \frac{2J_2^2+5J_2J_3+2J_3^2}{27|J_1|}},`$ (B6)
$`𝒥_c`$ $`=`$ $`{\displaystyle \frac{41J_2^2+100J_2J_3+36J_3^2}{1296|J_1|}},`$ (B7)
$`𝒥_d`$ $`=`$ $`{\displaystyle \frac{4(2J_2^2+5J_2J_3+2J_3^2)}{243|J_1|}}.`$ (B8)
Even if this effective Hamiltonian is not a well-known one, it is clear that there is no spin gap in first order, since then the system is effectively a nearest neighbor $`S=3/2`$ Heisenberg chain which is either gapless ($`J_2+2J_3>0`$) or ferromagnetic ($`J_2+2J_3<0`$).
If one neglects the $`𝒥_c`$ and $`𝒥_d`$ terms, one obtains a frustrated $`S=3/2`$ chain which has been investigated with DMRG and leads to a gap for $`𝒥_b/𝒥_a0.3`$ . It seems to be possible to obtain antiferromagnetic $`𝒥_a`$ and $`𝒥_b`$ in this region if $`J_2`$ and $`J_3`$ are chosen suitably and large (a region including the coupling constants determined in section II B). However, then one is not in the perturbative region anymore and the $`𝒥_c`$ and $`𝒥_d`$ terms may also become important. Further discussion is therefore needed for reliable conclusions about a gap on the basis of the Hamiltonian (B4) with coupling constants (B8).
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# A scalar model of inhomogeneous elastic and granular media
## I Introduction
Stress transmission in dry granular media is unusual because in these materials no simple relation exists between stress and strain. Physical ingredients that give rise to this are that there are no tensile forces, that the particle deformations are very small, and that the particles can rearrange . Over the last several years evidence has accumulated that force propagation in dry granular media could be fundamentally different than in elastic solids . Equations that have been proposed to describe stresses in lightly loaded granular media have the property that specification of boundary conditions at the top surface of the system is sufficient to determine the stresses throughout, in marked contrast to the elliptic equations of elasticity theory.
However, applying a large enough uniform pressure to a granular material will cause it to exhibit an elastic linear response to a small additional stress. This is because uniform pressure both inhibits rearrangements (because it suppresses Reynolds dilatancy) and compresses the contacts, so that the non-tensile constraint on the interparticle forces becomes irrelevant. Thus, if stress propagation in lightly loaded granular media is indeed substantially different than in elastic media, then subjecting the material to high pressures will change fundamentally the stress propagation characteristics.
This paper investigates theoretically the stress propagation in granular materials as they are subjected to increasing pressures. The goals of this work are to understand the physical mechanisms governing the evolution between granular and elastic behavior and to make specific experimental predictions for the behavior of granular media under increasing loads.
We study a two-dimensional model system and compare the results to molecular dynamics (MD) simulations of two-dimensional systems of slightly polydisperse discs. Numerical studies of statistical models of granular media, where geometrical complexity is modeled in terms of uncorrelated random variables, are much faster and simpler than molecular dynamics simulations. Models of this type hold promise as a means to obtaining insight into the physics underlying the force propagation in granular materials. Our model for the granular regime is the two-dimensional scalar $`q`$-model . Though the $`q`$-model has deficiencies , it is attractive because of its simplicity and its prediction of an exponential tail in the probability distribution of stress within a packing agrees with experiments and with simulations . Our model for the elastic regime is a network of springs with a regular topology, with disorder introduced via randomly chosen spring constants. To model the crossover between the two regimes, we exploit our observation that the $`q`$-model can be written as a scalar elastic network subject to certain constraints. Enforcing these constraints to an increasing degree, which causes the force propagation behavior to evolve from that of an elastic system to that of the $`q`$-model, models the crossover between elastic and granular behavior by a particulate assemblage subjected to decreasing pressure.
We test the lattice model by comparing the results from the model to those of our MD simulations of two-dimensional systems of slightly polydisperse discs, focusing primarily on the probability distribution of stresses and on the two-point stress-stress correlation functions. The results of this investigation are mixed. The crossover in the force histogram between the elastic network and the $`q`$-model is strikingly similar to the crossover observed in the molecular dynamics simulations as the pressure on the system is decreased. However, the lattice model and the molecular dynamics simulation exhibit qualitatively different trends in the behavior of the two-point correlation functions of the stress.
The paper is organized as follows: Section II defines the scalar networks that we investigate. Section III details the process of generation, solution, and analysis of these networks and discusses the generation of the molecular dynamics simulations of slightly polydisperse discs. Section IV reports the results of the force distributions and spatial correlation functions for both the scalar lattice model and the MD simulations. Section V compares the results of the scalar lattice model, the MD simulations, and relevant experiments. Section VI summarizes and interprets our results. Appendix A calculates a finite-size correction to the in-plane stress-stress spatial correlation function for the $`q`$-model that is relevant to the interpretation of our numerical results.
## II Scalar elastic networks and the $`q`$-model
This section discusses the relationship between the $`q`$-model and the elastic network studied in this paper. Both models are scalar and are defined on a two-dimensional lattice. A scalar model is appropriate for a spring network if either the the network is very highly stretched , or if the motions are constrained so that displacements are unidirectional . We consider the second situation and denote the direction along which the motion occurs as $`\widehat{y}`$, with positive $`y`$ pointing downwards.
Consider a network of nodes connected by springs on a diamond lattice as shown in Fig. 1, where the motion of every node is constrained to be along the vertical direction $`\widehat{y}`$. Each spring has the same unstretched length, so that in the limit of zero load the system forms a regular lattice. The springs connecting the nodes have spring constants that are chosen independently from a fixed probability distribution. Periodic boundary conditions are imposed in the horizontal direction, and the locations of the nodes at the top and bottom boundaries are fixed so that the vertical displacement of all the nodes in these rows relative to the unloaded configuration are identical. We index the nodes so that a node in column $`j`$ in a row $`i`$ with odd (even) $`i`$ lies along the same vertical line as the other nodes in column $`j`$ in rows with odd (even) indices.
Let $`y_{i,j}`$ be the position of the node in row $`i`$ and column $`j`$ measured relative to its location in the absence of a load, and let $`k_{i,j}^l`$ and $`k_{i,j}^r`$ be the spring constants of the springs emanating downward from the node at row $`i`$ and column $`j`$. Every spring obeys Hooke’s law, so that $`f_{i,j}^l`$ and $`f_{i,j}^r`$, the forces exerted on node $`(i,j)`$ by the left and right springs below the node, are $`k_{i,j}^l(y_{i,j}y_{i+1,j1})`$ and $`k_{i,j}^r(y_{i,j}y_{i+1,j})`$ for odd $`i`$ ($`k_{i,j}^l(y_{i,j}y_{i+1,j})`$ and $`k_{i,j}^r(y_{i,j}y_{i+1,j+1})`$ for even $`i`$). The system is then compressed by setting $`y_{1,j}=(L_y1)\mathrm{\Delta }Y`$ and $`y_{L_y,j}=0`$ for all $`j`$, where rows $`1`$ and $`L_y`$ are the top and bottom rows, respectively, and $`\mathrm{\Delta }Y`$ is the average strain. We define $`F_{i,j}`$ to be the total vertical force incident from above on node $`(i,j)`$, so that $`F_{i,j}=f_{i1,j1}^r+f_{i1,j}^l`$. The forces and displacements are determined by balancing the forces at every node, $`F_{i,j}=f_{i,j}^l+f_{i,j}^r`$, and requiring that each $`y_{i,j}`$ be well-defined. This latter condition can be written as $`𝐒=[d]𝐘`$; here, $`𝐒`$ is the strain and $`𝐘`$ is the displacement field .
In our spring networks, each spring constant has a value selected independently at random from various probability distributions that are described below. We obtain the forces and strains along each link of each network using the method outlined in Ref. .
This scalar elastic model is equivalent to a resistor network . Forces and strains in the elastic system correspond to currents and voltages, respectively, in the resistor network. The requirement that the vertical forces at each node balance is equivalent to Kirchhoff’s current law, while the requirement that the position of each node is well-defined is equivalent to Kirchhoff’s voltage law.
### A Comparison between the elastic model and the $`q`$-model.
The force $`F_{i,j}`$ incident from above on node $`(i,j)`$ is transmitted to the sites below in the two pieces $`f_{i,j}^l`$ and $`f_{i,j}^r`$. Because of force balance, one can always write
$$f_{i,j}^l=q_{i,j}F_{i,j},f_{i,j}^r=(1q_{i,j})F_{i,j}.$$
(1)
In a $`q`$-model, the $`q_{i,j}`$ are random variables that are chosen independently at every site. In an elastic network, Eq. (1) still holds, but the $`q_{i,j}`$ are determined by the configuration of random spring constants together with the requirement that the displacement field be single-valued. For spring constants that are chosen independently, the force along any branch will depend on the values of the spring constants throughout the system. Important consequences of this non-locality include the presence of spatial correlations between the $`q_{i,j}`$’s and a nontrivial relation between the distribution of spring constants and the distribution of the $`q`$’s, including possibly the presence of $`q`$’s that are negative, indicating the appearance of tensile forces in the network.
A key observation underlying our work is that the $`q`$-model is equivalent to an elastic network subject to the constraint that the strain on every spring in each row is identical. The strain need not be constant from one row to the next, but it is simplest to consider the case in which it is. Let the amount of strain be $`\mathrm{\Delta }Y`$. Given the total force incident on node $`(i,j)`$ from above, $`F_{i,j}`$, if one chooses the spring constants to be
$$k_{i,j}^l=\frac{q_{i,j}F_{i,j}}{\mathrm{\Delta }Y},k_{i,j}^r=\frac{(1q_{i,j})F_{i,j}}{\mathrm{\Delta }Y},$$
(2)
then the force exerted down the left link emanating from node $`(i,j)`$ is $`k_{i,j}^l\mathrm{\Delta }Y=q_{i,j}F_{i,j}`$, and the force exerted down the right link from node $`i,j`$ is $`k_{i,j}^r\mathrm{\Delta }Y=(1q_{i,j})F_{i,j}`$. This force redistribution rule is exactly that of of the $`q`$-model. Given the set of $`q_{i,j}`$ values and the forces at each node in the top row of the system, we can create an equivalent spring network in a layer-by-layer manner.
We do not implement explicitly a no-tensile force constraint in our networks, in contrast to the work of Refs. and . However, in the $`q`$-model limit, there are no tensile forces. Our molecular dynamics simulations of lightly loaded material yield force distributions much closer to that of the $`q`$-model than to those of the non-tensile elastic networks of Ref. .
To study the crossover between elastic and $`q`$-model behavior, we generate iteratively a sequence of networks that interpolate between the elastic and $`q`$-model limits. The procedure adjusts the spring constants to make the strain in the system more uniform while keeping the ratio of spring constants emanating from each node constant. At iteration $`n`$, the spring constants $`k_{i,j}^l(n)`$ and $`k_{i,j}^r(n)`$ are set to
$`k_{i,j}^l(n)`$ $`=`$ $`{\displaystyle \frac{F_{i,j}(n1)}{\mathrm{\Delta }Y}}q_{i,j},`$ (4)
$`k_{i,j}^r(n)`$ $`=`$ $`{\displaystyle \frac{F_{i,j}(n1)}{\mathrm{\Delta }Y}}(1q_{i,j}),`$ (5)
where $`F_{i,j}(n1)`$ is the force through node $`(i,j)`$ at iteration $`n1`$. The $`q_{i,j}`$ are kept fixed, and the iteration procedure is started with $`F_{i,j}(0)=1`$.
To characterize the crossover between elastic and $`q`$-model behavior as the iteration proceeds, we need to quantify the degree to which the constant-strain constraint is violated. We use as our measure of the spatial variation in the strain the dimensionless quantity
$$\delta 𝒮_N\frac{{\displaystyle \frac{1}{L_x(L_y1)}}{\displaystyle \underset{i=1}{\overset{L_y1}{}}}{\displaystyle \underset{j=1}{\overset{L_x}{}}}{\displaystyle \frac{1}{2}}\left((\delta Y_{i,j}^l\overline{\delta Y})^2+(\delta Y_{i,j}^r\overline{\delta Y})^2\right)}{\overline{\delta Y}^2}$$
(6)
where $`\delta Y_{i,j}^l=Y_{i,j}Y_{i+1,j1}`$ and $`\delta Y_{i,j}^r=Y_{i,j}Y_{i+1,j}`$ for odd $`i`$ ($`\delta Y_{i,j}^l=Y_{i,j}Y_{i+1,j}`$ and $`\delta Y_{i,j}^r=Y_{i,j}Y_{i+1,j+1}`$ for even $`i`$), and
$$\overline{\delta Y}=\frac{1}{L_x(L_y1)}\underset{i=1}{\overset{L_y1}{}}\underset{j=1}{\overset{L_x}{}}\frac{1}{2}(\delta Y_{i,j}^l+\delta Y_{i,j}^r)=\mathrm{\Delta }Y.$$
(7)
Here, $`L_y`$ and $`L_x`$ are the number of rows and columns, respectively. In the elastic limit $`\delta 𝒮_N0.2`$, and as discussed above, $`\delta 𝒮_N`$ is zero for the $`q`$-model.
## III Methods
### A Scalar lattice model
We consider diamond-shaped lattices with springs on each link, as shown in Fig. 1. The positions of the top and bottom node layers are fixed and periodic boundary conditions are imposed in the transverse direction. The forces along all the links depend on the choice of spring constants, $`\{k_{i,j}\}`$, and are calculated using the node-potential method described in Ref. . The overall strain for each network is scaled so that the average vertical force through each node is normalized to unity,
$$\overline{F}\frac{1}{L_xL_y}\underset{i=1}{\overset{L_y}{}}\underset{j=1}{\overset{L_x}{}}F_{i,j}=1,$$
(8)
where the sum is over the nodes in the network.
Networks of height $`L_y=500`$ are used, with analysis performed on separate groups of layers to distinguish between edge and bulk effects. The widths $`L_x=16`$ and $`128`$ are powers of 2 in order to take advantage of FFT techniques in the calculation of spatial correlation function values described below. The number of realizations averaged over varies from 10 to 50, depending on lattice size and number of iterations.
For the elastic regime, we use four different distributions of spring constants: uniform distribution of $`k^1`$ for $`k^1(0,1)`$, gaussian distribution of $`k^1`$ with the configuration average $`(\overline{k^1})=1`$ and standard deviation $`\sigma _{k^1}=0.5`$, uniform distribution of $`k`$ with $`k(0,1)`$, and gaussian distribution of $`k`$ with $`\overline{k}=1`$ and $`\sigma _k=0.5`$. We construct networks with $`L_x=16`$ and $`128`$ with 50 and 25 realizations, respectively.
For the $`q`$-model regime, a uniform distribution of $`q`$ with $`q(0,1)`$ is used. We implement the iterative scheme with networks of size $`L_x=16`$ and $`128`$ with 50 and 10 realizations, respectively, for 100 iterations.
The local stress redistribution in a real granular material depends on microscopic details such as particle shape, friction characteristics, and preparation history. Instead of attempting to model the local force redistribution rules microscopically, our statistical models treat them as random variables chosen from different probability distributions. Since these probability distributions are not known a priori, we wish to identify and study properties that are not sensitive to the choice of the probability distribution governing the local force redistribution in the model. We focus on $`P(F)`$, the probability distribution of stresses at the nodes; $`\stackrel{~}{P}(q)`$, the probability distribution of the redistribution fractions $`q`$; and the spatial correlation functions of the force fluctuations about the mean values ,
$$C_k(j)=\frac{1}{L_yk}\underset{l=1}{\overset{L_yk}{}}\left(\frac{{\displaystyle \underset{m=1}{\overset{L_x}{}}}\delta F_{l,m}\delta F_{l+k,m+j}}{{\displaystyle \underset{m=1}{\overset{L_x}{}}}\delta F_{l,m}^2}\right),\stackrel{~}{C}_k(j)=\frac{1}{L_yk}\underset{l=1}{\overset{L_yk}{}}\left(\frac{{\displaystyle \underset{m=1}{\overset{L_x}{}}}\delta q_{l,m}\delta q_{l+k,m+j}}{{\displaystyle \underset{m=1}{\overset{L_x}{}}}\delta q_{l,m}^2}\right),$$
(9)
where $`\delta F_{i,j}=F_{i,j}\overline{F}`$ and $`\delta q_{i,j}=q_{i,j}\overline{q}`$; $`\overline{F}`$ is the average force and $`\overline{q}`$ is the average $`q`$ value. The indices $`l`$ and $`m`$ in Eq. 9 label layers and columns, respectively, while $`k`$ and $`j`$ are the spatial separation in layers and columns. These correlation functions are normalized so that $`C_0(0)=1`$ and $`\stackrel{~}{C}_0(0)=1`$. Positive values (correlation) indicate a tendency for nodes separated by $`k`$ rows vertically and $`j`$ columns horizontally to be either both above or both below the mean, while negative values (anti-correlation) indicates opposite behavior of one above and one below the mean.
### B Molecular dynamics simulations
Here we discuss our molecular-dynamics (MD) simulations to used to generate 2-D packings of discs. Varying the ratio of external load pressure to particle stiffness induces crossover between granular and elastic behavior. We calculate the probability distributions and corresponding spatial correlation function values for forces and redistribution fractions $`q`$ that are analogous to those in the scalar model.
Our simulations employ a method similar to that used by Durian and collaborators for sheared foams , in addition incorporating kinetic friction, contact damping, and particle rotation, and using two different repulsive interparticle force laws (linear and Hertzian).
#### 1 MD interaction rules
The discs in our simulation are all of identical mass $`m_D=1`$ and interact via purely repulsive normal contact forces and kinetic friction. The interaction force between two discs whose centers are at positions $`\stackrel{}{r}`$ and $`\stackrel{}{r}_j`$ with radii $`a_i`$ and $`a_j`$ is non-zero only if their separation $`\delta r_{i,j}0`$, where
$$\delta r_{i,j}=\left|\stackrel{}{r}_i\stackrel{}{r}_j\right|(a_i+a_j).$$
(10)
The normal contact force $`_{i,j}`$ is calculated from the overlap $`|\delta r_{i,j}|`$. We examine two force laws. The first is a linear force law based on a spring-like restoring force that yields
$$_{i,j}=K_{i,j}\left|\delta r_{i,j}\right|,$$
(11)
with $`K_{i,j}=(1/K_i+1/K_j)^1`$, where $`K_d`$ is the spring constant for disc $`d`$. The second is a non-linear force law based on Hertzian contacts between spheres,
$$_{i,j}^{[\mathrm{HC}]}=D^1(\frac{1}{a_i}+\frac{1}{a_j})^{\frac{1}{2}}\left|\delta r_{ij}\right|^{\frac{3}{2}},$$
(12)
where $`D=\frac{3}{2}\left((1\sigma ^2)/E\right)`$, with $`\sigma `$ and $`E`$ being the material properties Poisson’s ratio and Young’s modulus, respectively. For both force laws, the forces are directed so as to separate the overlapping discs. To calculate forces generated by interactions with walls, we assume the walls to be discs of infinite radius.
Kinetic friction is incorporated into the disc interactions although static friction is not. The introduction of frictional forces causes the discs to rotate; however, the frictional force is zero at mechanical equilibrium. The kinetic friction force $`f_{i,j}`$ for contact between discs $`i`$ and $`j`$ is
$$f_{i,j}=\mu _{i,j},$$
(13)
where $`\mu `$ is the coefficient of kinetic friction, and is directed opposite to the contact point velocity $`\stackrel{}{v}_{i,j}^{\mathrm{cp}}`$. For disc $`i`$, this velocity $`\stackrel{}{v}_{i,j}^{\mathrm{cp}}`$ is related to disc velocities $`\stackrel{}{v}_i`$ and $`\stackrel{}{v}_j`$, the angular velocities $`\omega _i`$ and $`\omega _j`$, and directional vector $`\widehat{r}=(\stackrel{}{r}_i\stackrel{}{r}_j)/|\stackrel{}{r}_i\stackrel{}{r}_j|`$ by
$$\stackrel{}{v}_{i,j}^{\mathrm{cp}}=\stackrel{}{v}_{i,j}\left(\stackrel{}{v}_{i,j}\widehat{r}\right)\widehat{r}+\left(a_i\stackrel{}{\omega }_i+a_j\stackrel{}{\omega }_j\right)\times \widehat{r},$$
(14)
where $`\stackrel{}{v}_{i,j}=\stackrel{}{v}_i\stackrel{}{v}_j`$.
Damping during contact between discs $`i`$ and $`j`$ is used as an additional means of dissipating kinetic energy. It is generated by applying to disc $`i`$ a force $`_\mathrm{D}`$ and torque $`\mathrm{\Gamma }_\mathrm{D}`$ given by
$`_\mathrm{D}`$ $`=`$ $`\lambda _{\mathrm{trans}}v_i^{},`$ (16)
$`\mathrm{\Gamma }_\mathrm{D}`$ $`=`$ $`\lambda _{\mathrm{ang}}\omega _i^{},`$ (17)
where $`v_i^{}`$ is the translational velocity of disc $`i`$ relative to the interaction center of mass for the two discs $`i`$ and $`j`$ that are in contact and $`\omega _i^{}`$ is its angular velocity relative to the total angular momentum of the disc pair. $`\lambda _{\mathrm{trans}}`$ and $`\lambda _{\mathrm{ang}}`$ are damping constants. This process conserves both translational and angular momentum. Energy is directly removed from the system as opposed to being converted between translational and rotational motion.
The bottom and top walls have mass $`m_W`$ and are constrained to move only vertically. An inward force of magnitude $`F_{\mathrm{wall}}`$ is applied to each wall in order to compress the system. Damping of the wall motion suppresses volume oscillations and serves as the primary means of energy dissipation. The damping force $`_{\mathrm{WD}}`$ on a wall is
$$_{\mathrm{WD}}=\lambda _Wv_W,$$
(18)
where $`v_W`$ is the velocity of the wall and $`\lambda _W`$ is the wall damping constant.
#### 2 MD implementation
Ensembles of systems of $`N=1024`$ discs of average radius $`a_D`$ are generated by starting with triangular array of $`\sqrt{N}`$ rows and $`\sqrt{N}`$ discs per row placed in a horizontally periodic system with both height and width $`L=2.273a_D\sqrt{N}`$. For the data shown here, discs are placed in the system at positions $`(L(n_x+0.05)/\sqrt{N},L(n_y+0.5)/\sqrt{N})`$ for odd $`n_y`$ and $`(L(n_x+0.55)/\sqrt{N},L(n_y+0.5)/\sqrt{N})`$ for even $`n_y`$, with indices $`n_x`$ and $`n_y`$ running from 0 to $`\sqrt{N}1`$. In practice, discs with gaussian distributed polydispersity of $`\sigma _a=0.1a_D`$ placed on this triangular array do not overlap. The results obtained are not sensitive to initial disc placement. The system is then compressed by the application of an inward force on the top and bottom walls. All discs have the same spring constant $`K_dK=1`$. The coefficient describing wall damping is set to $`\lambda _W/m_W=1`$. Damping coefficients for translational and angular motion for disc contacts are set to $`\lambda _{\mathrm{trans}}/m_D=1`$ and $`\lambda _{\mathrm{ang}}/I_D=4.1`$, where $`I_D=\frac{1}{2}m_Da_D^2`$ is the moment of inertia for a disc with radius $`a_D`$. The coefficient of kinetic friction is set to $`\mu =0.2`$ for both disc-disc and disc-wall contacts. Comparisons with samples produced without disc-contact damping or kinetic friction revealed no measurable differences in force probability distributions or in the two-point force correlation function. Incorporating additional energy-dissipation mechanisms allows systems to reach mechanical equilibrium more rapidly. The end time for each compression stage is chosen so that the average residual kinetic energy for each disc is equivalent to translational movement of approximately or less than $`0.01a_D`$ in unit time. Because of the increased external energy input, systems at higher compressions are allowed a less restrictive limit of approximately $`0.05a_D`$. Visual inspection of final configurations do not reveal significant fluctuations in time in contact network topology or force magnitude in load-bearing structures. Comparisons with test systems with longer run times also do not show any significant quantitative differences.
For a system of fixed size $`L`$, the typical compression of the system can be controlled through variations in the disc spring constant $`K`$ or applied external force $`F_{\mathrm{wall}}`$. Typical relative particle deformations $`\delta `$ is given by
$$\delta \frac{1}{N_C}\underset{(i,j)}{}\frac{|\delta r_{i,j}|}{(a_i+a_j)}=\frac{1}{N_C}\underset{(i,j)}{}\frac{_{i,j}}{K_{i,j}}(a_i+a_j)^1\frac{F_{\mathrm{wall}}\left(\frac{L}{a_D}\right)^1}{\frac{K}{2}}(2a_D)^1=\frac{F_{\mathrm{wall}}}{LK}=\frac{\mathrm{\Pi }}{K},$$
(19)
where $`N_C`$ is the total number of contacts, the sums are over pairs of discs $`i`$ and $`j`$ in contact, and $`\mathrm{\Pi }F_{\mathrm{wall}}/L`$ is the external pressure. This estimate is approximate due to geometric factors and distributional fluctuations; however, the scaling of deformations to $`\mathrm{\Pi }/K`$ should hold generally. In our simulations, the disc spring constant $`K`$ is held fixed and the pressure $`\mathrm{\Pi }`$ is varied to induce crossover between granular and elastic behavior. We define the reference pressure $`\mathrm{\Pi }=\mathrm{\Pi }_0`$ such that the relative particle deformation $`\delta 6.25\times 10^4`$. The reference compression pressure $`\mathrm{\Pi }_0`$ yields a force histogram typical of the granular range, as discussed below in Sec. IV. After the initial compression with $`\mathrm{\Pi }=\mathrm{\Pi }_0`$, the applied pressure is increased in stages to $`\mathrm{\Pi }=100\mathrm{\Pi }_0`$, at which $`\delta 0.01`$. We also decrease the pressure from the initial $`\mathrm{\Pi }=\mathrm{\Pi }_0`$ configuration down to $`\mathrm{\Pi }=0.01\mathrm{\Pi }_0`$ ($`\delta 10^6`$) in order to approach the zero-deformation limit. Fig. 2 shows a sample MD system subjected to the pressures $`0.1\mathrm{\Pi }_0`$, $`\mathrm{\Pi }_0`$, $`10\mathrm{\Pi }_0`$, and $`50\mathrm{\Pi }_0`$.
For spheres with Hertzian contacts (using Eq. 12), the deformation can be approximated by $`\delta ^{[\mathrm{HC}]}\left(\frac{\mathrm{\Pi }D}{2a_D}\right)^{\frac{2}{3}}`$. For our simulations $`D`$ is chosen to yield deformations of the same order of magnitude as the linear contacts at the compression $`\mathrm{\Pi }=\mathrm{\Pi }_0`$. The pressures studied are the same as for the linear spring contact systems.
## IV Results
### A Scalar lattice model
Here we present our results for the scalar lattice models. We study how the probability distribution of total vertical force $`F`$ incident on a node from above $`P(F)`$ and the two-point force correlation function $`C_k(j)`$ characterize the behavior in scalar elastic lattice networks in which the constant-strain constraint is enforced to varying degrees. In the $`q`$-model both $`P(F)`$ and in-plane force-force correlation function $`C_0(j)`$ exhibit robust behavior for generic choices of probability distributions of $`q`$’s. We investigate the degree to which these quantities depend on the choice of spring constant distributions in the elastic networks, and discuss the crossover of $`P(F)`$ and $`C_0(j)`$ between the elastic and $`q`$-model behavior as the constant-strain constraint is implemented with increasing accuracy.
#### 1 Results for the $`q`$-model
In the $`q`$-model, the force histogram $`P(F)`$ decays exponentially at large forces and $`C_0(j)`$ is zero for non-zero $`j`$ . These properties hold for a wide variety of choices of the distribution of $`q`$-values.
Our results for the crossover from elastic to $`q`$-model behavior are obtained for the specific choice that the $`q`$’s are uniformly distributed in $`[0,1]`$. A two-dimensional $`q`$-model with this distribution of $`q`$’s yields
$$P(F)=4Fe^{2F}.$$
(20)
For a system of infinite lateral extent, the in-plane force-force correlation function $`C_0(j)=\delta _{j0}`$ where $`\delta _{j0}`$ is the Kronecker delta function . For a system of finite width $`L_x`$, force correlations must arise because all forces are positive and the total force through a layer is fixed. As discussed in Appendix A, assuming that this mechanism is the only one giving rise to correlations, one obtains that a 2-D system of lateral extent $`L_x`$ has $`C_0(j)`$ given by
$$C_0(j0)=(L_x1)^1.$$
(21)
This form for $`C_0(j)`$ agrees with our numerical results for the $`q`$-model on lattices of finite width.
#### 2 Elastic networks
For elastic networks with different distributions of spring constants, the probability distribution of vertical force $`P(F)`$, shown in Figs. 3, is narrower than that of the $`q`$-model. Its functional form depends on the choice of spring constant distribution. Choosing the spring constants $`k`$ from a distribution either uniform in $`k`$ or gaussian in $`k`$ yields $`P(F)`$’s that are roughly gaussian while the $`P(F)`$’s for networks for distributions uniformly distributed or gaussian in $`k^1`$ display a tail at large $`F`$ that is consistent with an exponential decay. Networks with gaussian distributed $`k`$ or $`k^1`$ exhibit narrower $`P(F)`$’s than their counterparts with uniformly distributed $`k`$ or $`k^1`$.
In contrast to the behavior of the force probability distribution $`P(F)`$, the force-force correlation function values $`C_k(j)`$ are quantitatively indistinguishable for all the distributions of spring constants that we examined, as shown in Fig. 4. For $`C_0(j)`$, the force-force correlation function within the same layer, we see a strong anti-correlation for $`j=1`$ of magnitude $`0.30`$ that decays within $`j8`$. For vertical separation $`k>0`$, we see a simultaneous reduction in peak magnitude (at $`j=0`$) and broadening of peak width but with the anti-correlation signature and decay joining the curve laid out by $`k=0`$.
The probability distributions of redistribution fraction $`q`$, $`\stackrel{~}{P}(q)`$, shown in Fig. 5, are roughly gaussian and peaked at $`q=\overline{q}0.5`$ for all distributions of spring constants examined. The widths of the $`\stackrel{~}{P}(q)`$ depend on the choice of distribution of spring constants, with the gaussian distributed $`k`$ and $`k^1`$ once again being narrower (standard deviation $`\sigma _q0.16`$ and $`0.15`$, respectively) than their uniformly distributed counterparts ($`\sigma _q0.25`$ for random $`k`$ and $`\sigma _q0.21`$ for random $`k^1`$). All of the elastic networks display significant correlations between $`q`$’s at different nodes as demonstrated in Fig. 6, which shows the correlation function $`\stackrel{~}{C}_k(j)`$ for all the random distributions. The correlations between $`q`$’s are an important factor in determining the statistical distribution of the forces in these systems; Fig. 7 shows that a $`q`$-model system with the same $`\stackrel{~}{P}(q)`$ as an elastic network with a uniform distribution of $`k`$ but with no correlations between the $`q`$’s yields a $`P(F)`$ markedly different from the elastic network.
No differences between the bulk (layers 201-300) and edge (layers 1-100 and 401-500) sections are detected in the distributions $`P(F)`$ and $`\stackrel{~}{P}(q)`$ or the correlation functions $`C_k(j)`$ and $`\stackrel{~}{C}_k(j)`$. The results for lattices with $`L_x=128`$ are the same within statistical errors to the results from $`L_x=16`$ lattices.
In the elastic networks, forces in less than 1% of the branches are tensile, and no node in any of the networks is subject to a tensile net force. Our results for $`P(F)`$ for uniformly distributed $`k`$’s are very similar to those reported in Sexton et al. , where a non-tensile force constraint is enforced.
#### 3 Iterated networks– the $`q`$-model limit
We now discuss the networks generated by our iterative algorithm for converting an elastic network to a $`q`$-model system. First, we verify that the generated networks do eventually converge to the $`q`$-model. After 100 iterations, the forces along the links of the iterated spring network are identical to those of the corresponding $`q`$-model to within $`10^4`$.
A subtle point in the method is that our iterative scheme yields a configuration in which the forces at the top and bottom boundaries of the iterated network may have nonzero spatial correlations, as the initial iteration $`n=0`$ system is elastic. As one proceeds away from the top and bottom boundaries, these correlations decay via a diffusive process that takes on the order of $`L_x^2`$ layers . Thus, forces at different sites in the same layer are effectively uncorrelated only for systems with large aspect ratios. This result is consistent with our numerical observation that in fully iterated systems, correlations between forces at different sites in the same layer are present throughout the $`L_x=128`$ systems, while they are only present in the top- and bottom-most 200 layers of $`L_x=16`$ systems.
#### 4 Iterated networks– crossover between elastic and $`q`$-model behavior
In the iterated networks the target values of the $`q_{i,j}`$ are fixed at the outset of iteration procedure. The initial state (zeroth iteration) is an elastic network with spring constants given by $`k_{i,j}^l=q_{i,j}/\mathrm{\Delta }Y`$ and $`k_{i,j}^r=(1q_{i,j})/\mathrm{\Delta }Y`$. Therefore, the initial probability of node forces $`P(F)`$ and the spatial correlation function $`C_0(j)`$ are those of elastic networks with spring constants chosen from a uniform distribution of $`k`$. The realized $`q`$ distribution $`\stackrel{~}{P}(q)`$ (as opposed to the distribution of the target $`q`$ values) is peaked at $`q=0.5`$ and its spatial correlation function $`\stackrel{~}{C}_k(j)`$ reveals slight nearest-neighbor correlations for $`k=0`$ and anti-correlations at $`j=0`$ for $`k>0`$, once again matching elastic-regime behavior.
The probability distribution of node forces, $`P(F)`$, is shown in Fig. 8(a) for different values of iteration number $`n`$. As the number of iterations is increased, the $`P(F)`$ develops an exponential tail at large forces. Fig. 8(b) shows the probability distribution of the $`q`$’s, $`\stackrel{~}{P}(q)`$, versus the number of iterations. $`\stackrel{~}{P}(q)`$ approaches the target form of a uniform distribution after roughly 10 iterations.
Fig. 9(a) shows our results for nearest-neighbor in-plane and vertical force-force correlation function values $`C_0(1)`$ and $`C_2(0)`$ as the number of iterations $`n`$ is increased. Fig. 9(b) shows the corresponding force-fraction $`q`$-$`q`$ correlations $`\stackrel{~}{C}_0(1)`$ and $`\stackrel{~}{C}_2(0)`$. While only about 10 iterations are necessary before the nearest-neighbor spatial correlations between $`q`$ values go to zero, in-plane force-force correlations are still present after 100 iterations although much reduced in magnitude from the initial elastic (iteration $`n=0`$) value and approaching asymptotically the expected zero-correlation value.
The quantity $`\delta 𝒮_N`$ that we use to characterize the crossover between elastic and $`q`$-model behavior is defined in Sec. II. As Fig. 10 demonstrates, we observe $`\delta 𝒮_N`$ decreasing as the number of iterations $`n`$ is increased according to a power law. A fit that assumes the dependence is of the form $`\delta 𝒮_Nn^\alpha `$ yields $`\alpha =1.68\pm 0.02`$.
### B Results of molecular dynamics simulations
Here we discuss the results of our MD simulations. Fig.11 shows $`P(F)`$, the probability distribution of vertical forces $`F=\stackrel{}{}\widehat{y}`$, for MD systems under various applied pressures $`\mathrm{\Pi }`$. As with the scalar model, $`F`$ has been normalized so that the average vertical force $`\overline{F}=1`$ for each system configuration. The progression of $`P(F)`$ as pressure is increased is very similar both qualitatively and quantitatively to the crossover from granular to elastic behavior in the scalar model lattice systems. We calculate the force-force correlation values $`C_0(j)`$, shown in Fig. 11(b), by defining discs to be in plane with a tolerance of $`\pm 0.10a_D`$ and $`j`$ in units of average disc diameter $`2a_D`$. In contrast with the scalar model behavior, the MD systems exhibit a significant nearest-neighbor anti-correlation for all applied external pressures. These results for $`P(F)`$ and $`C_0(j)`$ are independent of whether the samples are compressed in stages or directly at a fixed pressure $`\mathrm{\Pi }`$.
We define the $`q`$ value of a disc as the fraction of total vertical force received from its topward neighbors that is transferred to its bottom leftward neighbors. The probability distribution of $`q`$ values, $`\stackrel{~}{P}(q)`$ is shown in Fig.12. We also calculate the $`q`$-$`q`$ correlation values $`\stackrel{~}{C}_0(j)`$ and $`\stackrel{~}{C}_k(0)`$, although the large errors prevent the extraction of quantitative trends. Narrowing the statistical errors would be computationally prohibitive.
The number of contacts increases significantly with the pressure, as shown in Fig. 13. As the magnitude of the typical overlap increases, additional contacts are formed. The number of contacts at low pressures is below the theoretically predicted average of $`Z=2d`$, where $`d`$ is the dimension of the system, because the polydispersity in radii and the lack of gravity allow for the existence of “rattlers” which do not support any of the external load.
Our results for the Hertzian contact systems are indistinguishable from those of the linear springs throughout most of the range of pressures explored. At higher pressures $`\mathrm{\Pi }`$ ($`\mathrm{\Pi }35\mathrm{\Pi }_0`$), the added stiffness of the Hertzian contacts leads to the slower narrowing of $`P(F)`$.
## V Comparison of results of MD simulations and of scalar elastic networks
Here we compare the behavior observed in the MD simulations and in the scalar elastic networks. Because different schemes are used to induce the granular-elastic crossover in the two systems (iterations in the scalar networks and external pressure for MD), we need to establish a common measure to quantify a system’s position within the crossover region. As the evolution of the probability distribution of vertical forces, $`P(F)`$, is qualitatively and quantitatively similar in the network model and in the MD simulations, we use matches in its form to establish a relationship between iteration number $`n`$ and applied pressure $`\mathrm{\Pi }`$. Fig. 14(a) shows matches in form between linear-force-law MD packings and iterated scalar network systems for $`\mathrm{\Pi }/\mathrm{\Pi }_0=100`$ and iteration $`n=0`$, $`\mathrm{\Pi }/\mathrm{\Pi }_0=10`$ and $`n=10`$, and $`\mathrm{\Pi }/\mathrm{\Pi }_0=1`$ and $`n=100`$ systems. From these matchings, we map the iteration number $`n`$ in the scalar networks to the equivalent applied pressure $`\mathrm{\Pi }_{\mathrm{eq}}(n)`$ in the MD using the simple scaling :
$$\frac{\mathrm{\Pi }_{\mathrm{eq}}(n)}{\mathrm{\Pi }_0}=\frac{100}{n}.$$
(22)
We perform a check on this proposed scaling by considering the analogous quantities of deviation from constant strain in scalar systems $`\delta 𝒮_N`$, given by Eq. 6, and deviation from the infinitely hard, zero-deformation limit in MD systems calculated by
$$\delta 𝒮_{MD}\frac{1}{N_C}\underset{(i,j)}{}\frac{|\delta r_{i,j}|^2}{(a_i+a_j)^2}$$
(23)
where $`N_C`$ is the total number of contacts and the sum is over pairs of discs $`i`$ and $`j`$ in contact. We match $`\delta 𝒮`$ values for $`\mathrm{\Pi }/\mathrm{\Pi }_0=10`$ and $`n=10`$ by scaling the square deviation for the scalar network systems by a constant factor of 0.030. Fig. 14(b) shows that this scaling yields reasonable agreement between $`\delta 𝒮_N`$ and $`\delta 𝒮_{MD}`$ over the crossover region.
In contrast to the agreement in the trends of $`P(F)`$, qualitative differences exist between the scalar network model and the MD simulations in spatial correlation function values $`C_j(k)`$. Fig. 15(a) shows the nearest-neighbor in-plane and vertical force-force correlation values, $`C_0(1)`$ and $`C_2(0)`$, for the crossover between elastic and granular regimes. While the MD systems exhibit a significant in-plane nearest-neighbor anti-correlation throughout the crossover, a decrease in its magnitude is seen in the scalar networks as the systems change from elastic to granular. MD systems do not exhibit strong vertical correlations, in contrast with the scalar networks whose $`C_2(0)`$ value increases significantly as the granular limit is approached.
The large statistical uncertainties in our $`q`$-$`q`$ correlation functions for MD systems restrict us to making only qualitative behavior descriptions. The trend for in-plane nearest-neighbor correlation behavior $`\stackrel{~}{C}_0(1)`$ in both systems is similar. However, qualitative differences exist for vertical correlation value $`\stackrel{~}{C}_2(0)`$: the MD systems display consistent anti-correlation behavior, while the scalar networks display anti-correlation behavior in the elastic regime which decays rapidly to uncorrelated behavior as the granular limit is approached.
Our work indicates that experiments on granular media at high pressures should yield a force histogram that differs qualitatively from that observed at lower pressures. Experiments by Howell et al. as well as experiments and simulations by Makse et al. are in qualitative agreement with this result. Howell et al. control the transition between granular and elastic behavior of slowly sheared systems in a 2-D Couette geometry by varying the packing fraction $`\gamma `$ within a range $`0.77\gamma 0.81`$. The average force/length on a particle increases with $`\gamma `$. For lower values of $`\gamma `$, the distribution of large stresses is asymptotically exponential, while the distribution of stresses has a gaussian form at higher packing fractions $`\gamma `$. Makse et al. apply increasing pressure to three-dimensional packings of spherical glass beads to achieve the crossover between granular and elastic behavior and also perform MD simulations on 3-D systems. Makse et al. observe a crossover in the force histogram $`P(F)`$ in a pressure range that is consistent with our 2-D MD results.
An interesting question is whether the persistent in-plane nearest-neighbor anti-correlation in the forces that is observed in the MD simulations is present in experimental systems. Mueth et al. do not find evidence of correlations between different sites in the same horizontal layer; any nearest-neighbor anti-correlation in the experiment is smaller than the experimental resolution. However, they measure a different correlation function, $`K_1(r)`$, defined as
$$K_1(r)=\frac{{\displaystyle \underset{i=1}{\overset{N_B}{}}}{\displaystyle \underset{j=i+1}{\overset{N_B}{}}}\delta (r_{ij}r)f_if_j}{{\displaystyle \underset{i=1}{\overset{N_B}{}}}{\displaystyle \underset{j=i+1}{\overset{N_B}{}}}\delta (r_{ij}r)},$$
(24)
where the sums are over the $`N_B`$ particles in the bottom layer, $`f_i`$ is the force at position $`r_i`$ in the bottom layer and $`r_{ij}=|\stackrel{}{r}_i\stackrel{}{r}_j|`$. Calculation of $`K_1(r)`$ from the numerical data for our MD simulations yields values of the correlation function that are smaller than the error bars in the experiment. Comparison with Ref. is necessarily qualitative since the experiments measure the properties at the surface of a 3-D packing while our MD results are calculated using numerical data from the bulk of a 2-D system.
## VI Discussion
We have investigated the crossover between elastic and granular stress transmission in both a 2-D scalar lattice model and in molecular dynamics simulations of slightly polydisperse discs. The evolution of $`P(F)`$, the probability distribution of stresses, is very similar in the lattice model and in the MD. However, the behavior of the spatial correlation functions for stress, $`C_k(j)`$, differs qualitatively.
Our investigations of the scalar model have several implications for the development of granular media models. First, we have shown that implementing a local constraint can convert an elastic network to a $`q`$-model. This constraint has the natural physical interpretation that the strain in the system must be uniform; it is plausible that rearrangements would prevent strain gradients from forming. Second, implementing this constraint to increasing accuracy causes the force histogram $`P(F)`$ to evolve in a manner similar to that observed in the MD simulation as the pressure is decreased. $`P(F)`$ has a tail consistent with exponential decay at large forces in the granular limit, while the $`P(F)`$ for the highly compressed system is much narrower and decays more quickly at large forces. We note that implementing a non-tensile force constraint alone, as in Ref. , yields gaussian decay in $`P(F)`$ at large forces even at the lowest pressures, in qualitative disagreement with the MD results of us and others .
While this success in describing the evolution of the force histogram and the scalar model’s simplicity in both formulation and implementation make it an attractive platform for the study of media models, the discrepancy in the behavior of the correlation function behavior with the MD simulation results needs to be addressed. The scalar model assumes explicitly that in the granular regime the stress redistribution fractions $`q`$ at different sites are uncorrelated. The extent to which this condition is valid needs to be examined in more detail. Spatial correlations of the $`q`$’s can strongly affect the probability distribution of stress $`P(F)`$ but the degree to which these correlations exist in real packings has not been settled. A possible source of spatial correlations in the $`q`$’s is the constraint that non-tensile vector forces must be balanced. However, vector generalization of $`q`$-model systems that have been proposed to date have required arbitrary constraints to be imposed to limit the scale of stress components perpendicular to the direction of applied force . Clarification of the roles of vector force balance and contact formation is key to identifying and characterizing the processes governing stress transmission beyond those that have been implemented in the scalar model.
In conclusion, we have shown that similarities exist in the evolution of the probability distribution of stresses $`P(F)`$ in the crossover between elastic and granular regimes for a scalar lattice model and MD simulations of slightly polydisperse discs. However, the systems exhibit qualitative differences in the two-point force correlation function $`C_k(j)`$. Further investigation of the systematic influences leading to the spatial correlations between forces is necessary for the development of a successful model of stress transmission in granular media.
## Acknowledgements
We thank Alexei Tkachenko and Tom Witten for a key suggestion and Rick Clelland, Heinrich Jaeger, Dan Mueth, Sid Nagel, Joshua Socolar, and Bob Behringer for useful conversations. This research was supported by the MRSEC program of the NSF and by the Petroleum Research Fund of the American Chemical Society.
## A Finite-size correction to correlation function calculation
In the $`q`$-model in the limit of infinite size, forces at different sites in the same layer are completely uncorrelated. In a system of finite transverse extent, the requirements that the total force through every layer is identical and there are no tensile forces lead to a finite-size correction to the correlation function. This appendix discusses this correction.
We characterize the correlations between force fluctuations on sites in the same row using the correlation function
$$C_0(j)=\frac{1}{L_y}\underset{l=1}{\overset{L_y}{}}\left(\frac{{\displaystyle \underset{m=1}{\overset{L_x}{}}}\delta F_{l,m}\delta F_{l,m+j}}{{\displaystyle \underset{m=1}{\overset{L_x}{}}}\delta F_{l,m}^2}\right),$$
(A1)
where $`\delta F_{l,m}=F_{l,m}\overline{F}`$ is the deviation of the force at a site in row $`l`$ and column $`m`$ from the average force $`\overline{F}`$. For the $`q`$-model with a uniform distribution of $`q`$’s, in a system of infinite transverse extent this correlation function is
$$C_0(j)=\{\begin{array}{cc}1\hfill & j=0\hfill \\ 0\hfill & j0\hfill \end{array}.$$
(A2)
This result follows from the fact that $`P(F_\alpha ,F_\beta )`$, the probability that the force through node $`\alpha `$ is $`F_\alpha `$ and the force through node $`\beta `$ at the same horizontal level is $`F_\beta `$, is factorizable:
$`P(F_\alpha ,F_\beta )=P(F_\alpha )P(F_\beta ),`$
As a result, for $`\alpha \beta `$
$`\delta F_\alpha \delta F_\beta `$ $`=`$ $`\underset{L_x,L_y\mathrm{}}{lim}{\displaystyle \frac{1}{L_yL_x}}{\displaystyle \underset{l=1}{\overset{L_y}{}}}{\displaystyle \underset{m=1}{\overset{L_x}{}}}\delta F_{l,m}\delta F_{l,m+j}`$ (A3)
$`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑F_\alpha 𝑑F_\beta \delta F_\alpha \delta F_\beta P(F_\alpha ,F_\beta )`$ (A4)
$`=`$ $`{\displaystyle _0^{\mathrm{}}}{\displaystyle _0^{\mathrm{}}}𝑑F_\alpha 𝑑F_\beta \delta F_\alpha \delta F_\beta P(F_\alpha )P(F_\beta )`$ (A5)
$`=`$ $`\left({\displaystyle _0^{\mathrm{}}}𝑑F_\alpha \delta F_\alpha P(F_\alpha )\right)\left({\displaystyle _0^{\mathrm{}}}𝑑F_\beta \delta F_\beta P(F_\beta )\right)`$ (A6)
$`=`$ $`\left[{\displaystyle _0^{\mathrm{}}}𝑑F\delta FP(F)\right]^2`$ (A7)
$`=`$ $`\left[{\displaystyle _0^{\mathrm{}}}𝑑F(F\overline{F})P(F)\right]^2`$ (A8)
$`=`$ $`\left[\overline{F}\overline{F}\right]^2`$ (A9)
$`=`$ $`0.`$ (A10)
On a lattice of finite width ($`L_x`$ sites), the multipoint force probability distribution function must be consistent with the facts that first, the total force down every layer is fixed, and second, no force is negative. This implies that
* The maximum force on any node in any layer cannot be larger than $`F_{\mathrm{max}}=L_x\overline{F}`$, and
* The force $`F_\alpha `$ at a node $`\alpha `$ contributes to the total force along a layer and hence affects the sum of the forces through the remaining sites in the layer. Defining $`\stackrel{~}{F}`$ as the average force through all the sites in the layer other than site $`\alpha `$, we have
$`\stackrel{~}{F}`$ $`=`$ $`{\displaystyle \frac{L_x\overline{F}F_\alpha }{L_x1}}`$ (A11)
$`=`$ $`\overline{F}{\displaystyle \frac{F_\alpha \overline{F}}{L_x1}}`$ (A12)
$`=`$ $`\overline{F}{\displaystyle \frac{\delta F_\alpha }{L_x1}}.`$ (A13)
Assuming that the only correlations present in the finite system are those required to satisfy these conditions, the joint probability distribution in the system with finite $`L_x`$ can again be written
$$P(F_\alpha ,F_\beta )=\frac{1}{\overline{F}\stackrel{~}{F}}P(\frac{F_\alpha }{\overline{F}})\stackrel{~}{P}(\frac{F_\beta }{\stackrel{~}{F}}),$$
(A14)
but now the distributions are subject to the constraints
$`{\displaystyle _0^{L_x}}P(\omega )𝑑\omega =1`$ $`,`$ $`{\displaystyle _0^{L_x}}\omega P(\omega )𝑑\omega =1`$ (A16)
$`{\displaystyle _0^{L_x1}}\stackrel{~}{P}(\nu )𝑑\nu =1`$ $`,`$ $`{\displaystyle _0^{L_x1}}\nu \stackrel{~}{P}(\nu )𝑑\nu =1.`$ (A17)
In the limit of $`L_x\mathrm{}`$, we expect corrections to $`P(F)`$ to be of order $`1/L_x`$ as the fluctuations in the forces at the sites are of order unity. As we will see, with the assumptions that we have made, the finite size correction to the correlation function does not depend on the form of the probability distribution $`P`$. Taking the new constraints into account, for $`\alpha \beta `$ we have
$`\delta F_\alpha \delta F_\beta `$ $`=`$ $`{\displaystyle _0^{L_x}}𝑑F_\alpha {\displaystyle _0^{L_xF_\alpha }}𝑑F_\beta \delta F_\alpha \delta F_\beta P(F_\alpha ,F_\beta )`$ (A18)
$`=`$ $`{\displaystyle _0^{L_x}}𝑑F_\alpha {\displaystyle _0^{L_xF_\alpha }}𝑑F_\beta (F_\alpha \overline{F})(F_\beta \overline{F})\left({\displaystyle \frac{1}{\overline{F}\stackrel{~}{F}}}P\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right)\stackrel{~}{P}\left({\displaystyle \frac{F_\beta }{\stackrel{~}{F}}}\right)\right)`$ (A19)
$`=`$ $`\overline{F}^2{\displaystyle _0^{L_x}}d\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right)\left({\displaystyle \frac{F_\alpha }{\overline{F}}}1\right)P\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right){\displaystyle \frac{\stackrel{~}{F}}{\overline{F}}}{\displaystyle _0^{L_xF_\alpha }}d\left({\displaystyle \frac{F_\beta }{\stackrel{~}{F}}}\right)\left({\displaystyle \frac{F_\beta }{\stackrel{~}{F}}}{\displaystyle \frac{\overline{F}}{\stackrel{~}{F}}}\right)\stackrel{~}{P}\left({\displaystyle \frac{F_\beta }{\stackrel{~}{F}}}\right)`$ (A20)
$`=`$ $`\overline{F}^2{\displaystyle _0^{L_x}}d\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right)\left({\displaystyle \frac{F_\alpha }{\overline{F}}}1\right)P\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right){\displaystyle \frac{\stackrel{~}{F}}{\overline{F}}}{\displaystyle _0^{L_x1}}𝑑\nu \left(\nu {\displaystyle \frac{\overline{F}}{\stackrel{~}{F}}}\right)\stackrel{~}{P}(\nu )`$ (A21)
$`=`$ $`\overline{F}^2{\displaystyle _0^{L_x}}d\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right)\left({\displaystyle \frac{F_\alpha }{\overline{F}}}1\right)P\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right)\left[{\displaystyle \frac{\stackrel{~}{F}}{\overline{F}}}{\displaystyle _0^{L_x1}}𝑑\nu \nu \stackrel{~}{P}(\nu ){\displaystyle _0^{L_x1}}𝑑\nu \stackrel{~}{P}(\nu )\right]`$ (A22)
$`=`$ $`\overline{F}^2{\displaystyle _0^{L_x}}d\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right)({\displaystyle \frac{F_\alpha }{\overline{F}}}1)P\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right)\times `$ (A24)
$`\left[\left(1{\displaystyle \frac{1}{L_x1}}\left({\displaystyle \frac{F_\alpha }{\overline{F}}}\right)\right){\displaystyle _0^{L_x1}}𝑑\nu \nu \stackrel{~}{P}(\nu ){\displaystyle _0^{L_x1}}𝑑\nu \stackrel{~}{P}(\nu )\right]`$
$`=`$ $`\overline{F}^2{\displaystyle _0^{L_x}}𝑑\omega (\omega 1)P(\omega )\left[\left(1{\displaystyle \frac{\omega 1}{L_x1}}\right){\displaystyle _0^{L_x1}}𝑑\nu \nu \stackrel{~}{P}(\nu ){\displaystyle _0^{L_x1}}𝑑\nu \stackrel{~}{P}(\nu )\right]`$ (A25)
$`=`$ $`\overline{F}^2{\displaystyle _0^{L_x}}𝑑\omega (\omega 1)P(\omega )\left[{\displaystyle _0^{L_x1}}𝑑\nu (\nu 1)\stackrel{~}{P}(\nu ){\displaystyle \frac{\omega 1}{L_x1}}{\displaystyle _0^{L_x1}}𝑑\nu \nu \stackrel{~}{P}(\nu )\right]`$ (A26)
$`=`$ $`\overline{F}^2[{\displaystyle _0^{L_x}}d\omega (\omega 1)P(\omega ){\displaystyle _0^{L_x1}}d\nu (\nu 1)\stackrel{~}{P}(\nu )`$ (A28)
$`{\displaystyle \frac{1}{L_x1}}{\displaystyle _0^{L_x}}d\omega (\omega 1)^2P(\omega ){\displaystyle _0^{L_x1}}d\nu \nu \stackrel{~}{P}(\nu )]`$
$`=`$ $`{\displaystyle \frac{\overline{F}^2}{L_x1}}{\displaystyle _0^{L_x}}𝑑\omega (\omega 1)^2P(\omega ).`$ (A29)
As the correlation function is normalized with respect to the average fluctuation size,
$`C_0(j0)`$ $`=`$ $`\underset{L_y\mathrm{}}{lim}{\displaystyle \frac{1}{L_y}}{\displaystyle \underset{l=1}{\overset{L_y}{}}}\left({\displaystyle \frac{{\displaystyle \underset{m=1}{\overset{L_x}{}}}\delta F_{l,m}\delta F_{l,m+j}}{{\displaystyle \underset{m=1}{\overset{L_x}{}}}\delta F_{l,m}^2}}\right)`$ (A30)
$`=`$ $`{\displaystyle \frac{\delta F_\alpha \delta F_\beta _{\alpha \beta }}{\delta F_\alpha ^2}}`$ (A31)
$`=`$ $`{\displaystyle \frac{\frac{\overline{F}^2}{L_x1}_0^{L_x}𝑑\omega (\omega 1)^2P(\omega )}{\overline{F}^2_0^{L_x}𝑑\omega (\omega 1)^2P(\omega )}}`$ (A32)
$`=`$ $`(L_x1)^1`$ (A33)
which is independent of the form of the $`P(F)`$. As $`L_x\mathrm{}`$, $`C_0(j0)0`$, as expected.
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# Precise Measurement of 𝐵⁰-(𝐵⁰)̄ Mixing Parameters at the Υ(4𝑆)
## Abstract
We describe a measurement of $`B^0\overline{B^0}`$ mixing parameters exploiting a method of partial reconstruction of the decay chains $`\overline{B}D^\pm \pi ^{}`$ and $`\overline{B}D^\pm \rho ^{}`$. Using 9.6 $`\times 10^6B\overline{B}`$ pairs collected at the Cornell Electron Storage Ring, we find $`\chi _d=0.198\pm 0.013\pm 0.014`$, $`|y_d|<0.41`$ at 95% confidence level, and $`|\mathrm{}e(ϵ_B)|<0.034`$ at 95% confidence level.
preprint: CLNS 00-1668 CLEO 00-6
B. H. Behrens,<sup>1</sup> W. T. Ford,<sup>1</sup> A. Gritsan,<sup>1</sup> J. Roy,<sup>1</sup> J. G. Smith,<sup>1</sup> J. P. Alexander,<sup>2</sup> R. Baker,<sup>2</sup> C. Bebek,<sup>2</sup> B. E. Berger,<sup>2</sup> K. Berkelman,<sup>2</sup> F. Blanc,<sup>2</sup> V. Boisvert,<sup>2</sup> D. G. Cassel,<sup>2</sup> M. Dickson,<sup>2</sup> P. S. Drell,<sup>2</sup> K. M. Ecklund,<sup>2</sup> R. Ehrlich,<sup>2</sup> A. D. Foland,<sup>2</sup> P. Gaidarev,<sup>2</sup> L. Gibbons,<sup>2</sup> B. Gittelman,<sup>2</sup> S. W. Gray,<sup>2</sup> D. L. Hartill,<sup>2</sup> B. K. Heltsley,<sup>2</sup> P. I. Hopman,<sup>2</sup> C. D. Jones,<sup>2</sup> D. L. Kreinick,<sup>2</sup> M. Lohner,<sup>2</sup> A. Magerkurth,<sup>2</sup> T. O. Meyer,<sup>2</sup> N. B. Mistry,<sup>2</sup> E. Nordberg,<sup>2</sup> J. R. Patterson,<sup>2</sup> D. Peterson,<sup>2</sup> D. Riley,<sup>2</sup> J. G. Thayer,<sup>2</sup> P. G. Thies,<sup>2</sup> B. Valant-Spaight,<sup>2</sup> A. Warburton,<sup>2</sup> P. Avery,<sup>3</sup> C. Prescott,<sup>3</sup> A. I. Rubiera,<sup>3</sup> J. Yelton,<sup>3</sup> J. Zheng,<sup>3</sup> G. Brandenburg,<sup>4</sup> A. Ershov,<sup>4</sup> Y. S. Gao,<sup>4</sup> D. Y.-J. Kim,<sup>4</sup> R. Wilson,<sup>4</sup> T. E. Browder,<sup>5</sup> Y. Li,<sup>5</sup> J. L. Rodriguez,<sup>5</sup> H. Yamamoto,<sup>5</sup> T. Bergfeld,<sup>6</sup> B. I. Eisenstein,<sup>6</sup> J. Ernst,<sup>6</sup> G. E. Gladding,<sup>6</sup> G. D. Gollin,<sup>6</sup> R. M. Hans,<sup>6</sup> E. Johnson,<sup>6</sup> I. Karliner,<sup>6</sup> M. A. Marsh,<sup>6</sup> M. Palmer,<sup>6</sup> C. Plager,<sup>6</sup> C. Sedlack,<sup>6</sup> M. Selen,<sup>6</sup> J. J. Thaler,<sup>6</sup> J. Williams,<sup>6</sup> K. W. Edwards,<sup>7</sup> R. Janicek,<sup>8</sup> P. M. Patel,<sup>8</sup> A. J. Sadoff,<sup>9</sup> R. Ammar,<sup>10</sup> A. Bean,<sup>10</sup> D. Besson,<sup>10</sup> R. Davis,<sup>10</sup> N. Kwak,<sup>10</sup> X. Zhao,<sup>10</sup> S. Anderson,<sup>11</sup> V. V. Frolov,<sup>11</sup> Y. Kubota,<sup>11</sup> S. J. Lee,<sup>11</sup> R. Mahapatra,<sup>11</sup> J. J. O’Neill,<sup>11</sup> R. Poling,<sup>11</sup> T. Riehle,<sup>11</sup> A. Smith,<sup>11</sup> J. Urheim,<sup>11</sup> S. Ahmed,<sup>12</sup> M. S. Alam,<sup>12</sup> S. B. Athar,<sup>12</sup> L. Jian,<sup>12</sup> L. Ling,<sup>12</sup> A. H. Mahmood,<sup>12,</sup><sup>*</sup><sup>*</sup>*Permanent address: University of Texas - Pan American, Edinburg, TX 78539. M. Saleem,<sup>12</sup> S. Timm,<sup>12</sup> F. Wappler,<sup>12</sup> A. Anastassov,<sup>13</sup> J. E. Duboscq,<sup>13</sup> K. K. Gan,<sup>13</sup> C. Gwon,<sup>13</sup> T. Hart,<sup>13</sup> K. Honscheid,<sup>13</sup> D. Hufnagel,<sup>13</sup> H. Kagan,<sup>13</sup> R. Kass,<sup>13</sup> T. K. Pedlar,<sup>13</sup> H. Schwarthoff,<sup>13</sup> J. B. Thayer,<sup>13</sup> E. von Toerne,<sup>13</sup> M. M. Zoeller,<sup>13</sup> S. J. Richichi,<sup>14</sup> H. Severini,<sup>14</sup> P. Skubic,<sup>14</sup> A. Undrus,<sup>14</sup> S. Chen,<sup>15</sup> J. Fast,<sup>15</sup> J. W. Hinson,<sup>15</sup> J. Lee,<sup>15</sup> N. Menon,<sup>15</sup> D. H. Miller,<sup>15</sup> E. I. Shibata,<sup>15</sup> I. P. J. Shipsey,<sup>15</sup> V. Pavlunin,<sup>15</sup> D. Cronin-Hennessy,<sup>16</sup> Y. Kwon,<sup>16,</sup>Permanent address: Yonsei University, Seoul 120-749, Korea. A.L. Lyon,<sup>16</sup> E. H. Thorndike,<sup>16</sup> C. P. Jessop,<sup>17</sup> H. Marsiske,<sup>17</sup> M. L. Perl,<sup>17</sup> V. Savinov,<sup>17</sup> D. Ugolini,<sup>17</sup> X. Zhou,<sup>17</sup> T. E. Coan,<sup>18</sup> V. Fadeyev,<sup>18</sup> Y. Maravin,<sup>18</sup> I. Narsky,<sup>18</sup> R. Stroynowski,<sup>18</sup> J. Ye,<sup>18</sup> T. Wlodek,<sup>18</sup> M. Artuso,<sup>19</sup> R. Ayad,<sup>19</sup> C. Boulahouache,<sup>19</sup> K. Bukin,<sup>19</sup> E. Dambasuren,<sup>19</sup> S. Karamov,<sup>19</sup> G. Majumder,<sup>19</sup> G. C. Moneti,<sup>19</sup> R. Mountain,<sup>19</sup> S. Schuh,<sup>19</sup> T. Skwarnicki,<sup>19</sup> S. Stone,<sup>19</sup> G. Viehhauser,<sup>19</sup> J.C. Wang,<sup>19</sup> A. Wolf,<sup>19</sup> J. Wu,<sup>19</sup> S. Kopp,<sup>20</sup> S. E. Csorna,<sup>21</sup> I. Danko,<sup>21</sup> K. W. McLean,<sup>21</sup> Sz. Márka,<sup>21</sup> Z. Xu,<sup>21</sup> R. Godang,<sup>22</sup> K. Kinoshita,<sup>22,</sup>Permanent address: University of Cincinnati, Cincinnati, OH 45221 I. C. Lai,<sup>22</sup> S. Schrenk,<sup>22</sup> G. Bonvicini,<sup>23</sup> D. Cinabro,<sup>23</sup> S. McGee,<sup>23</sup> L. P. Perera,<sup>23</sup> G. J. Zhou,<sup>23</sup> E. Lipeles,<sup>24</sup> M. Schmidtler,<sup>24</sup> A. Shapiro,<sup>24</sup> W. M. Sun,<sup>24</sup> A. J. Weinstein,<sup>24</sup> F. Würthwein,<sup>24,</sup><sup>§</sup><sup>§</sup>§Permanent address: Massachusetts Institute of Technology, Cambridge, MA 02139. D. E. Jaffe,<sup>25</sup> G. Masek,<sup>25</sup> H. P. Paar,<sup>25</sup> E. M. Potter,<sup>25</sup> S. Prell,<sup>25</sup> V. Sharma,<sup>25</sup> D. M. Asner,<sup>26</sup> A. Eppich,<sup>26</sup> J. Gronberg,<sup>26</sup> T. S. Hill,<sup>26</sup> R. J. Morrison,<sup>26</sup> H. N. Nelson,<sup>26</sup> and R. A. Briere<sup>27</sup>
<sup>1</sup>University of Colorado, Boulder, Colorado 80309-0390
<sup>2</sup>Cornell University, Ithaca, New York 14853
<sup>3</sup>University of Florida, Gainesville, Florida 32611
<sup>4</sup>Harvard University, Cambridge, Massachusetts 02138
<sup>5</sup>University of Hawaii at Manoa, Honolulu, Hawaii 96822
<sup>6</sup>University of Illinois, Urbana-Champaign, Illinois 61801
<sup>7</sup>Carleton University, Ottawa, Ontario, Canada K1S 5B6
and the Institute of Particle Physics, Canada
<sup>8</sup>McGill University, Montréal, Québec, Canada H3A 2T8
and the Institute of Particle Physics, Canada
<sup>9</sup>Ithaca College, Ithaca, New York 14850
<sup>10</sup>University of Kansas, Lawrence, Kansas 66045
<sup>11</sup>University of Minnesota, Minneapolis, Minnesota 55455
<sup>12</sup>State University of New York at Albany, Albany, New York 12222
<sup>13</sup>Ohio State University, Columbus, Ohio 43210
<sup>14</sup>University of Oklahoma, Norman, Oklahoma 73019
<sup>15</sup>Purdue University, West Lafayette, Indiana 47907
<sup>16</sup>University of Rochester, Rochester, New York 14627
<sup>17</sup>Stanford Linear Accelerator Center, Stanford University, Stanford, California 94309
<sup>18</sup>Southern Methodist University, Dallas, Texas 75275
<sup>19</sup>Syracuse University, Syracuse, New York 13244
<sup>20</sup>University of Texas, Austin, TX 78712
<sup>21</sup>Vanderbilt University, Nashville, Tennessee 37235
<sup>22</sup>Virginia Polytechnic Institute and State University, Blacksburg, Virginia 24061
<sup>23</sup>Wayne State University, Detroit, Michigan 48202
<sup>24</sup>California Institute of Technology, Pasadena, California 91125
<sup>25</sup>University of California, San Diego, La Jolla, California 92093
<sup>26</sup>University of California, Santa Barbara, California 93106
<sup>27</sup>Carnegie Mellon University, Pittsburgh, Pennsylvania 15213
The discovery of $`B^0\overline{B^0}`$ mixing in 1987 signaled a large top quark mass and allows the anticipation of observable CP violating asymmetries in the $`B^0`$ meson in the near future. Well-known values of the parameters describing mixing will be necessary to extract precise values of CP violating parameters, as planned at asymmetric $`B`$ factories . In addition, the size of $`B^0\overline{B^0}`$ mixing is characterized by the mass difference parameter $`\mathrm{\Delta }m_d`$ and is proportional to the square of $`|V_{tb}^{}V_{td}|`$, the magnitude of one side of the “Unitarity Triangle” which describes some of the mathematical constraints imposed by unitarity upon the elements of the CKM matrix . Accurate measurements of $`\mathrm{\Delta }m`$ for $`B^0`$ and $`B_s^0`$ mesons therefore provide an independent check on our understanding of CP violation in the Standard Model.
Mixing in the $`B^0\overline{B^0}`$ system may be described by the parameters $`x_d`$, $`y_d`$, $`p`$, and $`q`$. The parameter $`x_d=\mathrm{\Delta }m_d`$/$`\mathrm{\Gamma }_d`$, where $`\mathrm{\Delta }m_d`$ is the mass difference between the heavy and light eigenstates $`B_H`$ and $`B_L`$ and $`\mathrm{\Gamma }_d`$ is the average natural decay width. Similarly, $`y_d`$ is the normalized lifetime difference between the two eigenstates, and can be written as $`\mathrm{\Delta }\mathrm{\Gamma }_d`$/$`2\mathrm{\Gamma }_d`$. The parameters $`p`$ and $`q`$ describe the $`B^0`$ and $`\overline{B^0}`$ amplitudes, respectively, in the eigenstates $`B_H`$ and $`B_L`$. When the $`\mathrm{{\rm Y}}(4S)`$ is produced with symmetric energy electron positron collisions, the experimentally accessible quantity is $`\chi _d`$, the time-integrated probability for an initially produced $`B^0`$ or $`\overline{B^0}`$ meson to decay as its CP conjugate. It may be written in terms of $`x_d`$ and $`y_d`$:
$`\chi _d={\displaystyle \frac{\mathrm{\Gamma }(B^0\overline{B^0})}{\mathrm{\Gamma }(B^0B^0)+\mathrm{\Gamma }(B^0\overline{B^0})}}{\displaystyle \frac{x_d^2+y_d^2}{2(1+x_d^2)}}.`$ (1)
Under certain assumptions, a measurement of $`\chi _d`$ can be combined with direct determinations of $`\mathrm{\Delta }m_d`$ and the $`B`$ meson lifetime in order to extract $`y_d`$. The ancillary variable, $`ϵ_B`$, is analogous to the $`K^0`$-mixing parameter $`ϵ`$, and is defined through the relation $`p=(1+ϵ_B`$)/$`\sqrt{2(1+|ϵ_B|^2)}`$. Limits on $`ϵ_B`$ can be extracted by searching for a CP violating asymmetry in the events where mixing has occurred. The mixing parameters $`y_d`$ and $`ϵ_B`$ are both expected to be of order $`10^2`$ with considerable uncertainty.
In this letter, we report new measurements of the $`B^0\overline{B^0}`$ mixing parameters measured at the $`\mathrm{{\rm Y}}(4S)`$ resonance. We attempt to determine the beauty quantum number (or flavor) at decay for both of the $`B`$ mesons produced in the $`\mathrm{{\rm Y}}(4S)`$ decay using a novel method subject to systematic uncertainties very different from previous measurements ,. When the decay flavors of the two $`B`$ mesons in the event coincide, it indicates that the second $`B`$ has undergone mixing in the interval between the decays of the two $`B`$ mesons. The flavor of one $`B`$ meson at decay is tagged by a high-momentum lepton originating from the decay chain $`BX\mathrm{}\nu `$. The flavor at decay of the remaining $`B`$ meson is determined through partial reconstruction of the decay chain $`\overline{B^0}D^+h_W^{}`$ (charge conjugate modes implied), where $`h_W^{}`$ refers either to a $`\pi ^{}`$ or $`\rho ^{}`$. The electric charge of the $`h_W`$ identifies (tags) the value of the $`B`$ flavor at the time of its decay. (We assume the double Cabibbo suppressed decay $`\overline{B^0}D^{}h_W^+`$ is negligible.) By employing the hadronic flavor tag for one $`B^0`$ in the event, this method sacrifices statistical accuracy relative to methods where a semileptonic decay is used to tag the flavor of both $`B`$ mesons in the event. However, the systematic error due to the uncertainty in the charged to neutral $`B`$ meson production ratio at the $`\mathrm{{\rm Y}}(4S)`$ that dominates measurements of $`\chi _d`$ at the $`\mathrm{{\rm Y}}(4S)`$ using dileptons is substantially reduced . As a result this method results in a significant improvement in precision over previous measurements of $`\chi _d`$ at the $`\mathrm{{\rm Y}}(4S)`$.
Four charge combinations of hadrons and leptons are possible: $`h_W^+\mathrm{}^+`$, $`h_W^{}\mathrm{}^{}`$, $`h_W^+\mathrm{}^{}`$, and $`h_W^{}\mathrm{}^+`$. In the absence of backgrounds or mistags, these correspond to the four flavor combinations $`B^0B^0`$, $`\overline{B^0}\overline{B^0}`$, $`B^0\overline{B^0}`$, and $`\overline{B^0}B^0`$, respectively. Then,
$`\chi _d={\displaystyle \frac{h_W^+\mathrm{}^++h_W^{}\mathrm{}^{}}{h_W^+\mathrm{}^++h_W^{}\mathrm{}^{}+h_W^+\mathrm{}^{}+h_W^{}\mathrm{}^+}}.`$ (2)
In practice, the raw counts recorded in each of the combinations must be corrected for processes that incorrectly tag the $`B`$ decay flavor (mistags). Mistags may be due either to leptons not arising from the primary decay $`BX\mathrm{}\nu `$ or hadrons not arising from the hypothesized decay chain $`\overline{B^0}D^+h_W^{}`$. Backgrounds (non-$`\overline{B^0}`$ events), which can contribute either to the denominator or numerator, must be subtracted.
The data were recorded at the Cornell Electron Storage Ring (CESR) with two configurations of the CLEO detector called CLEO II and CLEO II.V. In the CLEO II.V configuration, the innermost wire chamber was replaced with a precision three-layer silicon vertex detector (SVX) . The results presented here are based upon an integrated luminosity of 9.1 $`\mathrm{fb}^1`$ of $`e^+e^{}`$ data taken at the $`\mathrm{{\rm Y}}(4S)`$ energy and 4.4 $`\mathrm{fb}^1`$ taken an average of 60 MeV below $`B\overline{B}`$ threshold. The Monte Carlo simulation of the CLEO detector was based upon GEANT and simulated events were processed in the same manner as the data.
The method of partial reconstruction used in this Letter has been described in detail elsewhere . By observing the $`h_W`$ and the soft pion from the decay $`D^+D^0\pi _s^+`$, we deduce the kinematics of the decay chain $`\overline{B^0}D^+h_W^{}`$ without reconstruction of the $`D^0`$. We denote the charged pion from the $`h_W`$ as $`\pi _f`$.
The hadronic $`B^0`$ decay may be described by three angles. The angle formed, in the $`B`$ ($`D^{}`$) rest frame, between the $`D^{}(D^0)`$ flight direction and the direction of the lab frame, is called $`\theta _B^{}(\theta _D^{}^{})`$. A larger value of $`\mathrm{cos}\theta _B^{}(\mathrm{cos}\theta _D^{}^{})`$ corresponds to a higher momentum of the $`h_W`$($`\pi _s`$). The third angle, $`\varphi `$, is the angle between the plane of the $`\overline{B^0}D^+h_W^{}`$ decay and the plane that contains the $`h_W^{}`$ and the $`\pi _s^+`$, as shown in Fig. 1. All three angles have distinctive distributions for signal and background. The $`\mathrm{cos}\theta _B^{}`$ distribution is constant for signal because the $`B`$ meson is a scalar particle. The distribution of $`\mathrm{cos}\theta _D^{}^{}`$ shows the 100% polarization in the $`D^+\pi ^{}`$ mode from conservation of angular momentum, and shows the 87% polarization that has been measured in the $`D^+\rho ^{}`$ mode. The distribution of $`\mathrm{cos}\varphi `$ is a combination of the $`\pi _s`$ and $`h_W`$ momenta and the angle between them. For most signal events it reconstructs inside the physical region $`1<\mathrm{cos}\varphi <1`$. For non-signal events as well as signal events with imperfect measurements of the pion momenta, $`\mathrm{cos}\varphi `$ can be calculated but may fail to have a physical value.
In order to select hadronic $`D^+h_W^{}`$ decay candidates for the analysis, both the $`\pi _s`$ and $`\pi _f`$ candidate tracks must be well-reconstructed and consistent with originating at the $`e^+e^{}`$ interaction point, and must not be identified as a lepton. We reconstruct $`\rho ^\pm `$ candidates from $`\pi _f^\pm \pi ^0`$ combinations, where the $`\pi ^0`$ is formed from a pair of photon candidates.
The $`h_W`$ momentum is required to fall in the kinematically allowed range for $`\overline{B^0}D^+h_W^{}`$ decays, assuming $`E_{B^0}=E_{beam}`$. We require the momentum of the $`\pi _s`$ to be below 300 MeV/$`c`$. The $`\pi _s`$ and $`\pi _f`$ are required to have opposite electrical charges. We require $`|\mathrm{cos}\varphi |<7`$ and $`\mathrm{cos}\theta _{h_W\pi _s}<0.8`$, where $`\theta _{h_W\pi _s}`$ is the angle between the $`h_W`$ and $`\pi _s`$. In 8% of selected events, there is more than one combination of charged tracks that satisfy these criteria. In this case, we select the one that is reconstructed as $`D^{}\pi `$ rather than $`D^{}\rho `$. If more than one combination still remains, we choose the one for which the value of $`\mathrm{cos}\varphi `$ is nearest 0.6 (the peak of the signal distribution.) The resulting bias in the $`\mathrm{cos}\varphi `$ shape has a negligible effect upon our measurement of the mixing parameters. The events satisfying these criteria are used to determine the sample composition. For the events we use to measure yields (and therefore the mixing parameters), we use more restrictive criteria, requiring $`\mathrm{cos}\theta _{h_W\pi _s}<0.95`$, $`|\mathrm{cos}\varphi |<2`$, and $`|\mathrm{cos}\theta _D^{}^{}|<1`$.
Lepton candidates are selected by requiring that the track is well-reconstructed, consistent with originating from the $`e^+e^{}`$ interaction point, and well identified as either an electron or a muon. We require that the momentum of the lepton candidate is greater than 1.4 GeV/$`c`$ and we use the angle between the lepton and the $`\pi _s`$ to suppress semileptonic decays from the unreconstructed $`D^0`$. We veto leptons from the decays $`J/\psi \mathrm{}^+\mathrm{}^{}`$. If the event has been partially reconstructed as $`D^+\pi ^{}`$, we require that the lepton form a large angle with the thrust axis of the remainder of the event in order to suppress $`e^+e^{}q\overline{q}`$ backgrounds, where $`q=\{u,d,s,c\}`$. If more than one lepton candidate in an event satisfies the criteria, we select the highest momentum candidate.
For the events selected by these criteria, the $`B`$ decay modes contributing to the $`D^+h_W^{}`$ candidates can be divided into five categories: signal ($`D^+\pi ^{}`$ and $`D^+\rho ^{}`$), other two-body and semi-leptonic $`\overline{B^0}`$ decays (such as $`\overline{B^0}D^+\rho ^{}`$ and $`\overline{B^0}D^+\mathrm{}^{}\overline{\nu }`$), two-body and semi-leptonic $`B^\pm `$ decays (such as $`B^{}D^0\pi ^{}`$ and $`B^{}D^0\mathrm{}^{}\overline{\nu }`$), random combinatoric backgrounds, and events of the type $`e^+e^{}q\overline{q}`$ (continuum). For the distributions of two-body $`B^0`$ and $`B^+`$ decays in $`\mathrm{cos}\theta _D^{}^{}`$ and $`\mathrm{cos}\varphi `$, we rely on the simulation. We include 10 two-body and semi-leptonic decay modes of the $`B^0`$ in the definition of two-body $`B^0`$ decays, and 12 in the definition of the two-body $`B^+`$ decays . These decays are well-measured and are reliably modeled by the simulation. Combinations of $`h_W`$ and $`\pi _s`$ that satisfy the analysis requirements, yet originate from neither the signal decays nor the two-body decays, are considered random combinatoric backgrounds. To model these, we use a synthetic distribution in $`\mathrm{cos}\theta _D^{}^{}`$ and $`\mathrm{cos}\varphi `$, generated by combining track pairs drawn at random from the observed spectrum of all $`B\overline{B}`$ track momenta in data. The cosine of the angle between them, $`\mathrm{cos}\theta _{h_W\pi _s}`$, is distributed uniformly. The simulation predicts that the distribution generated with this procedure provides an excellent approximation to the distributions of these decays, and indicates that $`40\%`$ of the combinatoric background comes from $`B^+B^{}`$ events. Distributions from continuum $`q\overline{q}`$ production are directly measured in data taken below the $`\mathrm{{\rm Y}}(4S)`$ resonance.
In order to determine the composition of our event sample, we divide the sample into two subsets based on the value of $`\mathrm{cos}\theta _{h_W\pi _s}`$. Events for which $`0.9<\mathrm{cos}\theta _{h_W\pi _s}<0.8`$ are the sideband sample; events for which $`\mathrm{cos}\theta _{h_W\pi _s}<0.95`$ make up the signal sample. We perform a binned two-dimensional maximum likelihood fit simultaneously to the $`\mathrm{cos}\varphi `$ and $`\mathrm{cos}\theta _D^{}`$ distributions of the signal sample and the sideband sample. The fit determines the normalization of the first four categories of events; that of the continuum is fixed by the relative luminosity of the data taken at and below the $`\mathrm{{\rm Y}}(4S)`$.
The results of the fit are shown in Table I. The two projections of the fits in the mixing sample (more restrictive selection criteria) are shown in Fig. 2. We find that the fit has a Baker-Cousins $`\chi ^2`$ of 129.7 for 151 degrees of freedom. From simulations we find that this corresponds to a confidence level of 91%.
We also show in Fig. 3 that the distributions of the subset of events in Fig. 2 that contribute to the numerator of Eqn. 2, ($`h^\pm l^\pm `$ events), are well described by the fit.
The partially reconstructed hadronic tag may incorrectly identify the flavor of the decaying $`B`$ meson. The dominant source of mistagged events is $`B^0`$ candidates formed from random combinations of tracks in $`B^+B^{}`$ or $`B^0\overline{B^0}`$ events. We determine that the mistag rate of combinatoric events is (21$`\pm `$12)% (where the uncertainty is statistical only), using a separate sample of fully-reconstructed $`BD^{}\mathrm{}\nu `$ decays in the data. We combine the composition as determined by the fit (Table I) with the mistag rates for each individual component to determine the hadronic mistag rate, also shown in Table I. We calculate a total hadronic mistag rate of (3.1$`\pm `$1.2)%. The uncertainty includes the statistical uncertainties from the fit and in the random combinatoric mistag rate.
The lepton in the event may also mistag the flavor of the decaying $`B`$ meson for several reasons. Leptons may arise from the secondary decay chain $`BDX,DX\mathrm{}\nu `$. The magnitude of this source is well-constrained by measurements of the $`BDX`$ spectrum and the known form factors governing $`D`$ semileptonic decays . We also correct for leptons from $`q\overline{q}`$ events, $`B\psi X`$, $`\psi \mathrm{}^+\mathrm{}^{}`$ events , misidentified hadrons , leptons from $`D_s^+`$ and other upper vertex ($`b\overline{c}`$) production , in-flight decays, $`\pi ^0`$ Dalitz decays, $`\gamma `$ conversions, and $`\delta `$ rays . Altogether we find that (3.6$`\pm `$0.5)% of electrons and (3.8$`\pm `$0.5)% of muons incorrectly tag the $`B`$ decay flavor. The uncertainties are the total systematic uncertainties obtained by adding in quadrature the uncertainties associated with each of the input branching fractions, spectra, and fake rates.
The yields for the mixing sample (more restrictive selection criteria) in the possible charge combinations, and subsequent corrections, are summarized in Table II. We correct the continuum subtracted raw yields of Table II for the mistag levels that we have determined for the leptonic and hadronic tags, then subtract $`B^+B^{}`$ background, which contributes to the denominator even when the beauty quantum number is correctly reconstructed. The total charged $`B`$ background is (13.3$`\pm 2.5)\%`$, which is the sum of the two-body $`B^+B^{}`$ decays and the $`40\%`$ of the combinatoric background that is attributed to $`B^+B^{}`$. The fully-corrected result is $`\chi _d=0.198\pm 0.013\pm 0.014`$. The systematic uncertainties are listed in Table III.
The largest systematic uncertainty in $`\chi _d`$ is due to the uncertainty in the total hadronic mistag rate which in turn is dominated by the uncertainty in the mistag rate of combinatoric background events. The systematic uncertainty in the combinatoric background mistag rate is determined by comparing the mistag rates of energetic pions in data and simulated events, using samples in which the $`B`$ flavor has been tagged using the decay $`D^+\mathrm{}^{}\overline{\nu }`$. Smaller contributions to the total mistag rate come from the statistical uncertainty of the fit and uncertainty in the two-body mistag rates. We evaluate mistag rates for hadronic and leptonic tags independently, assuming no correlation. The difference in efficiencies and mistag rates for mixed and unmixed events are both found to be consistent with zero in large samples of simulated signal events, and the uncertainty in $`\chi _d`$ reflects the statistical uncertainty of the finite simulation samples. We assign a systematic uncertainty in $`\chi _d`$ due to the uncertainty in the mistag rate totaling 9%.
The uncertainty in the $`B^+B^{}`$ background is dominated by uncertainty in the percentage of random combinations that are due to $`B^+B^{}`$ decays. The assigned uncertainty allows the fraction of random combinations arising from $`B^+B^{}`$ decays to vary uniformly from 0 to 100%.
The uncertainty due to the distributions in $`\mathrm{cos}\theta _D^{}^{}`$ and $`\mathrm{cos}\varphi `$ of two-body decays is evaluated by repeating the analysis, varying in turn each two-body decay mode’s weight according to the experimental uncertainty in its branching fraction and the limited statistics of the simulation.
The uncertainty due to the shape of the signal distribution is assessed by examining the variation of $`\chi _d`$ as the data are refit with modified signal distributions. Modifications include variations in $`D^{}\rho `$ polarization, overall tracking efficiency, solution multiplicity, and beam energy. The uncertainty due to the fitting distribution of combinatoric decays is evaluated by considering possible variations in the momentum spectrum of random tracks. We consider three variations, corresponding to decreasing the average momentum of soft tracks, increasing the average momentum of high-momentum tracks and changing the shape of the low-momentum spectrum as if there were 50% more $`D^{}`$’s than expected.
By comparing the yield of $`B^0B^0`$ candidate events to the yield of $`\overline{B^0}\overline{B^0}`$ candidates, we can limit $`ϵ_B`$ from this measurement . We compare $`\chi _d`$ in events with positively charged leptons to $`\chi _d`$ with negatively charged leptons, defining $`\chi _\pm `$= $`h_W^\pm \mathrm{}^\pm `$/($`h_W^{}\mathrm{}^\pm +h_W^\pm \mathrm{}^\pm `$) and $`A_{CP}`$ = ($`\chi _+\chi _{}`$)/($`\chi _++\chi _{}`$), where $`A_{CP}=4\mathrm{}e(ϵ_B)`$. Charge asymmetries in lepton identification cancel with this method. From studies of detection asymmetries for hadrons and from measurements of hadronic fake contributions, we find negligible systematic bias in the measurement and estimate a systematic uncertainty of 1.4% on $`A_{CP}`$. We determine $`A_{CP}`$ = 0.017$`\pm `$0.070$`\pm `$0.014, corresponding to $`\mathrm{}e(ϵ_B)=0.004\pm 0.018\pm 0.003`$, or $`|\mathrm{}e(ϵ_B)|<0.034`$ at the 95% confidence level.
We are also able to provide a non-trivial limit on $`y_d`$ using Eqn.(1) under two assumptions. We assume that any possible $`\mathrm{\Delta }\mathrm{\Gamma }_d`$ has negligible impact upon the extraction of $`\mathrm{\Delta }m_d`$ from the experimental results listed in , and we assume that indirect CP violation is not present ($`|p/q|=1`$). We combine our measurement of $`\chi _d`$ with the values $`\mathrm{\Delta }m_d=0.464\pm 0.018\mathrm{ps}^1`$ and $`\tau _{B^0}=1.56\pm 0.04\mathrm{ps}`$ , to find $`|y_d|<0.41`$ at 95% confidence level.
We have described a measurement of $`B^0\overline{B^0}`$ mixing parameters $`\chi _d`$ and $`ϵ_B`$, and have combined our result with direct measurements of $`\mathrm{\Delta }m_d`$ to extract limits on $`y_d`$. We exploit a method of partial reconstruction of the decay chains $`BD^\pm \pi ^{}`$ and $`BD^\pm \rho ^{}`$ subject to systematic uncertainties different from previous measurements. We note that this result is independent of previous CLEO mixing analyses which used leptons to tag the beauty quantum numbers at decay for both $`B`$ mesons in the event . This measurement provides the first non-trivial limits on $`y_d`$.
We gratefully acknowledge the effort of the CESR staff in providing us with excellent luminosity and running conditions. This work was supported by the National Science Foundation, the U.S. Department of Energy, the Research Corporation, the Natural Sciences and Engineering Research Council of Canada, the A.P. Sloan Foundation, the Swiss National Science Foundation, the Texas Advanced Research Program, and the Alexander von Humboldt Stiftung.
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# X-ray Fluctuations from the Slim Disk
## 1 Introduction
Recent X-ray observations report that not only stellar black hole candidates (SBHCs) in their low states (=faint sources), but also NLS1s and stellar SLJSs (=bright sources) exhibit X-ray fluctuations (variability). The fluctuations of bright sources are made in the optically thick ADAD, what we call SD (=bright disk). However, the time evolution of the SD has not been investigated so far. The numerical simulation by Manmoto et al. (1996) is well known as a time-evolution calculation of the optically thin ADAD (=faint disk). The disturbance added into the optically thin ADAD falls into the central star. The disk luminosity increases when the disturbance falls, and the light curve of this process is in good agreement with the X-ray shot configuration of the SBHC Cyg X-1 in its low state.
I investigate how the luminosities of disturbed SDs vary. I add a similar disturbance as Manmoto et al. (1996) into the SD, and as a result, a similar light curve to that of the optically thin ADAD is obtained.
The basic equations and numerical procedures are described in §2. The resultant time evolution and discussion will be presented in §3.
## 2 Basic Equations
I calculate the evolution of a one-dimensional axisymmetric disk. The basic equations are the same as those of Manmoto et al. (1996), and are those of mass conservation, momentum conservation, angular momentum conservation, and energy flow:
$$\frac{}{t}(r\mathrm{\Sigma })+\frac{}{r}(r\mathrm{\Sigma }v_\mathrm{r})=0,$$
(1)
$$\frac{}{t}(r\mathrm{\Sigma }v_\mathrm{r})+\frac{}{r}(r\mathrm{\Sigma }v_\mathrm{r}^2)=r\frac{W}{r}+r^2\mathrm{\Sigma }(\mathrm{\Omega }^2\mathrm{\Omega }_\mathrm{K}^2)W\frac{d\mathrm{ln}\mathrm{\Omega }_\mathrm{K}}{d\mathrm{ln}r},$$
(2)
$$\frac{}{t}(r^2\mathrm{\Sigma }v_\phi )+\frac{}{r}(r^2\mathrm{\Sigma }v_\mathrm{r}v_\phi )=\frac{}{r}(r^2\alpha W),$$
(3)
and
$$\frac{}{t}(r\mathrm{\Sigma }e)+\frac{}{r}(r\mathrm{\Sigma }ev_\mathrm{r})=\frac{}{r}(rWv_\mathrm{r})\frac{}{r}(r\alpha Wv_\phi )rF$$
(4)
where $`\mathrm{\Sigma }(\rho 𝑑z)`$ is the surface density, $`W(p𝑑z)`$ is the vertically integrated pressure, and $`e`$ is the internal energy of the accreting gas. $`\mathrm{\Omega }(=v_\phi /r)`$ and $`\mathrm{\Omega }_\mathrm{K}[=(GM/r)^{1/2}/(rr_\mathrm{S})`$\] are the angular frequency of the gas flow and the Keplerian angular frequency in the pseudo-Newtonian potential (Paczyński & Witta 1980), respectively, where $`M`$ is the mass of the central black hole and $`r_\mathrm{S}`$ is the Schwarzschild radius. I set the viscosity parameter to be $`\alpha =0.1`$.
To evaluate the radiative cooling rate, $`F`$, I consider black body radiation:
$$F=\frac{8acT^4}{3\tau _\mathrm{R}/2+\sqrt{3}},$$
(5)
where $`a`$ is the radiation constant, $`T`$ is the gas temperature, and $`\tau _\mathrm{R}`$ is the Rosseland optical depth.
I obtain the steady state solution of equations 14 and perform the time-evolution calculation using the steady state solution as the initial state. As for the outer boundary condition, I set all quantities to be those of the standard disk by Shakura & Sunyaev (1973). The inner boundary is set at $`r_{\mathrm{in}}=2.7r_\mathrm{S}`$, where a free boundary is adopted. I add a mass of $`\dot{M}\mathrm{}t`$ into the disk through the outer boundary at every time step ($`\mathrm{}t`$). I set $`\dot{M}=100L_{\mathrm{Edd}}/c^2`$ where $`L_{\mathrm{Edd}}`$ is the Eddington luminosity.
I add a perturbation (disturbance) to the initial state of the disk as
$$\frac{\delta \rho }{\rho }=k\mathrm{exp}\left[\left(\frac{rr_0}{\lambda /2}\right)^2\right].$$
(6)
Throughout these calculations, I assign the wavelength of the perturbation, $`\lambda =20r_\mathrm{S}`$. I set two parameters, the radius of the center of the perturbation, $`r_0`$, and the ratio between $`\delta \rho `$ and $`\rho `$ at $`r=r_0`$, $`k`$. These parameters are varied, with three calculations being performed. The parameter sets of the three calculations are listed in Table 1.
## 3 Results and Summary
Fig. 1 presents the light curves of Models A–C. All models show shot-like light curves.
Fig. 2 plots how the surface density, $`\mathrm{\Sigma }`$, varies after the perturbation is added for Model A. The dotted and solid lines represent the initial state and the later evolution, respectively. The perturbation does not change its configuration very much, and it does not rapidly decay. The time for the perturbation to propagate corresponds to the free-fall time. After the perturbation has damped, the disk structure is not globally changed.
The response of SDs to local disturbances has been examined by one-dimensional numerical simulations. It is generally believed that SDs are thermally stable. I, however, find that disturbances added into the accretion flow do not damp rapidly and decay with roughly the free-fall time. After the disturbance has damped, the global disk structure of the disk is not greatly modified. This can account for the persistent X-ray emission with substantial variations observed in NLS1s and SLJSs. When a perturbation is made in the SD, it decays and exhibits one X-ray shot. Since the structure of the SD does not globally vary much after a perturbation propagates, the next perturbation produced in the disk can also exhibit an X-ray shot. Repeating such processes continuously can make the substantial variability seen in the X-ray luminosities of SLJSs and NLS1s.
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# Correlated V/R and IR photometric variations in the Be/X-ray binary LS I +61∘ 235/RX J0146.9+6121
## 1 Introduction
RX J0146.9+6121 is one of the slowest high mass X-ray pulsar systems (Mereghetti et al. 1993; Hellier 1994). Its optical counterpart is the $`V=11.2`$, B1III-V star LS I +61 235 located at an estimated distance of 2.3$`\pm `$0.5 kpc (Motch et al. 1997; Coe et al. 1993; Reig et al. 1997a, hereafter R97). R97 derived the astrophysical parameters of the optical counterpart and reported V/R variations with a quasi-period of $`3`$ years. The line profile variability was attributed to the prograde precession of a one-armed mode confined in the Be star’s disc.
Be/X-ray binaries comprise approximately 70% of the more general class of high mass X-ray binaries, the other $``$ 30% containing evolved (luminosity class I and II) primaries. In a Be/X-ray binary the optical companion, a Be star, is characterised by an emission line spectrum and an infrared excess when compared to normal B-type stars of the same spectral type. These two observational properties have their origin in the cool gaseous quasi-Keplerian disc which lies on the equatorial plane of the central star. A neutron star revolves around the 10-20 $`\mathrm{M}_{}`$ primary in a rather eccentric orbit (e $``$ 0.2–0.8) and accretes material expelled by the Be star from the disc – in the form of a low velocity, high density equatorial wind – giving rise to the X-rays.
The physical properties of the Be star’s disc in Be/X-ray binaries have traditionally been considered to be the same as those of isolated Be stars. Indeed, the long-term variability characterised as disc-loss phases and V/R variations have been seen in both, isolated and Be/X-ray binaries (Okazaki 1997; Negueruela et al. 1998). However, it is not yet clear what role the compact companion in Be/X-ray binaries may play in the onset and subsequent development of the perturbation that gives rise to the asymmetric profiles or in the formation and loss of the disc. There is growing evidence that the circumstellar disc surrounding isolated Be stars and Be star in X-ray binaries may not share, on average, the same physical properties (Reig et al. 1997b; Negueruela et al. 1998; 1999)
In this paper we analyse optical, infrared and X-ray data and search for correlations between the characteristics of the radiation at these wavelengths. Our ultimate goal is to assess the validity of the global one-armed oscillation model and whether the compact companion has any effect on the V/R variability.
## 2 Observations and results
### 2.1 Optical data
Optical spectroscopic observations were made from the 2.6m telescope at the Crimean Astronomical Observatory (CAO), the 2m RCC telescope of the Bulgarian National Astronomical Observatory ”Rozhen” (BNAO), the 1.3m telescope of the Skinakas Observatory (SKI), in Crete (Greece), the 1.0m Jacobus Kapteyn Telescope (JKT), the 2.5m Isaac Newton Telescope (INT) in service mode and the 4.2m William Herschel Telescope (WHT). The last three telescope are located in the Observatorio del Roque de Los Muchachos (La Palma, Spain). Table 1 shows the journal of the observations and the instrument set-up. For a description of the observations not mentioned in the table see R97.
Fig 1 shows a selected sample of H$`\alpha `$ line profiles of LS I +61 235 covering the period 1991 August – 1999 September. Part of the new data (all observations taken after 1996 February plus those from Rozhen BNAO) are plotted together with some of the spectra presented in R97. The evolution of the V/R ratio can be seen in Fig 2$`a`$. Different shapes of the line have been represented by different symbols. Stars are used for $`V>R`$ points, triangles for $`V<R`$ phases, dots for single-peak lines and squares for shell profiles (i.e. when the central absorption goes beyond the continuum). Note that the shell phase is very brief and always occurs during the transition from $`V>R`$ to $`V<R`$. In contrast, the transition from $`V<R`$ to $`V>R`$ is separated by single-peak profiles. Strictly speaking there are not symmetric single-peak lines in our observations (as observed in pole-on stars) since one can always see flank inflections revealing a second peak. Nevertheless, we will refer to as single-peak phase the profiles seen during the transition $`V<R`$ to $`V>R`$. By fitting a sine function to the V/R curve we refined the V/R quasi-period to a value of 1240$`\pm `$30 days.
We have also measured the separation of the blue and red peaks by fitting two Gaussian functions to the line profile (Fig 2$`d`$). The peak separation gives a measure of the velocity field, assuming it to be Keplerian. The blue-dominated profiles seem to sample a wider range of velocities than red-dominated profiles: $`V>R`$ points spread over a velocity range 220–360 km s<sup>-1</sup>, whereas $`V<R`$ points distribute around 250–320 km s<sup>-1</sup>. Also, blue-dominated profiles with high values of the peak separation occur only at the end of the $`V>R`$ cycle, just before the shell phase.
The density wave does not only affect the H$`\alpha `$ line but also the Paschen series and He I $`\lambda `$6678 (Fig 3). Since these lines are formed at different regions inside the equatorial disc the perturbation must extend over a very wide region, hence confirming its global nature.
### 2.2 Infrared data
The infrared observations were taken with the Continuous Variable Filter (CVF) infrared photometer, using the 1.5-m Carlos Sánchez telescope, located at the Teide Observatory, in Tenerife (Spain). The data were reduced following the procedure described by Manfroid (1993). Instrumental values were transformed to the TCS standard system (Alonso et al. 1998). Table 2 shows the results of the infrared observations. For earlier observations the reader is referred to R97.
The combination of the old and new infrared data yields a very interesting result, namely the correlation between the V/R variations and the infrared magnitudes. Fig 2 shows the J and K light curves and the evolution of the colour index J-K, covering the period 1991 August to 1999 October. The light curves were rebinned into 30 day bins. The errors were calculated from the photometric errors depending on the number of points, $`N`$, in each bin: when the corresponding bin contained only one point we simply took the observation error, if $`1<N5`$ then we defined the error as $`|x_ix_m|/N`$ and if $`N>5`$ then we considered the standard deviation $`\sqrt{(x_ix_m)^2/N}`$.
As can be seen in Fig 2, there is a distinct modulation with an amplitude, from maximum to minimum, of $``$ 0.3 mag. The period of this modulation is 3.1–3.8 years, in good agreement with the optical V/R variability of 3.4 years. Interestingly, the infrared maxima occur in coincidence with the optical ($``$ MJD 49250) and X-ray ($``$ MJD 50630) outbursts. On the other hand, the infrared colours do not change drastically over the period of the observations. That is, the slope of the infrared continuum remained the same over the time covered by the observations.
### 2.3 X-ray data
LS I +61 235 is a persistent low-luminosity Be/X-ray binary. These systems are characterised by long pulse periods, low X-ray variability ($`L_{max}/L_{min}\mathrm{¡}\mathrm{}10`$) and low but permanent levels of X-ray emission (Reig & Roche 1999).
RX J0146.9+6121 was observed with the RXTE Proportional Counter Array (PCA) on March 21, 1998 (02:01–11:54 UT). Good time intervals were defined by removing data taken at low Earth elevation angle ($`<`$ 10) and during times of high particle background. An offset of only 0.02 between the source position and the pointing of the satellite was allowed, to ensure that any possible short stretch of slew data at the beginning and/or end of the observation was removed. All five PCA units were functioning during the entire observation. The total net exposure was 19223 s. Due to the relatively wide field of view of the PCA instrument (1 FWHM) the nearby X-ray source 4U 0142+61 (White et al. 1996) also contributed to the total flux. After correcting for collimator efficiency the contribution of 4U 0142+61 to the total observed flux in the energy range 2–30 keV was estimated to be of $``$ 10%. Because of this no X-ray spectral analysis was attempted.
The pulse period was determined by correcting the data to the solar system barycentre and using the epoch folding technique, i.e, we folded the data over a range of periods and searching for a maximum $`\chi ^2`$ as a function of period. The pulse period found was 1404.5$`\pm `$0.5 s, which is virtually the same as the 1404.2$`\pm `$1.2 s pulse period obtained in another RXTE observation nine months earlier (Haberl et al. 1998).
## 3 Discussion
### 3.1 Global $`m=1`$ oscillations
In R97 we investigated the different models that had been put forward to explain the V/R variability in Be stars and concluded that the model which best accounted for the observational data in LS I +61 235 was the Global One-armed Oscillation model (Okazaki 1991, 1997; Papaloizou et al. 1992). This model suggests that the long-term V/R variations are caused by global $`m=1`$ oscillations in the cool equatorial disc of the Be star. In other words, an enhanced density perturbation develops on one side of the disc, which slowly precesses. The precession time being that associated with the V/R quasi-period. The density perturbation is confined within a few stellar radii in the disc and the precession period turns out to be fairly insensitive to the size of the disc (Savonije & Heemskerk 1993).
One prediction of the model is that no changes in the slope of the infrared continuum are expected. The reason is that the V/R variations are not the result of changes in the radial gradient of the circumstellar gas. The slope of the infrared continuum is a measure of the radial density distribution but since matter in the disc does not move radially no changes in the shape of the infrared continuum are expected. This is exactly the behaviour that we find in the case of LS I +61 235. While the individual infrared photometric bands changed ($`\mathrm{\Delta }J\mathrm{\Delta }H\mathrm{\Delta }K0.3`$ magnitudes) the infrared colours remained unchanged (Fig. 2).
In principle, the issue of whether the motion of the perturbation occurs in the same sense (prograde rotation) or opposite sense (retrograde rotation) to the stellar rotation can be found out from the observations. Telting et al. (1994) realised that a prograde revolution implies that the $`V>R`$ phase must be followed by a shell profile and a similar profile but with much less pronounced absorption feature (or possibly a single peak line if the inclination is low) during the transition from $`V<R`$ to $`V>R`$. These characteristic line shapes must translate into noticeable photometric variations. According to Mennickent et al. (1997), we should expect a minimum of brightness when $`V=R`$ prior to the $`V<R`$ phase if the motion is prograde and $`V>R`$ after $`V=R`$ if retrograde. In LS I+61 235 the minimum of brightness in the infrared photometric bands occurred during the shell phase ($`V=R`$) before the $`V<R`$ phase began, confirming the prograde nature of the precession inside the disc.
However, models of one-armed global density waves cannot reproduce strong shell-non-shell transitions like the ones seen in LS I +61 235 which are reminiscent of the so-called spectacular variations (Doazan et al. 1983). Such shell events seem to be rare a phenomenon which have been reported for only three Be stars: $`\gamma `$ Cas, 59 Cyg and Pleione (Hummel 1998 and references therein), all of which are either binaries or suspected binaries. One possible explanation might be a thick disc in the region where the perturbation lies. When the perturbation is behind the star and for the right inclination angle, no shell event is seen because the disc in between the central star and the observer is thin and does not occult the star nor the perturbation. The shell phase would occur when the perturbation is at inferior conjunction since the thicker disc would hide the central star from the observer. An alternative explanation is given by Hummel (1998) who suggested a tilted or warped circumstellar disc with precessing nodal line in addition to the density wave (see also Porter 1998).
### 3.2 A dense circumstellar disc in LS I +61 235
The works by Dachs et al. (1986) and Hanuschik et al. (1988) have shown that the equivalent width of H$`\alpha `$ line emission for Be stars increases with the effective disc radius. Since for rotationally dominated profiles $`\mathrm{\Delta }V/(2v\mathrm{sin}i)`$ can be regarded as a measure of the radius of the H$`\alpha `$ emitting region (Huang 1972), we expect a correlation between the peak separation and the H$`\alpha `$ equivalent width. Hanuschik et al. (1988) derived the following law
$$\mathrm{log}\left(\frac{\mathrm{\Delta }V}{2v\mathrm{sin}i}\right)=a\mathrm{log}(EW(H\alpha ))+b$$
(1)
where $`v\mathrm{sin}i`$ is the projected rotational velocity and EW(H$`\alpha `$) is given in anstrongs. $`a`$ and $`b`$ are related to the rotational law index $`a=j/2`$ ($`j`$=0.5 for Keplerian rotation and $`j`$=1 for conservation of angular momentum) and with the disc electron density, respectively. A least square fit to the LS I +61 235 data gave $`a=0.23\pm 0.10`$, i.e. $`j0.5`$ and $`b=+0.1\pm 0.1`$. These values are to be compared with the average values $`a=0.4`$ and $`b=0.1`$ found by Hanuschik et al. (1988) for a sample of 26 isolated Be stars. The higher value of $`b`$ in LS I +61 235 implies a denser disc than those of isolated Be stars.
The main (and probably only) difference between a Be/X-ray binary and an early-type isolated Be star is the presence of a neutron star in the former. It then seems natural to attribute the dissimilarity in the properties of the circumstellar envelopes to the influence of such compact companion. The neutron star trims the disc to a certain radius and prevents its free growth, making it denser. Disc truncation has been suggested in other Be/X-ray binaries like V0332+53 (Negueruela et al. 1999). This idea would support the hypothesis that the neutron star plays a fundamental role in the evolution and properties of the equatorial disc in Be/X-ray binaries, as proposed by Reig et al. (1997b).
### 3.3 X-ray/optical/IR correlations
In R97 the observation of an optical outburst around 1993 September-October (MJD $``$ 49260) was reported. During the outburst the H$`\alpha `$ equivalent width, EW(H$`\alpha `$), changed by $``$ 10 Å, in about 270 days, decreasing to pre-outburst values ($`8`$Å) in about the same amount of time (Fig 2$`c`$). This increased coincided with a single-peak phase of the H$`\alpha `$ profile. The question of whether this outburst was an isolated event or was associated with the V/R cycle remained open due to the short coverage of the data. The new observations show no new EW(H$`\alpha `$) maximum. After the outburst the EW(H$`\alpha `$) increased slowly up to a level of $`12`$ Å, considerably lower than the peak of 1993 September. The new single-peak phase should have occurred during 1997 March-May. Unfortunately, the star was too close to the Sun to be observed. However, we notice that EW(H$`\alpha `$) seems to have reached a maximum value just before and after the period when we would expect the single-peak phase (around MJD 50500). Thus, we are inclined to think that the higher EW(H$`\alpha `$) associated with single-peak profiles reported in R97 may reflect the fact that the fluxes of both components are adding up and are not affected by the absorption feature present in the other profiles. In this context then, the increase in EW(H$`\alpha `$) is real and would be an event related to the motion of the density pattern in the disc, rather than an isolated episode.
In 1997 July (MJD 50634) LS I +61 235 underwent a small X-ray outburst (Haberl et al. 1998). The X-ray luminosity increased by nearly a factor of five in one week reaching 3.45 $`\times `$ 10<sup>35</sup> erg s<sup>-1</sup> in the energy range 0.5–10 keV. This outburst is marked with and arrow in Fig 2$`c`$. Interestingly, the outburst occurred at the time of the expected infrared and EW(H$`\alpha `$) maxima. One is then tempted to attribute this correlated X-ray/optical/IR behaviour to the high density perturbation where most of the Balmer emission is formed. If the inclination of the system is less than 90, when the high-density part of the equatorial disc is behind the star it offers the largest geometric area (especially so if the disc is thicker in this region) and the highest optical and infrared emission. If the neutron star happens to be close to the Be star it will accrete from this high-density material and the X-ray emission will be enhanced. New optical and X-ray observations around the next expected maximum are needed to solve this issue.
## 4 Conclusion
Optical spectroscopic observations confirm the presence of global one-armed oscillations in the circumstellar disc of LS I +61 235. These oscillations manifest themselves as quasi-periodic variations in the shape of the H$`\alpha `$ line, whose asymmetric double peak profile alternates between red- (V$`<`$R) and blue-dominated (V$`>`$R) emission. The system also goes through Be-Be shell transitions, which might indicate an asymmetric vertical structure of the disc in the form of a thick or a tilted disc. The V/R quasi-period is determined to be $``$ 1240$`\pm `$30 days. We have found a correlation between the infrared emission and the V/R variations. This is the first time that such correlation is reported in a Be/X-ray binary. From the pattern traced by the IR light curves in relation to the V/R ratio we conclude that the one-armed disk oscillations are prograde. The Be star’s disc in LS I +61 235 is found to be denser than isolated Be stars, which may be connected to the presence of the neutron star. From a Rossi X-ray Timing Explorer observation we derive a spin period of the neutron star of 1404.5$`\pm `$0.5 s.
### Acknowledgements
We thank Chris Moran for providing us with one of the INT spectra and Dr E. V. Paleologou for helping us with the spectroscopic observations at Skinakas Observatory. Skinakas Observatory is a collaborative project of the University of Crete, the Foundation for Research and Technology-Hellas and the Max-Planck-Institut f r Extraterrestrische Physik. P. Reig acknowledges support via the European Union Training and Mobility of Researchers Network Grant ERBFMRX/CT98/0195. IN is supported by an ESA external fellowship. R. Zamanov acknowledges support from Dirección general de relaciones culturales y científicas, Spain. Some observations were taken as part of the ING service observing programme. We are grateful to the referees, Dr D. Baade and Dr T. Rivinius, for useful comments.
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# A Nonlinear Approximation of Operator Equation 𝑉^∗𝑄𝑉=𝑄 : Nonspectral Decomposition of Nonnormal Operator and Theory of Stability
## 1 Introduction
Throughout this paper $`H`$ will denote a Hilbert space with scalar product $`<,>`$, $`V`$ denotes a linear bounded bijection $`H`$ onto $`H`$,
$$r(T):=\text{spectral radious of }T.$$
We will discuss the structure of the next four sets:
$$Stab_+(V):=\{xH|a>1C0N0V^NxCa^N\}.$$
$$Stab(V):=\{xH|V^nxM\text{ for some real }M\text{ and every }n=0,1,2,\mathrm{}\}$$
$$Stab_0(V):=\{xH|V^nx0(n\mathrm{})\}$$
$$l_2(V):=\{xH|Vx^2+V^2x^2+\mathrm{}+V^nx^2+\mathrm{}<\mathrm{}\}$$
Recall $`V`$ is similar to an unitary operator iff there exists a bounded uniformly positive operator $`Q`$ such that
$$V^{}QV=Q$$
With this equation we shall consider an ’approximation’ equation ( parametrized by real $`t`$)
$$Q=V^{}\frac{Q+t}{I+tQ}V,Q0,0<t1$$
$`(^{})`$
(hereinafter $`t`$ denotes always a real number such that $`0<t1`$ and if no confusion can occur we shall often write $`t`$ instead of $`tI`$, $`I`$ is identity operator).
The interest in this equation can be motivated by the next
Example 1. Let $`V`$ be normal. Routine, though tedious calculation shows that
$$Q_t:=\left[\frac{IV^{}V}{2}+\sqrt{\left(\frac{IV^{}V}{2}\right)^2+t^2V^{}V}\right]\frac{1}{t}$$
is uniformly positive solution of $`(^{})`$, there holds
$$Q_t^1=\left[\frac{I(V^{}V)^1}{2}+\sqrt{\left(\frac{I(V^{}V)^1}{2}\right)^2+t^2(V^{}V)^1}\right]\frac{1}{t}$$
and there exist
$$X_0:=stronglimtQ_t=(V^{}VI)E(1,\mathrm{})(t+0)$$
$$Y_0:=stronglimtQ_t^1=((V^{}V)^11)E(0,1)(t+0)$$
Besides that it is fast evident that $`Q_tE([1])=E([1]),Q_tE(1,\mathrm{})`$ is monotone increasing (with $`t+0`$), $`Q_tE(0,1)`$ is monotone decreasing and there exists an
$$R_0:=stronglim(I+Q_t)^1=E(0,1)+E([1])/2(t+0).$$
here $`E(\mathrm{\Delta })`$ denotes the spectral function of the selfadjoint operator $`V^{}V`$. Note:
$$IR_0=E([1])/2+E(1,\mathrm{})$$
$$Ker(IR_0)=E(0,1)H$$
$$RanX_0E(1,\mathrm{})H=\overline{RanX_0}$$
$$RanY_0E(0,1)H=\overline{RanY_0}$$
$`\mathrm{}`$
Also, in the above considered case of the normal $`V`$ it is estableshed that the operators $`X_0,Y_0,R_0`$ define (in essential) the spectral subspaces of $`V`$ (with $`V`$ together one can consider $`aVb,b/aspectrumV`$) . In this article we shall show that the similar situation holds for the arbitrary bounded bijection $`V`$.
We follow standards of \[RS\] when we apply mathematical concepts and sometimes we apply P.A.M. Dirac’s ‘bra-ket’ syntax.
We will often cite some assertions and propositions of \[Ch1-4\]. For the most convenient and accesible way to do it, we collect them together and resume them here as
Theorem 1.
(i) the solution of $`(^{})`$ exists and it is unique; denote it by $`Q_t`$
(ii) $`Q_t`$ is bounded selfadjoint uniformly positive and there are satisfied inequalities:
$$tV^{}VQ_tV^{}V/t$$
$$tQ_tsQ_s(0<ts1)$$
denote $`X_t:=tQ_t`$
(iii) $`Q_t^1`$ is (unique) solution of the analogous equation:
$$Q_t^1=V^1\frac{Q_t^1+t}{I+tQ_t^1}(V^1)^{}$$
so there are satisfied inequalities:
$$tV^1(V^1)^{}Q_t^1V^1(V^1)^{}/t$$
$$tQ_t^1sQ_s^1(0<ts1)$$
denote $`Y_t:=tQ_t^1`$.
(iv) Let
$$X_0:=stronglimX_t,Y_0:=stronglimY_t(t+0)$$
The operators $`X_0,Y_0`$ are bounded positive and they are maximal solutions of the equations
$$X=V^{}\frac{X}{I+X}V,X0,\text{resp.}VYV^{}=\frac{Y}{I+Y},Y0$$
(’maximal’ denotes ’maximal with respect to usual partial order for bounded operators on Hilbert space’)
(v) There hold the formulae
$$Y_0=stronglim(V^{}V+V^2V^2+\mathrm{}+V^nV^n)^1,(n\mathrm{})$$
$$X_0=stronglim((V^{}V)^1+(V^2V^2)^1+\mathrm{}+(V^nV^n)^1)^1,(n\mathrm{})$$
(vi) Denote $`R_t:=(I+Q_t)^1`$. Then $`0R_tI`$, $`Q_t=R_t^1I`$ and the equation $`(^{})`$ is equivalent to the equation
$$[R_t+t(IR_t)](V^1)^{}(IR_t)=[(IR_t)+tR_t]VR_t$$
(vii) Let $`R_0`$ be a weak operator limit point of the net $`\{R_t,t+0\}`$ (it is clear that $`R_0`$ exists and $`0R_0I`$). Then
$$VKer(IR_0)=Ker(IR_0)$$
$$V^1KerR_0=KerR_0$$
In particular,
$$V\overline{RanR_0}=\overline{RanR_0}\mathrm{}$$
## 2 Equation $`Q=V^{}{\displaystyle \frac{Q+t}{I+tQ}}V`$ . General Properties.
Hereinafter F denotes an arbitrary ultra filter, which majorizes usual convergence to $`+0`$. We will write
$$t+0\text{ instead of }t\stackrel{𝐅}{}+0$$
if no confusion can occur.
Definition 1.
$$FinQ:=\{x𝐇|<x,Q_tx>M_x\text{ for some real }M_x\text{ and for almost every }t\text{ resp. }𝐅\}$$
$$KerQ_0:=\{x𝐇|<x,Q_tx>0\text{ for }t+0\}$$
Theorem 1.
$$RanY_0^{1/2}=l_2(V)Stab_0(V)KerQ_0Ker(IR_0)\overline{RanR_0}KerX_0$$
$`()`$
$$Stab(V)FinQ\overline{RanR_0}KerX_0$$
$`()`$
Every set of these series is $`V`$-surinvariant. (recall, some $`L`$ is said to be $`T`$-surinvariant, iff $`TL=L`$).
In addition
$$FinQRanR_0^{1/2}$$
Observation 0. It is well-known and evident that
$$\begin{array}{cccccc}A^{}AB^{}B& & AxBx(xH)& & RanA^{}RanB^{}& \\ & & & \text{(}A\text{ and }B\text{ are bounded)}.& & \end{array}$$
Corollary 0. Let $`xH`$ , $`Y`$ be selfadjoint and let $`Y0`$. Then
$$xRanY^{1/2}xx|cY\text{ for some real }c\text{}$$
Proof.
Proof of $`xx|cYxRanY^{1/2}`$:
Take into account Observation 0. Then obtain
$$xx|cYRanxx|RanY^{1/2}xRanY^{1/2}.$$
Now proof of $`xRanY^{1/2}xx|cY`$ :
By the definition of $`Ran`$
$$xRanY^{1/2}x=Y^{1/2}y\text{ for an }y\text{.}$$
Hence
$$xRanY^{1/2}xx|Y^{1/2}yY^{1/2}y|=Y^{1/2}yy|Y^{1/2}y^2Y.$$
Now denote $`y^2`$ by $`c`$.
$`\mathrm{}`$
Observation 1. Let $`\{A_t\}_t`$ be a net of selfadjoint positive bounded operators. Suppose $`A_taI`$ for some positive number $`a`$ (and every $`t`$), $`A_0`$ be a weak operator limit point of this net (clear: $`A_0`$ exists and $`0A_0aI`$ ).
Then
$$<x,A_tx>0A_tx0xKerA_0$$
Proof. Clear (see e.g. \[Ch1\]).
$`\mathrm{}`$
Observation 2. Let $`Q`$ be bounded selfadjoint positive,$`0<t1`$ . Then
$$\frac{Q}{I+tQ}\frac{Q+t}{I+tQ}Q+t$$
Corollary.
$$\overline{)VFinQ=FinQ,VKerQ_0=KerQ_0.}$$
Observation 2’. Let $`Q`$ be bounded selfadjoint positive,$`0<t1`$ . Then
$$\frac{1t}{1+t}Q+\frac{2t}{1+t}\frac{Q+t}{I+tQ}=\frac{t(1t)}{1+t}\frac{(QI)^2}{I+tQ}$$
In particular
$$\frac{Q+t}{I+tQ}\frac{1t}{1+t}Q+\frac{2t}{1+t}$$
Denote $`z:=(1t)/(1+t)`$. Then $`1z=2t/(1+t)`$ and with these denotations
$$\frac{Q+t}{I+tQ}zQ+(1z)$$
Note
$$t+0z10.$$
Proof. Clear.
$`\mathrm{}`$
Observation 3 . For $`Q_t`$ the just mentioned inequality gives
$$Q_tV^{}(zQ_t+(1z))V$$
and with iterating this inequality one can obtain
$$Q_t(1z)[V^{}V+zV^2V^2+\mathrm{}+z^{n1}V^nV^n]+z^nV^nQ_tV^n$$
In particular, given numbers $`z,M_0`$ and an $`xH`$ such that
$$0<z<1,M_00,V^nx^2M_0(n=1,2,3,\mathrm{})$$
then
$$<x,Q_tx>(1z)(Vx^2+zV^2x^2+z^2V^3x^2+\mathrm{})M_0.$$
In particular
$$\overline{)Stab(V)FinQ}$$
Now suppose $`V^nx0`$ ($`n\mathrm{}`$) for some $`xH`$ . Let $`t+0`$. Then $`z10`$ and hence $`<x,Q_tx>0`$ since $`V^nx0`$ and $`<x,Q_tx>0`$.
In particular
$$\overline{)Stab_0(V)KerQ_0}$$
$`\mathrm{}`$
Observation 4. Let $`xH`$, $`x=1`$. Then
a)
$$xFinQM0\alpha 0:xx|(M+\alpha )(Q_t+\alpha )^1$$
b) Given some $`M,\alpha 0`$, such that
$$xx|(M+\alpha )(Q_t+\alpha )^1$$
then
$$xFinQ$$
Proof
$`<x,Q_tx>M`$
$``$ $`<x,(Q_t+\alpha )x>M+\alpha `$
$``$ $`(Q_t+\alpha )^{1/2}x^2M+\alpha `$
$``$ $`(Q_t+\alpha )^{1/2}x(Q_t+\alpha )^{1/2}x|M+\alpha `$
$``$ $`(Q_t+\alpha )^{1/2}xx|(Q_t+\alpha )^{1/2}M+\alpha `$
$``$ $`xx|(M+\alpha )(Q_t+\alpha )^1`$
$`\mathrm{}`$
Corollary. $`FinQRanR_0^{1/2}\overline{RanR_0}`$
Proof First apply Observation 4 for $`\alpha :=1`$ and let $`t+0`$. Then obtain
$$xx|(M+1)R_0;$$
Now apply Observation 0 or Corollary 0.
$`\mathrm{}`$
Observation 5. Let $`xH`$ . Then
$$xRanY_0^{1/2}xx|cY_0<x,Q_tx>ct$$
Proof
$$xRanY_0^{1/2}x=Y_0^{1/2}yxx|Y_0^{1/2}yY_0^{1/2}y|y^2Y_0$$
Besides,
$$xx|cY_0Ranxx|RanY_0^{1/2}xRanY_0^{1/2}$$
Now recall that
$$Y_0Y_t\text{ und }Y_t\stackrel{s}{}Y_0.$$
Hence
$`xRanY_0^{1/2}`$ $``$ $`xx|cY_0`$
$``$ $`xx|cY_t`$
$``$ $`xx|ctQ_t^1`$
$``$ $`Q_t^{1/2}xQ_t^{1/2}x|ct`$
$``$ $`Q_t^{1/2}x^{1/2}ct`$
$``$ $`<x,Q_tx>ct`$
$`\mathrm{}`$
Proof of $`KerQ_0Ker(IR_0).`$
Note
$$IR_t=Q_t(I+Q_t)^1Q_t$$
Now suppose $`xKerQ_0`$ i.e. $`<x,Q_tx>0`$ and apply the Observation 1 to $`A_t=IR_t`$.
$`\mathrm{}`$
Proof of $`Ker(IR_0)\overline{RanR_0}.`$
Note $`Ker`$ is closed and use definitions of $`Ker,Ran`$.
$`\mathrm{}`$
Proof of $`\overline{RanR_0}KerX_0.`$
Note
$$R_tX_t=t(I+Q_t)^1=X_tR_t$$
Then
$$R_tX_t=X_tR_t\stackrel{}{}0.$$
Recall $`R_t,X_t`$ are selfadjoit, positive, bounded and $`X_t\stackrel{s}{}X_0`$, $`R_t\stackrel{w}{}R_0`$. Hence $`R_0X_0=0=X_0R_0`$.
The rest is evident.
$`\mathrm{}`$
Proof of $`RanY_0^{1/2}l_2(V).`$
We have
$$VY_0V^{}=\frac{Y_0}{I+Y_0}$$
$`(^{})`$
Hence $`Y_0^{1/2}V^{}x=(I+Y_0)^{1/2}Y_0^{1/2}xY_0^{1/2}x`$ and the relation
$$W:Y_0^{1/2}xY_0^{1/2}V^{}x$$
defines a contraction $`RanY_0^{1/2}RanY_0^{1/2}`$. This contraction has extension to a contraction $`\overline{RanY_0^{1/2}}\overline{RanY_0^{1/2}}`$. It will be denote by $`W`$ too. Clear, there hold
$$\begin{array}{cc}(i)WY_0^{1/2}=Y_0^{1/2}V^{}\hfill & (iii)W^{}W=(I+Y_0)^1|\overline{RanY_0^{1/2}}\hfill \\ (ii)Y_0^{1/2}W^{}P=VY_0^{1/2}\hfill & (iv)Y_0P=((W^{}W)^1I)P\hfill \end{array}$$
and $`WP=PWP,W^{}P=PW^{}P`$ .Here is
$$P:=orthoprojectiononto\overline{RanY_0^{1/2}}$$
Hence
$$Y_0^{1/2}V^nV^nY_0^{1/2}=W^nY_0W^nP=PW^nY_0W^nP$$
$$=PW^n((W^{}W)^1I)W^nP$$
$$=P(W^{n1}W^{n1}W^nW^n)P.$$
and
$$\underset{n=2}{\overset{N}{}}V^nY_0^{1/2}x^2=WPx^2W^NPx^2$$
Hence
$$\underset{0}{\overset{\mathrm{}}{}}V^nY_0^{1/2}x^2<\mathrm{}\mathrm{}$$
$`\mathrm{}`$
Remark to this proof.
It follows from $`(^{})`$ and $`Y_00`$ that
$$VRanY_0=RanY_0,V^{}KerY_0=KerY_0,V\overline{RanY_0}=\overline{RanY_0}$$
Besides that,
$$\begin{array}{ccc}\hfill (iii)& =>& Wxx/(1+Y_0)(xD_W)\hfill \\ \hfill (ii)& =>& VRanY_0^{1/2}RanY_0^{1/2}\hfill \end{array}$$
$$\text{(}i\text{, or definition of }W\text{)}=>WRanY_0^{1/2}=RanY_0^{1/2}\mathrm{}$$
Proof of $`l_2(V)RanY_0^{1/2}.`$
Let $`xl_2(V)`$ and set
$$\begin{array}{ccc}\hfill c& :=& Vx^2+V^2x^2+V^3x^2+\mathrm{}\hfill \\ \hfill P& :=& orthoprojectionmappingontospan\{x\}\hfill \end{array}$$
Then
$$<x,Y_t^1x>=<x,\frac{1}{t}Q_tx>\frac{1z}{t}c=\frac{2c}{1+t}\mathrm{\hspace{0.33em}2}c$$
Hence
$$x^2P\mathrm{\hspace{0.17em}2}cY_t,x^2P\mathrm{\hspace{0.17em}2}cY_0,RanPRanY_0^{1/2}$$
$`\mathrm{}`$
Remark to this proof.
$$\begin{array}{cc}\hfill (i)& Y_t\mathrm{\hspace{0.17em}0},Y_0\mathrm{\hspace{0.17em}0},P=P^{1/2}\mathrm{\hspace{0.17em}0}\hfill \\ \hfill (ii)& AxBx(xH)=>RanA^{}RanB^{}\hfill \\ & \text{(}A\text{ and }B\text{ are bounded)}.\mathrm{}\hfill \end{array}$$
Corollary. $`l_2(V)=RanY_0^{1/2}.`$
## 3 Nonspectral Decomposition
Observation 1.
It follows from the definitions of $`X_t,Y_t`$ that $`X_tY_t=Y_tX_t=t^2`$ . Hence $`X_0Y_0=Y_0X_0=\mathrm{\hspace{0.33em}0}`$ and $`\overline{RanX_0}KerY_0`$ , $`\overline{RanY_0}KerX_0`$ . But $`X_0,Y_0`$ are selfadjoint. Thus we obtain an orthogonal decomposition
$$H=\overline{RanY_0}+(KerX_0KerY_0)+\overline{RanX_0}$$
such that
1) first component is $`V`$-surinvariant;
2) third component is $`V^1`$-surinvariant.
(recall, some $`L`$ is said to be $`T`$-surinvariant, iff $`TL=L`$). Moreover, denote $`j_t:=(I+tQ_t^1)^{1/2}/(I+tQ_t)^{1/2}`$ ,then $`j_t`$ is uniformly positive, bounded, there hold
$$(1+V^2)^{1/2}j_t(1+V^1^2)^{1/2},$$
and $`(j_tV)^{}Q_t(j_tV)=Q_t`$. In particular $`j_tV`$ is similar to an unitary operator.
For a moment suppose $`dimH<\mathrm{}`$. It is clear that now the restriction of $`V`$ onto $`\overline{RanY_0}`$ and the restriction of $`V^1`$ onto $`\overline{RanX_0}`$ are similar to uniform contractions (see theorems 2.1 with remarks to the proof of $`l_2(V)=RanY_0^{1/2}`$). In addition, if
$$RanY_0=\{0\}=RanX_0,$$
then the restriction of $`V`$ onto $`KerX_0KerY_0`$ has the unit spectrum. $`\mathrm{}`$
This motivates the
Definition 1. We shall say, a linear bounded operator $`T`$ is near similar to uniform contraction , iff there exists a bounded operator $`Y>\mathrm{\hspace{0.33em}0}`$ such that
$$TYT^{}=\frac{Y}{I+Y}$$
We shall say, T is s-approximately similar to an unitary , iff there exists a net $`\{j_t\}_t`$ of bounded uniformly positive operators such that
1) $`1/Mj_tM`$ for some real $`M`$ (and every $`t`$ ),
2) $`stronglimj_t=I`$ (hence $`stronglimj_t^1=I`$),
3) for every fixed $`t`$ the operator $`j_tT`$ is similar to an unitary operator. $`\mathrm{}`$
With this definition one can resume the section as follows:
Theorem 1. There exists an orthogonal decomposition
$$H=H_<+H_=+H_>$$
such that
1) $`H_<`$ is $`V`$\- surinvariant and the restriction of $`V`$ onto $`H_<`$ is near similar to an uniform contraction;
2) $`H_>`$ is $`V^1`$ \- surinvariant and the restriction of $`V^1`$ onto $`H_<`$ is near similar to an uniform contraction;
3) If $`H_<=\{0\}=H_>`$ then $`V`$ is s-aproximately similar to an unitary;
4) $`H_<\overline{l_2(V)}`$, $`H_>\overline{l_2(V^1)}`$.
Proof. Set $`H_<:=\overline{RanY_0},H_=:=KerX_0KerY_0,H_>:=\overline{RanX_0}`$
and apply the text.
$`\mathrm{}`$
## 4 Stability of $`U`$ in terms of $`Q_t`$
Proposition 1. Let
$$\frac{1}{M}IQ_tMI$$
for some number $`M>0`$ (and every $`t`$). Then $`V`$ is similar to an unitary operator.
Proof. Let $`Q_0`$ be a weak limit point of the net $`Q_t,t+0`$. It exists and
$$\frac{1}{M}IQ_0MI$$
$$Q_0=V^{}Q_0V$$
(see Observation 2.2). Hence, $`V`$ is similar to an unitary operator.
$`\mathrm{}`$
Proposition 2. Let $`V`$ be similar to an unitary operator. Then
$$\frac{1}{M}IQ_tMI$$
for some number $`M>0`$ (and every $`t`$).
Proof. For assumed $`V`$ there exists a number $`M>0`$ such that for every natural $`n`$ there hold
$$V^n^2M,(V^{})^nM$$
Take arbitrary $`xH`$, number $`z`$, $`0<z<1`$ and apply the Observation 2.3:
$$\begin{array}{ccc}\hfill <x,Q_tx>& & (1z)(Vx^2+zV^2x^2+\mathrm{}+z^{n1}V^nx^2+\mathrm{})\hfill \\ & & (1z)\frac{M}{1z}<x,x>=M<x,x>\hfill \end{array}$$
Hence, $`Q_tMI`$.
Now apply the Observation 2.3 to $`Q_t^1`$ and $`V^1`$:
$$\begin{array}{ccc}\hfill <x,Q_t^1x>& & (1z)(V^1x^2+z(V^1)^2x^2+\mathrm{})\hfill \\ & & M<x,x>\hfill \end{array}$$
Hence, $`Q_t^1M`$ and $`\frac{1}{M}Q_t`$.
$`\mathrm{}`$
## 5 When Spectrum has Dichotomy.
Return us to the Example 1.1, which was called motivating . We remarked there, that $`R_0`$ is reminiscent of one of the spectral projector of the operator $`V`$ ( this operator was taken there to be normal ).
Now we will show that somewhat similar situation holds always, especially when the spectrum of the operator $`V`$ does not intersect the unit circle.
Observation 1. $`Y_0=Y_0R_0=R_0Y_0=R_0^{1/2}Y_0R_0^{1/2}`$ In addition $`R_0`$ acts on $`RanY_0`$, hence on $`\overline{l_2(V)}`$, as identity operator
Proof. Recall $`Y_t=tQ_t^1`$, $`Y_0=slim_{t+0}Y_t`$, $`R_t=(I+Q_t)^1`$, $`R_0wlimpt(R_t)`$, all these operators are selfadjoint. What is more, the straightforward calculation shows that
$$Y_t(1R_t)=\frac{t}{I+Q_t}=(1R_t)Y_t.$$
The rest is obvious.
$`\mathrm{}`$
Theorem 1. Suppose there is an $`V`$-invariant subspace, $`L`$ say, such that $`|spectrum(V|L)|<1`$; let $`P`$ denote orthoprojector onto $`L`$ ,
Then
$$LRanY_0^{1/2}Ker(IR_0)$$
$$(IR_t)P0(t+0)$$
Proof.
a) Note $`Ll_2(V)`$ and apply theorem 2.1.
b)
Return to the Observation 2.3:
$$0Q_t(1z)[V^{}V+zV^2V^2+\mathrm{}+z^{n1}V^nV^n]+z^nV^nQ_tV^n$$
Let $`V_P`$ denote $`PV|L`$. Since $`VP=PVP`$ , we can deduce that
$$0PQ_tP(1z)[V_P^{}V_P+zV_P^2V_P^2+\mathrm{}+z^{n1}V_P^nV_P^n]+z^nV_P^nPQ_tPV_P^n$$
Next we adopt the spectrum argument. We make it in the same manner that standard practice suggests: there are some real positive $`ϵ`$, $`M`$ such that
$$r(V_P)+ϵ<1,V_P^nM(r(V)+ϵ)^n(n=1,2,\mathrm{}),$$
Hence
$$z^nV_P^nPQ_tPV_P^n0(n\mathrm{}).$$
and it is routine matter to verify that
$$PQ_tP\frac{(1z)}{1z(r(V_P)+ϵ)^2}M^2(r(V_P)+ϵ)^2.$$
Note that $`M`$ and $`ϵ`$ does not depend on $`z`$ . So, we obtain $`PQ_tP0`$ for $`t+0`$.
( One can show moreover : the serie
$$(1z)[V_P^{}V_P+zV_P^2V_P^2+\mathrm{}+z^{n1}V_P^nV_P^n\mathrm{}]$$
is norm-convergent.)
Now note that
$$PR_tPP(I+Q_t)^1PP(I+PQ_tP)^1P\stackrel{}{}0$$
$$P(IR_t)PPQ_t(I+Q_t)^1PPQ_tP(I+PQ_tP)^1\stackrel{}{}0$$
Recall
$$0IR_tI$$
Hence
$$0P(IR_0)^2P=P(IR_0)^{1/2}(IR_0)(IR_0)^{1/2}PP(IR_0)P$$
So, we can now establish that
$$(IR_t)P^2=P(IR_t)(IR_t)PP(IR_t)P0,$$
It was to be proved.
$`\mathrm{}`$
Corollary. Suppose that the spectrum of the operator $`V`$ does not intersect the unit circle; let $`P`$ denote orthoprojector onto spectral subspace $`L`$ corresponded to the set $`spectrum(V)\{z𝐂||z|<1\}`$ .
Then
$$R_0=P.$$
Proof. By Theorem 1
$$(IR_0)P=0.$$
Hence
$$P=R_0P.$$
Next note that the equation
$$Q=V^{}\frac{Q+t}{I+tQ}V,Q0,0<t1$$
$`(^{})`$
is equivalent to the equation
$$Q^1=V^1\frac{Q^1+t}{I+tQ^1}V^1,Q0,0<t1.$$
$`(^{})`$
(for details see \[Ch1,2\])
For a moment introduce for the (unique) solution of (\*) a longer denotation:
$$Q_t(V).$$
It is straightforward to deduce now that
$$Q_t(V^1)=Q_t(V)^1,R_t(V)=IR_t(V^1),Y_t(V)=X_t(V^1)\mathrm{}\text{etc.}$$
Last recall the Standard Spectrum Theorems (see e.g. \[RS\]) and apply Theorem 1 to the operator $`V^1`$. Then obtain
$$R_0(IP)=0.$$
Hence
$$R_0=R_0P.$$
To complete the proof let compare the second displayed formula with the last one in the current period.
$`\mathrm{}`$
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# Naturalness of the Coleman-Glashow Mass Relation in the 1/𝑁_𝑐 Expansion: an Update
## Abstract
A new measurement of the $`\mathrm{\Xi }^0`$ mass verifies the accuracy of the Coleman-Glashow relation at the level predicted by the $`1/N_c`$ expansion. Values for other baryon isospin mass splittings are updated, and continue to agree with the $`1/N_c`$ hierarchy.
preprint: UCSD/PTH 00-12 JLAB-THY-00-13
The recent measurement of the $`\mathrm{\Xi }^0`$ mass $`1314.82\pm 0.06\pm 0.2`$ MeV by the NA48 collaboration represents a significant improvement over the 30-year-old value $`1314.9\pm 0.6`$ MeV. The $`\mathrm{\Xi }^0`$ mass now is known to an uncertainty comparable to that of the other baryons of the lowest-lying spin-$`1/2`$ octet. This improvement makes it possible to test the precision of the famous Coleman-Glashow (CG) mass relation
$$\mathrm{\Delta }_{\mathrm{CG}}=(pn)(\mathrm{\Sigma }^+\mathrm{\Sigma }^{})+(\mathrm{\Xi }^0\mathrm{\Xi }^{})=0.$$
(1)
Using the old and new experimental values for $`\mathrm{\Xi }^0`$ yields $`\mathrm{\Delta }_{\mathrm{CG}}=0.39\pm 0.61`$ and $`0.29\pm 0.26`$ MeV, respectively: For the first time, $`\mathrm{\Delta }_{\mathrm{CG}}`$ has been measured to have a nonzero value, though only at the one-sigma level. It is of theoretical interest to understand the size of this breaking. In this note, we observe that the experimental value agrees with the theoretical accuracy of the CG relation as predicted in the $`1/N_c`$ expansion of QCD .
The mass spectrum of the baryon spin-$`1/2`$ octet and spin-$`3/2`$ decuplet was analyzed in Ref. in a combined expansion in $`1/N_c`$ and flavor-symmetry breaking. It was found that all of the baryon mass splittings have a natural explanation in terms of powers of $`1/N_c`$, $`SU(3)`$ breaking $`ϵ`$, and isospin breaking $`ϵ^{}`$ (from $`m_dm_u`$) or $`ϵ^{\prime \prime }`$ (from electromagnetic effects). Our analysis differs from the standard flavor-symmetry breaking analysis in that it incorporates the enhanced symmetry of baryons present in the large-$`N_c`$ limit. Large-$`N_c`$ baryons respect an exact $`SU(6)`$ spin-flavor symmetry .For a recent review, see Ref. and references therein. For arbitrary $`N_c`$, the ground state baryons fill the $`N_c`$-quark completely symmetric representation of the spin-flavor algebra, which for $`N_c=3`$ reduces to the usual $`\mathrm{𝟓𝟔}`$-plet of $`SU(6)`$. The spin-flavor symmetry is broken by corrections of subleading order in $`1/N_c`$, while flavor symmetry is broken in the usual manner. Our analysis in Ref. showed that the CG mass combination is $`O(ϵ^{}ϵ/N_c^2)`$ relative to the average mass of the baryon $`\mathrm{𝟓𝟔}`$ spin-flavor multiplet, which is of order $`N_c\mathrm{\Lambda }_{\mathrm{QCD}}`$. For $`N_c=3`$, this result implies that the CG mass combination is predicted to be an order of magnitude smaller than expected from an $`SU(3)`$ flavor symmetry-breaking analysis alone.
In this work, we update the experimental values of mass combinations affected by the new mass measurement of the $`\mathrm{\Xi }^0`$. First, we briefly review notation introduced in Ref. : The isospin $`I`$ combinations of baryon masses are denoted by a subscript $`I`$. Thus, the $`I=0`$ and $`I=1`$ mass combinations of the $`\mathrm{\Xi }^0`$ and $`\mathrm{\Xi }^{}`$ masses are denoted by
$`\mathrm{\Xi }_0`$ $`=`$ $`{\displaystyle \frac{1}{2}}(\mathrm{\Xi }^0+\mathrm{\Xi }^{}),`$ (2)
$`\mathrm{\Xi }_1`$ $`=`$ $`(\mathrm{\Xi }^0\mathrm{\Xi }^{}),`$ (3)
respectively. Using the new value of the $`\mathrm{\Xi }^0`$ mass changes the experimental values of these mass combinations to
$`\mathrm{\Xi }_0`$ $`=`$ $`1318.07\pm 0.12(\mathrm{was}1318.11\pm 0.31)\mathrm{MeV},`$ (4)
$`\mathrm{\Xi }_1`$ $`=`$ $`6.50\pm 0.25(\mathrm{was}6.4\pm 0.6)\mathrm{MeV}.`$ (5)
The improvement in the experimental value of the $`I=0`$ mass combination $`\mathrm{\Xi }_0`$ is small, and does not appreciably affect the numerical evaluation performed in Ref. of $`I=0`$ mass combinations. Thus, we restrict our attention here to $`I=1`$ mass combinations. The remaining $`I=1`$ mass combinations are denoted by
$`N_1=(pn),\mathrm{\Sigma }_1=(\mathrm{\Sigma }^+\mathrm{\Sigma }^{}),\mathrm{\Delta }_1=(3\mathrm{\Delta }^{++}+\mathrm{\Delta }^+\mathrm{\Delta }^03\mathrm{\Delta }^{}),`$ (6)
$`\mathrm{\Sigma }_1^{}=(\mathrm{\Sigma }^+\mathrm{\Sigma }^{}),\mathrm{\Xi }_1^{}=(\mathrm{\Xi }^0\mathrm{\Xi }^{}),`$ (7)
and the $`\mathrm{\Lambda }`$-$`\mathrm{\Sigma }^0`$ mixing parameter. In terms of these definitions, the CG mass combination is given by $`\mathrm{\Delta }_{\mathrm{CG}}=N_1\mathrm{\Sigma }_1+\mathrm{\Xi }_1`$. There are large uncertainties in the $`\mathrm{\Delta }`$ isospin mass splittings, so the $`I=1`$ mass combination $`\mathrm{\Delta }_1`$ does not figure in the mass combinations we consider.
As in Ref. , we define the relative accuracy $`R`$ of a linear combination of masses written in the form $`\mathrm{}r`$ (where the combinations $`\mathrm{}`$ and $`r`$ are uniquely defined to contain baryon masses with only positive coefficients) by $`R|\mathrm{}r|/[(\mathrm{}+r)/2]`$. The quantity $`R`$ yields a scale-independent measure of the breaking of the relation compared to the average baryon mass. The theoretical expectation $`R_T`$ for $`R`$ of a particular mass combination is given by the combined flavor and $`1/N_c`$ suppressions of the mass combination, which are listed in Table II of Ref. , divided by $`N_c`$ since the average baryon mass is $`O(N_c^1)`$. As an example, $`R_T=ϵ^{}ϵ/N_c^2`$ for the CG combination. If the $`1/N_c`$ expansion is natural, then $`R/R_T`$ should be a number of order unity.
In Table I we present values of $`R`$ and $`R_T`$ for the four mass combinations depending upon $`\mathrm{\Xi }_1`$ in Table II of Ref. . We obtain numerical values for $`R_T`$ by taking $`ϵ1/4`$ and $`ϵ^{}1/3^5`$ (and of course $`N_c=3`$); these are typical values one finds for flavor breaking in the meson mass spectrum, but one could also in principle fit to them using the observed baryon masses. One sees first that only the CG combination central value changes substantially from the improvement of the $`\mathrm{\Xi }^0`$ mass measurement, although three of the four uncertainties drop significantly. Most importantly, one observes that the combined $`1/N_c`$ and flavor expansion continues to explain the size of the mass combinations in a natural way. It is also clear from Table I that the agreement of $`R`$ and $`R_T`$ would simply fail without the explicit $`1/N_c`$ factors.
The improvement in the measured value of $`\mathrm{\Xi }_1`$ also permits a better estimate of the $`\mathrm{\Lambda }`$-$`\mathrm{\Sigma }^0`$ mixing parameter. Using Eq. (4.10) of Ref. , we find
$$\mathrm{\Lambda }\mathrm{\Sigma }^0=\frac{1}{2\sqrt{3}}(\mathrm{\Xi }_1N_1)=1.50\pm 0.07(\mathrm{was}1.47\pm 0.17)\mathrm{MeV},$$
(8)
up to a theoretical uncertainty of $`O(ϵ^{}ϵ/N_c^2)`$ times the average baryon mass, which yields a comparable theoretical error.
In summary, the new $`\mathrm{\Xi }^0`$ mass measurement leads to a one-sigma determination of the magnitude of the Coleman-Glashow mass combination. The current experimental value of the CG mass combination is naturally explained in the $`1/N_c`$ expansion, which yields an additional suppression factor of $`1/N_c^2`$ beyond flavor symmetry-breaking factors; an $`SU(3)`$ flavor symmetry-breaking analysis alone fails to explain the observed accuracy of the CG mass combination. Further testing of the mass hierarchy predicted in the combined $`1/N_c`$ and flavor-symmetry breaking expansion is possible by improving the measurements of isospin mass splittings in the decuplet, particularly those of the $`\mathrm{\Delta }`$ baryon.
Acknowledgments
EJ was supported by the Department of Energy under Contract No. DOE-FG03-97ER40546 and by the National Young Investigator program through Grant No. PHY-9457911 from the National Science Foundation. RFL was supported by the Department of Energy under Contract No. DE-AC05-84ER40150.
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# References
December 1999
UTHEP-417
hep-th/0005169
On string theory in $`AdS_3`$ backgrounds<sup>*</sup><sup>*</sup>* Talk given at YITP workshop ‘Developments in Superstring and M-theory’, Kyoto, Japan, October 27-29, 1999.
Yuji Satoh
Institute of Physics, University of Tsukuba
Tsukuba, Ibaraki 305-8571, Japan
ysatoh@het.ph.tsukuba.ac.jp
Abstract
We discuss the string theory on $`AdS_3`$. In the first half of this talk, we review the $`SL(2,R)`$ and the $`SL(2,C)/SU(2)`$ WZW models which describe the strings on the Lorentzian and Euclidean $`AdS_3`$ without RR backgrounds, respectively. An emphasis is put on the fundamental issues such as the unitarity, the modular invariance and the closure of the OPE. In the second half, we discuss some attempts at clarifying such problems. In particular, we discuss the modular invariance of the $`SL(2,R)`$ WZW model and the calculation of the correlation functions of the $`SL(2,C)/SU(2)`$ WZW model using the path-integral approach.
1 Introduction : why string theory on $`AdS_3`$
In this talk, I would like to discuss the string theory on $`AdS_3`$, namely, $`SL(2,R)`$ or its Euclidean analog $`SL(2,C)/SU(2)`$. Besides the recent intensive studies, this string theory has in fact been investigated for more than a decade from various interests -.
First of all, string theory in backgrounds with curved time is not well understood. There are several such models which are relatively well studied , but the analysis is essentially reduced to that of the free theory. In this respect, the string theory on $`SL(2,R)`$ seems to give the simplest truly interacting model. This is because $`SL(2,R)`$ is a very simple space-time with maximal symmetries and the corresponding model is described by the $`SL(2,R)`$ WZW model when there are no RR charges.
Second, related to the above, little is known about non-compact (non-rational) CFTs . Again, the $`SL(2,R)`$ WZW model or its Euclidean analog, $`SL(2,C)/SU(2)`$ WZW model, gives the simplest one.
Third, it is known that the $`SL(2,R)`$ WZW model is closely related to the string models in various black hole backgrounds. For instance, the $`SL(2,R)/U(1)`$ WZW model, which is obtained by a coset, describes the strings in two-dimensional black hole backgrounds . An orbifold of the $`SL(2,R)`$ WZW model gives the string model in the three-dimensional BTZ black hole geometry . When a five-dimensional black hole corresponding to the D1-D5 system is lifted to six dimensions, its near-horizon geometry becomes $`AdS_3\times S^3`$ (precisely speaking (BTZ black hole)$`\times S^3`$). By further taking an S-dual, the system is described by the $`SL(2,R)\times SU(2)`$ WZW . Similarly, $`AdS_3`$ or the BTZ black hole appears quite generally as the near-horizon geometry of the black strings obtained by lifting charged black holes in generic dimension .
Finally, closely related to the above D1-D5 system, the string theory on $`AdS_3`$ gives the simplest case of the AdS/CFT correspondence . This aroused the renewed interest and many works have been devoted to the study of the strings on $`AdS_3`$ in the cases both with and without - RR charges.
However, in spite of recent progress, it seems that there still remain open questions about the string theory on $`AdS_3`$ itself at the fundamental level. Such a state of the problem was recently discussed in . In this talk, we will focus on the cases without RR charges. In the next two sections, we will review the $`SL(2,R)`$ WZW model and its Euclidean analog, the $`SL(2,C)/SU(2)=H_3^+`$ WZW model. We will see that our understanding is still incomplete on the fundamental consistency conditions of string theory such as the unitarity, the modular invariance and the closure of the operator product expansions. Hence we will discuss some attempts towards better understanding in the following sections. In section 4, we will discuss modular invariance of the $`SL(2,R)`$ WZW model and obtain some important information about the spectrum . In section 5, we will discuss the calculation of the correlation functions of the $`H_3^+`$ WZW model using a path-integral approach . We will conclude with a brief summary.
2 $`SL(2,R)`$ WZW model
Let us start with the discussion of Lorentzian $`AdS_3`$. It is defined by the following metric and the embedding equation,
$`ds^2`$ $`=`$ $`dx_0^2dx_1^2+dx_2^2+dx_3^2,`$
$`l^2`$ $`=`$ $`x_0^2x_1^2+x_2^2+x_3^2.`$ (2.1)
This is a maximally symmetric space with negative constant curvature and a solution to the three-dimensional Einstein’s equations with a negative cosmological term $`l^2`$,
$`R_{\mu \nu }`$ $`=`$ $`2l^2g_{\mu \nu }.`$ (2.2)
The space-time defined in the above is the same as the group manifold $`SL(2,R)`$. Hence without RR charges the (bosonic) string theory in this background is described by the $`SL(2,R)`$ WZW model. Its action is given by
$`S`$ $`=`$ $`{\displaystyle \frac{k}{8\pi }}{\displaystyle _\mathrm{\Sigma }}\mathrm{Tr}\left(dgdg^1\right)+{\displaystyle \frac{ik}{12\pi }}{\displaystyle _B}\mathrm{Tr}\left(g^1dg\right)^3,`$ (2.3)
where $`g(z)SL(2,R)`$, $`k`$ is the level, $`\mathrm{\Sigma }`$ is a two-dimensional surface (world-sheet) and $`B`$ is a three-dimensional manifold satisfying $`B=\mathrm{\Sigma }`$. The action has the $`\widehat{sl}(2,R)_L\times \widehat{sl}(2,R)_R`$ current algebra symmetry. The corresponding currents are
$`J(z)={\displaystyle \frac{ik}{2}}gg^1,`$ $`\stackrel{~}{J}(\overline{z})={\displaystyle \frac{ik}{2}}g^1\overline{}g.`$ (2.4)
Here we have denoted the quantity in the right sector by tilde. In the following we will omit the expressions in the right sector unless they are necessary. The model has the conformal symmetry and its energy-momentum tensor is given by the Sugawara form,
$`T(z)`$ $`=`$ $`{\displaystyle \frac{1}{k2}}\eta _{ab}J^a(z)J^b(z),`$ (2.5)
where $`\eta _{ab}=`$ diag $`(1,1,1)`$. $`J^a(z)`$ are defined through $`J(z)=\eta _{ab}\tau ^aJ^b(z)`$ with $`\tau ^asl(2,R)`$. In terms of the modes of $`J^a(z)`$, those of $`T(z)`$ are written as
$`L_n`$ $`=`$ $`{\displaystyle \frac{1}{k2}}{\displaystyle \underset{m𝐙}{}}:{\displaystyle \frac{1}{2}}J_{nm}^+J_m^{}+{\displaystyle \frac{1}{2}}J_{nm}^{}J_m^+J_{nm}^0J_m^0:.`$ (2.6)
These currents and the energy momentum tensor satisfy the following commutation relations
$`[J_n^0,J_m^0]`$ $`=`$ $`{\displaystyle \frac{1}{2}}kn\delta _{n+m},[J_n^0,J_m^\pm ]=\pm J_{n+m}^\pm ,`$
$`[J_n^+,J_m^{}]`$ $`=`$ $`2J_{n+m}^0+kn\delta _{n+m},`$
$`[L_n,J_m^a]`$ $`=`$ $`mJ_{n+m}^a,`$ (2.7)
$`[L_n,L_m]`$ $`=`$ $`(nm)L_{n+m}+{\displaystyle \frac{c}{12}}n(n^21)\delta _{n+m},`$
with $`c`$ being the central charge given by
$`c`$ $`=`$ $`{\displaystyle \frac{3k}{k2}}.`$ (2.8)
Because of the symmetry, the states at the lowest grade are classified by the representations of $`SL(2,R)`$. They are labeled by the values of $`J_0^0`$ and $`\stackrel{}{J}^2=\frac{1}{2}J_0^+J_0^{}+\frac{1}{2}J_0^{}J_0^+J_0^0J_0^0`$. These operators act as
$`\stackrel{}{J}^2|j,m=j(j+1)|j,m,`$ $`J_0^0|j,m=m|j,m.`$ (2.9)
Since $`j(j+1)`$ is invariant under $`jj1`$, one can always bring the values of $`j`$ into the region Re $`j1/2`$ and Im $`j0`$. We will take this convention. A generic state in the left sector is obtained by acting on $`|j,m`$ with $`J_n^a`$ $`(n0)`$ and takes the form
$`\left(J_{n_1}^{a_1}J_{n_2}^{a_2}\mathrm{}\right)|j,m.`$ (2.10)
A generic states in the model is obtained by tensoring (2.10) and a similar expression of the right sector.
Since we expect the model to be unitary, we choose the unitary $`SL(2,R)`$ representations for the zero-mode part. There are five classes of such representations. For the universal covering group of $`SL(2,R)`$, they are
(1) Identity representation $`𝒟_{\mathrm{id}}`$ : the trivial representation with $`\stackrel{}{J}^2=J_0^0=0`$.
(2) Principal continuous series $`𝒟_{\mathrm{pc}}`$: representations with $`m=m_0+n,\mathrm{\hspace{0.17em}0}m_0<1`$, $`n𝐙`$
and $`j=1/2+i\rho ,\rho >0`$.
(3) Supplementary series $`𝒟_{\mathrm{sup}}`$: representations with $`m=m_0+n,\mathrm{\hspace{0.17em}0}m_0<1,n𝐙`$ and
min$`\{m_0,m_01\}<j1/2`$.
(4) Highest weight discrete series $`𝒟_{\mathrm{hw}}`$ : representations with $`m=M_{\mathrm{max}}n`$, $`n=0,1,2,\mathrm{},`$
$`j=M_{\mathrm{max}}1/2`$ and the highest weight state satisfying $`J_0^+|j,j=0`$.
(5) Lowest weight discrete series $`𝒟_{\mathrm{lw}}`$: representations with $`m=M_{\mathrm{min}}+n`$, $`n=0,1,2,\mathrm{},`$
$`j=M_{\mathrm{min}}1/2`$ and the lowest weight state satisfying $`J_0^{}|j,j=0`$.
If we do not take the universal covering group, the parameters are restricted to $`m_0=0,1/2`$ in (2), $`m_0=0`$ in (3) and $`j=`$ (half integers) in (4) and (5).
The harmonic analysis on $`SL(2,R)`$ shows that the square-integrable functions are decomposed into the representations of $`𝒟_{\mathrm{pc}}`$, $`𝒟_{\mathrm{hw}}`$ and $`𝒟_{\mathrm{lw}}`$. Schematically,
$`L^2\left(SL(2,R)\right)`$ $``$ $`{\displaystyle \underset{j<1/2}{}}\left(j{\displaystyle \frac{1}{2}}\right)\left(𝒟_{\mathrm{hw}}^j𝒟_{\mathrm{lw}}^j\right){\displaystyle _0^{\mathrm{}}}𝑑\rho f(\rho )𝒟_{\mathrm{pc}}^{1/2+i\rho },`$ (2.11)
where $`f(\rho )`$ is a certain measure (for details, see, e.g., ).
Ghost problem
Soon after the study of the string theory on $`SL(2,R)`$ was initiated, it turned out that the model contains negative-norm physical states, namely, ghosts . In the flat case, the original model (in the conformal gauge) also contains negative-norm states because of the time direction. However, the physical state conditions $`(L_n\delta _n)|\mathrm{\Psi }=0`$ $`(n0)`$ are sufficient to remove such states. The result in indicates that this does not work in the $`SL(2,R)`$ case. In fact, it is easy to find the ghosts.
To see this, we first note the on-shell condition
$`\left(L_01\right)|\mathrm{\Psi }=0,`$ $`L_0={\displaystyle \frac{j(j+1)}{k2}}+N,`$ (2.12)
where $`N`$ is the grade. This means that the spin at the zero-mode part, $`|j,m`$, should be
$`j`$ $`=`$ $`j(N){\displaystyle \frac{1}{2}}\left(\mathrm{\hspace{0.17em}1}+\sqrt{1+4(k2)(N1)}\right),`$ (2.13)
which corresponds to $`𝒟_{\mathrm{hw}}`$ or $`𝒟_{\mathrm{lw}}`$ for $`k>2`$ and $`N>1`$. Next, we consider a set of states
$`\left\{J_0^a\mathrm{}J_0^b|E_N\right\},`$ $`|E_N=(J_1^+)^N|j(N),j(N).`$ (2.14)
We then find that all the above states are physical but form a non-unitary representation of $`SL(2,R)`$ for a sufficiently large $`N`$. This is because $`|E_N`$ behaves like a highest weight state of an $`SL(2,R)`$ representation with $`j=m=j(N)+N>0`$. Thus we have found the ghosts.
Having found that the model contains ghosts, one might think that the $`SL(2,R)`$ WZW model is sick. However, there are several pieces of evidence that the model should make sense. First of all, in the weak curvature limit the model becomes the flat model and hence one should be able to get a sensible model at least at weak curvature. Second, the authors of studied the particle limit of the model but did not find any pathologies. Third, the effective action of the bosonic $`\sigma `$-model was studied in . There it was found that the effective action has an extremal point corresponding to $`AdS_3`$ and the model is unitary at one-loop. Finally, as discussed in the introduction, the near-horizon geometry of the D1-D5 system is described by the $`SL(2,R)\times SU(2)`$ WZW model after an S-dual transformation. Since the D1-D5 system is unitary (at least at weak coupling), we expect that the $`SL(2,R)\times SU(2)`$ WZW model should also be unitary.
Resolution of the ghost problem
The above argument implies that we might be missing something important and it might be possible to get a sensible theory by finding out an appropriate treatment. There are actually two types of the proposals for the resolution of the ghost problem.
In one proposal , the discrete series $`𝒟_{\mathrm{hw}}`$ and $`𝒟_{\mathrm{lw}}`$ are used and the claim is that if we truncate the spectrum so that the spin and the level are restricted to
$`{\displaystyle \frac{1}{2}}j<{\displaystyle \frac{k}{2}},k>2,`$ (2.15)
one can remove the ghosts. We call this the unitarity bound. This bound seems natural if we recall the argument of the $`SU(2)`$ WZW model. In that case, to maintain the unitarity or the modular invariance, one needs to truncate the $`SU(2)`$ spin so that
$`0j{\displaystyle \frac{k}{2}}.`$ (2.16)
Such a truncation is compatible with the closure of the OPE and the Ward identities. Thus it is completely sensible.
However, in the $`SL(2,R)`$ case, it is not clear if the truncation (2.15) is compatible with other consistency conditions of string theory. This is because such consistency conditions are not well understood either.(Regarding the discussion of the OPE in the Euclidean case, see .) Moreover, from the on-shell condition (2.12), the unitarity bound means the truncation of the string excitation $`N`$. This seems physically unnatural. In addition, the dimensions of the primaries $`L_0=j(j+1)/(k2)`$ are negative for the discrete series when $`k>2`$.
In the other proposal , the principal continuous series $`𝒟_{\mathrm{pc}}`$ is used. One way to understand this argument is to start with a Wakimoto-like representation of $`\widehat{sl}(2,R)`$ using one free boson $`\varphi `$ and the $`\beta `$-$`\gamma `$ system. We then bozonize the $`\beta `$-$`\gamma `$ by two free bosons. One of the points there is that an additional zero-mode is introduced through this bosonization. Here, it may be useful to recall that the primary states $`|j,m`$ have only two zero-modes whereas those in the three-dimensional flat theory have three as $`|p^0,p^1,p^2`$ (though the total zero-modes in the left and the right sectors are three in both cases). Thus it seems natural to incorporate another zero-mode if we expect that the $`SL(2,R)`$ model smoothly leads to the flat model in the weak curvature limit. The added zero-mode turn out to give the sector satisfying
$`{\displaystyle \frac{\gamma }{\gamma }}`$ $``$ $`0,`$ (2.17)
which is called the long string sector . Then by carefully treating the zero-mode part, we find that the on-shell condition picks up the spins $`j=1/2+i\rho `$ which precisely correspond to $`𝒟_{\mathrm{pc}}`$. Finally, using the expressions in terms of the free bosons, the no-ghost theorem is shown similarly to the flat case. In this proposal, the smooth flat limit is achieved by taking $`k\mathrm{}`$. In addition, the applications to the black holes discussed in the introduction appear to be straightforward .
Nevertheless, as in the previous case, it is not clear if this proposal is compatible with the other consistency conditions. (For the discussions of the OPE and the modular invariance in this case, see and respectively.)
In fact, we must say that there is no agreement about how to construct the sensible theory of the $`SL(2,R)`$ strings. Therefore, to clarify this issue it is very important to further investigate the fundamental problems such as the modular invariance, the closure of the OPE, how to choose the spectrum and how to calculate the correlators.
3 $`SL(2,C)/SU(2)`$ WZW model
In the previous section, we discussed the fundamental open questions about the $`SL(2,R)`$ WZW model. Now let us turn to the discussion of the $`SL(2,C)/SU(2)=H_3^+`$ WZW model. The precise formulation of the AdS/CFT correspondence requires Euclidean anti-de Sitter spaces . Euclidean $`AdS_3`$ is called $`H_3^+`$ and given by
$`ds^2`$ $`=`$ $`dx_0^2+dx_1^2+dx_2^2+dx_3^2,`$
$`l^2`$ $`=`$ $`x_0^2+x_1^2+x_2^2+x_3^2.`$ (3.1)
Note the sign-flips compared with the Lorentzian case (2.1). This space is also a maximally symmetric space with negative constant curvature.
To get a string background, one needs to introduce the NS $`B_{\mu \nu }`$ field. In some parametrization, the action takes the form,
$`S`$ $`=`$ $`{\displaystyle \frac{k}{\pi }}{\displaystyle d^2\sigma \left(\varphi \varphi +e^{2\varphi }\overline{\gamma }\overline{}\gamma \right)}.`$ (3.2)
Here $`\overline{\gamma }=\gamma ^{}`$ and $`\varphi +\mathrm{}`$ corresponds to the boundary of $`H_3^+`$. If $`\gamma `$ and $`\overline{\gamma }`$ are independent, the geometry becomes Lorentzian $`AdS_3`$. This action is obtained also by substituting $`g(z)=h(z)h^{}(z)`$ with
$`h`$ $`=`$ $`\left(\begin{array}{cc}1& \gamma \\ 0& 1\end{array}\right)\left(\begin{array}{cc}e^{\varphi /2}& 0\\ 0& e^{\varphi /2}\end{array}\right)SL(2,C)`$ (3.7)
into (2.3) . From this construction, the coset structure $`SL(2,C)/SU(2)`$ is obvious and one finds that the model is actually a WZW model. In terms of $`\varphi `$, $`\gamma `$ and $`\overline{\gamma }`$, the functional measure takes a non-trivial form
$`𝒟\varphi 𝒟(e^\varphi \gamma )𝒟(e^\varphi \overline{\gamma }).`$ (3.8)
The action has the current algebra symmetry $`\widehat{sl}(2,C)\times \widehat{sl}(2,C)^{}`$. In this case, the left and the right symmetries are the complex conjugate to each other. The currents of the global symmetry acting on the zero-mode part are realized by
$`J_0^{}`$ $`=`$ $`_\gamma ,J_0^0=\gamma _\gamma {\displaystyle \frac{1}{2}}_\varphi ,`$
$`J_0^+`$ $`=`$ $`\gamma ^2_\gamma \gamma _\varphi e^{2\varphi }_{\overline{\gamma }}.`$ (3.9)
Note that the last term in the second line. This does not affect the commutation relations but is necessary to assure the invariance of the action.
A convenient way to generate the primary fields is to use the following functionals
$`V^j`$ $`=`$ $`\left[(\gamma x)(\overline{\gamma }\overline{x})e^\varphi +e^\varphi \right]^{2j},`$ (3.10)
where $`x`$ and $`\overline{x}`$ are some parameters. By expanding $`V^j`$ in terms of $`x^{j+m}`$ and $`\overline{x}^{j+\overline{m}}`$, one gets the primary fields $`V_{m,\overline{m}}^j`$ with the definite eigenvalues of $`J_0^0,\overline{J}_0^0`$ and the left and the right Casimirs. It turns out that $`x`$ and $`\overline{x}`$ are interpreted as the coordinates of the boundary CFT .
Similarly to the Lorentzian case, the Hilbert space is decomposed into the representations of $`SL(2,C)`$. Schematically ,
$`L^2(H_3^+)`$ $``$ $`{\displaystyle _0^{\mathrm{}}}𝑑\rho \rho ^2𝒟_{\mathrm{pc}}^{1/2+i\rho }.`$ (3.11)
We remark that only the principal continuous series $`𝒟_{\mathrm{pc}}`$ appear and there are no discrete series.<sup>1</sup><sup>1</sup>1In section 2, we considered the representations of $`SL(2,R)`$. Here we are considering the corresponding representations of $`SL(2,C)`$. In this case, the spectrum for $`j=1/2+i\rho `$ is given by $`m=(ip+n)/2,\overline{m}=(ipn)/2`$ with $`p𝐑,n𝐙`$.
Furthermore, by (i) introducing auxiliary fields $`\beta `$ and $`\overline{\beta }`$, (ii) taking into account the non-trivial measure (3.8) and (iii) rescaling $`\varphi `$, one obtains the following action,
$`S`$ $`=`$ $`{\displaystyle \frac{1}{2\pi }}{\displaystyle d^2\sigma \left(\varphi \varphi +\beta \overline{}\gamma +\overline{\beta }\overline{\gamma }\beta \overline{\beta }e^{2\varphi /\alpha _+}\frac{2}{\alpha _+}\varphi \sqrt{\widehat{g}}\widehat{R}\right)},`$ (3.12)
where $`\alpha _+=\sqrt{2(k2)}`$ and $`\widehat{g}`$ and $`\widehat{R}`$ are the background metric and curvature of the world-sheet, respectively. In this expression, the interaction term $`\beta \overline{\beta }e^{2\varphi /\alpha _+}`$ drops out in the limit $`\varphi \mathrm{}`$. Thus we obtain a free theory in that limit, namely, near the boundary of $`H_3^+`$. The last term in $`J_0^+`$ in (3.9) also drops out and we get the free-field expression in that limit.
Puzzles
This WZW model has been intensively studied recently - in relation to the AdS/CFT correspondence. We may need to be careful about to what extent and how the string theory without RR charges is relevant to the AdS/CFT correspondence (see, for example, ). However, if the correspondence is naively taken, one finds some puzzles. They are summarized in the following table of the correspondence,
| string (WZW model) | supergravity | CFT |
| --- | --- | --- |
| discrete series (non-normalizable) | KK mode | chiral primary |
| continuous series (normalizable) | ?? | ?? |
Namely, although the Hilbert space of the $`H_3^+`$ WZW model consists of the principal continuous series, we do not find the corresponding objects on the supergravity and the CFT sides. Moreover, following the argument in , the scaling dimension of a boundary CFT operator and the $`sl_2`$ spin of the corresponding operator of the $`H_3^+`$ WZW model are related by $`h=j`$. If this is valid also for the principal continuous series, the dimension of the corresponding boundary CFT operator becomes complex. Thus it is hard to interpret the correspondence for the continuous series even if it exists.
These puzzles might not lead to an immediate contradiction because, as discussed in , the $`H_3^+`$ WZW model is a non-compact CFT and hence there might not be the state-operator correspondence as in the Liouville theory. Nevertheless, in order to complete the correspondence, it seems necessary to further investigate these puzzles. To this end, we may need to study the $`H_3^+`$ WZW model in detail. Again, the fundamental consistency conditions play the role of the guideline there.
How are then the precise discussions of the $`H_3^+`$ WZW model possible? One way is to use the free field approximation. By this approach, we can get much information but this is valid only near the boundary $`\varphi \mathrm{}`$. Another way is to use the generating functionals of the primary fields (3.10) following . In this approach, the full analysis beyond the free field approximation is possible but it tends to be semi-classical (see, however, regarding the full quantum analysis based on the bootstrap). Therefore it would be nice to have a description beyond the free field or the semi-classical treatment. We will return to this point later.
4 Modular invariance
In the preceding sections, we put an emphasis on the importance of the further investigations of the fundamental problems. Here we would like to discuss some attempts at clarifying the modular invariance of the $`SL(2,R)`$ WZW model. Although the modular invariance in the non-compact case is not well understood, there are several arguments in the $`SL(2,R)`$ case. For example, the modular invariants are discussed in by using the $`\widehat{sl}(2,R)`$ characters based on the discrete series $`𝒟_{\mathrm{hw}}`$ and $`𝒟_{\mathrm{lw}}`$ and by incorporating some new sectors corresponding to winding modes. They are also discussed in using the characters for $`𝒟_{\mathrm{pc}}`$ along the line of .
In this section, we will focus on the possibility of constructing the modular invariants from the characters for $`𝒟_{\mathrm{hw}}`$ and $`𝒟_{\mathrm{lw}}`$ without incorporating any additional sectors as in . For details, see . This issue is also discussed in .
Let us start with the definition of the characters. For the current algebras based on compact Lie groups, the characters are naturally defined using three variables. With this in mind, we define the characters for the discrete series by
$`\mathrm{ch}_j(z,\tau ,u)`$ $``$ $`e^{2\pi iku}{\displaystyle e^{2\pi iJ_0^0z}e^{2\pi i\tau (L_0\frac{c}{24})}}.`$ (4.1)
The summation is taken over the entire module of the current algebra representations. The plus sign in the first factor $`e^{+2\pi iku}`$ is due to the change $`kk`$ compared with the compact case. To calculate these characters, one needs to know about singular vectors. For a generic highest (or lowest) weight representation, which is not necessarily $`𝒟_{\mathrm{hw}}`$ or $`𝒟_{\mathrm{lw}}`$, the current module has singular vectors when one of the following conditions is satisfied :
$`(1)`$ $`2j+1=s+(k2)(r1),`$
$`(2)`$ $`2j+1=sr(k2),`$ (4.2)
$`(3)`$ $`k2=0,`$
where $`r,s`$ are positive integers.
The characters $`\mathrm{ch}_j`$ in generic cases seem unknown. However, when there are no singular vectors, they are given by
$`\chi _\mu ^{\mathrm{hw}}(z,\tau ,u)`$ $`=`$ $`e^{2\pi iku}e^{2\pi i\mu z}q^{\frac{\mu ^2}{k2}}i\vartheta _1^1(z|\tau ),`$ (4.3)
for $`𝒟_{\mathrm{hw}}`$ and $`\chi _\mu ^{\mathrm{lw}}(z,\tau ,u)=\chi _\mu ^{\mathrm{hw}}(z,\tau ,u)`$ for $`𝒟_{\mathrm{lw}}`$. Here $`q=e^{2\pi i\tau }`$, $`\mu j+1/2`$ and
$`\vartheta _1(z|\tau )`$ $`=`$ $`2q^{1/8}\mathrm{sin}(\pi z){\displaystyle \underset{n=1}{\overset{\mathrm{}}{}}}(1q^n)(1q^ne^{2\pi iz})(1q^ne^{2\pi iz}).`$ (4.4)
Since $`\chi _\mu ^{\mathrm{lw}}(z,\tau ,u)=\chi _\mu ^{\mathrm{hw}}(z,\tau ,u)`$, $`\chi _\mu ^{\mathrm{hw}}`$ with $`\mu 0`$ $`(j1/2)`$ are regarded as $`\chi _\mu ^{\mathrm{lw}}`$ with $`\mu 0`$. Thus we will use only $`\chi _\mu ^{\mathrm{hw}}`$ and drop the superscript hw. We remark that one cannot consider the specialized characters $`\chi _\mu (0,\tau ,0)`$ since they diverge in the limit $`z0`$ because of the infinite degeneracy with respect to $`L_0`$.
In our normalization of $`(z,\tau ,u)`$, the modular transformations are generated by
$`S:`$ $`(z,\tau ,u)({\displaystyle \frac{z}{\tau }},{\displaystyle \frac{1}{\tau }},u+{\displaystyle \frac{z^2}{4\tau }}),`$
$`T:`$ $`(z,\tau ,u)(z,\tau +1,u).`$ (4.5)
Under $`T`$-transformation, the characters just get phases,
$`\chi _\mu (z,\tau +1,u)`$ $`=`$ $`e^{2\pi i\left(\frac{\mu ^2}{k2}+\frac{1}{8}\right)}\chi _\mu (z,\tau ,u).`$ (4.6)
For $`k2<0`$, the $`S`$-transformation of $`\chi _\mu (z,\tau ,0)`$ is given in . In our case with three variables, it reads as
$`\chi _\mu ({\displaystyle \frac{z}{\tau }},{\displaystyle \frac{1}{\tau }},u+{\displaystyle \frac{z^2}{4\tau }})`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{2k}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\nu e^{4\pi i\frac{\mu \nu }{k2}}\chi _\nu (z,\tau ,u).`$ (4.7)
For $`k2>0`$, the right-hand side of (4.7) does not converge on the upper half plane of $`\tau `$. Instead, after some calculation, we get a slightly different result,
$`\chi _\mu ({\displaystyle \frac{z}{\tau }},{\displaystyle \frac{1}{\tau }},u+{\displaystyle \frac{z^2}{4\tau }})`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{k2}}}{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\nu e^{4\pi \frac{\mu \nu }{k2}}\chi _{i\nu }(z,\tau ,u).`$ (4.8)
Note that an imaginary $`\mu =j+1/2`$ corresponds to a spin of the principal continuous series but $`\chi _{i\mu }`$ are not the characters for those representations.
As a simpler case, we will first discuss the possibility of constructing modular invariants using finite number of the discrete series characters. Given the explicit forms of $`\chi _\mu `$, we can then show that it is impossible to make modular invariants from finite number of the discrete series characters without singular vectors, i.e., from $`\chi _\mu `$. Similarly, since the characters with singular vectors are obtained by subtracting states from $`\chi _\mu `$, the above statement is extended to some cases including singular vectors. In fact, we can show that, for $`k>2`$, it is impossible to construct modular invariants from finite number of the characters based on either $`𝒟_{\mathrm{hw}}`$ or $`𝒟_{\mathrm{lw}}`$. The arguments are simple applications of Cardy’s for $`c>1`$ CFT and we will omit them. To further extend the latter statement to the cases including both $`𝒟_{\mathrm{hw}}`$ and $`𝒟_{\mathrm{lw}}`$, the explicit forms of the characters with singular vectors seem to be necessary. In addition, we notice that a similar statement does not hold for $`k<2`$. To see this, we note that the arguments do not use any special properties of the discrete unitary series and hence it is the same as for a generic highest (or lowest) weight $`\widehat{sl}(2,R)`$ representations. However, when $`k<2`$, modular invariants using finite number of the characters are actually known for the so-called admissible representations .
Next, let us move on to the case in which infinitely many characters are allowed. For the time being, we will discuss the modular invariants using $`\chi _\mu `$ only. In such a case, using their modular properties we can show that it is impossible to construct modular invariants only from $`\chi _\mu `$ with $`\mu `$ belonging to a finite interval $`\mu [\mu _1,\mu _2]`$ even if infinitely many $`\chi _\mu `$ are used. This might seem obvious from the $`S`$-transformation of $`\chi _\mu `$ since the right-hand sides of (4.7) and (4.8) does not close within $`\chi _\mu `$ with $`\mu [\mu _1,\mu _2]`$. However, we need to be careful because we are considering an infinite dimensional space of the characters $`\chi _\mu `$. For instance, it is not clear which $`\chi _\mu `$ are independent and whether or not the expressions (4.7) and (4.8) are unique. In fact, it may be possible to get different expressions by deforming the integration contours in (4.7) and (4.8). In any case, the detailed argument is given in .
Since, for $`k>2`$, $`\chi _\mu `$ become divergent for $`\mathrm{Im}\tau >0`$ as $`|\mu |\mathrm{}`$, the above statement means that for $`k>2`$ it is impossible to construct modular invariants only from $`\chi _\mu `$. Thus the possibility of constructing modular invariants from $`\chi _\mu `$ is limited to the case where $`k<2`$ and $`\chi _\mu `$ with $`|\mu |\mathrm{}`$ are included. In this case, we can actually construct a modular invariant,
$`Z_{\mathrm{diag}}(z,\tau ,u)`$ $`=`$ $`{\displaystyle _{\mathrm{}}^{\mathrm{}}}𝑑\mu |\chi _\mu |^2={\displaystyle _{\mathrm{}}^0}𝑑\mu \left(|\chi _\mu ^{\mathrm{hw}}|^2+|\chi _\mu ^{\mathrm{lw}}|^2\right)`$ (4.9)
$`=`$ $`{\displaystyle \frac{1}{2}}e^{4\pi k\mathrm{Im}u}e^{(2k)\pi \frac{(\mathrm{Im}z)^2}{\mathrm{Im}\tau }}\sqrt{{\displaystyle \frac{2k}{\mathrm{Im}\tau }}}|\vartheta ^2(z|\tau )|.`$
The diagonal partition function with $`u=0`$, i.e., $`Z_{\mathrm{diag}}(z,\tau ,0)`$, was discussed in . In our case, it is straightforward to check that $`Z_{\mathrm{diag}}(z,\tau ,u)`$ is modular invariant owing to the presence of $`u`$. Although it may be interpreted as a kind of a twisted partition function, its physical meaning is still unclear (recall that we cannot set $`u=z=0`$).
As pointed out also in , $`Z_{\mathrm{diag}}(z,\tau ,0)`$ was discussed in in the context of a path-integral approach to the $`H_3^+`$ WZW model. Since this model has the $`\widehat{sl}(2,C)\times \widehat{sl}(2,C)^{}`$ symmetry, the diagonal partition function may be understood also as the partition function of this model. However, in different spectrum seems to be summed up. It is interesting to consider the precise relationship between the approach here and the one in .
In order to discuss a generic case including the characters with singular vectors, we may again need the explicit forms of such characters. Nevertheless, it turns out that the case without singular vectors covers physically interesting cases and gives important implication to the unitarity bound (2.15). This is because the condition of the singular vectors (4.2) implies that there are no singular vectors within (2.15). Furthermore, since the spins $`j`$ in that bound belong to a finite interval, our results indicate that one cannot construct modular invariants only from the discrete series characters based on the representations satisfying the unitarity bound (2.15). This means that one cannot make a consistent string theory on $`SL(2,R)=AdS_3`$ only from the spectrum within (2.15). This was already discussed in , but we believe that at least we have refined the argument a little.
Since there exist ghosts for the discrete series outside the unitarity bound, simply adding such spectrum may not give a consistent theory. Therefore, the possibilities for a consistent theory seem (a) to use the discrete series satisfying (2.15) but include some new sectors with different characters from $`\chi _\mu `$ as in , and/or (b) to use the spectrum of other representations as in . To settle down this problem, further investigations are necessary.
5 Correlation functions
In the previous section, we discussed the modular invariance and saw that it gives an important information about the spectrum of the strings on $`SL(2,R)`$. Finally in this section, we will discuss the calculation of the correlations functions of the $`H_3^+`$ WZW model . This is important not only by itself but also for studying the OPE and hence the spectrum of the model.
The outstanding feature of the $`H_3^+`$ WZW model is that it allows us the Lagrangian approach . It is alternative and complementary to the current algebra approach which is often used. Actually in the Lagrangian approach we may be able to get a description beyond the free field and the semi-classical approximations.
To see this, let us first recall the action (3.2) and the functional measure (3.8). Surprisingly, with (3.2) and (3.8) it is possible to carry out the path-integrals for some correlators . For this purpose, we need (i) the ‘partition function’ obtained after integrating out $`\gamma `$ and $`\overline{\gamma }`$,
$`\mathrm{exp}\left[S(\varphi )\right]`$ $`=`$ $`\mathrm{exp}\left[{\displaystyle \frac{1}{\pi }}{\displaystyle d^2\sigma \left((k2)\varphi \varphi \frac{1}{4}\varphi \sqrt{\widehat{g}}\widehat{R}\right)}\right],`$ (5.1)
and (ii) the propagator,
$`\gamma (z)\overline{\gamma }(\overline{w})`$ $`=`$ $`{\displaystyle \frac{1}{k\pi }}{\displaystyle d^2y\frac{e^{2\varphi (y)}}{(zy)(\overline{w}\overline{y})}}.`$ (5.2)
The ‘partition function’ implies that the resulting effective theory of $`\varphi `$ is a free theory with a background charge. Since the the propagator contains the factor $`e^{2\varphi (y)}`$, it plays a similar role to the screening charges in the free field approach. Using (5.1) and (5.2), one can calculate the correlators of the form
$`{\displaystyle \underset{i}{}}e^{a_i\varphi (z_i)}\gamma ^{b_i}(z_i)\overline{\gamma }^{c_i}(\overline{z}_i).`$ (5.3)
We are interested in the correlation functions of the primary fields in the discrete series. In the free field approximation, the primaries are given by
$`V_{m,\overline{m}}^{j(\mathrm{free})}`$ $`=`$ $`e^{2j\varphi }\gamma ^{j+m}\overline{\gamma }^{j+\overline{m}}.`$ (5.4)
In the full theory, they are obtained by expanding the functionals (3.10) and take the form
$`V_{m,\overline{m}}^j`$ $`=`$ $`{\displaystyle \underset{j^{},m^{},\overline{m}^{}}{}}C_{m^{},\overline{m}^{}}^j^{}V_{m^{},\overline{m}^{}}^{j^{}(\mathrm{free})},`$ (5.5)
where $`C_{m^{},\overline{m}^{}}^j^{}`$ are some coefficients.
We would like to calculate the correlators among the above primary fields. A similar calculation was actually carried out in the case of the finite dimensional representations . In our infinite dimensional case, it seems that we need to carefully choose the ‘conjugate’ fields paired with the above primaries. Once we have obtained the correlation functions, we can extract important information about the $`H_3^+`$ WZW model and in turn this gives useful insights into the AdS/CFT correspondence. We would like to report progress in this direction elsewhere .
6 Summary
In this talk, we first noted that the string theory on $`AdS_3`$ is important in various respects besides in relation to the AdS/CFT correspondence. We then reviewed the $`SL(2,R)`$ and the $`H_3^+`$ WZW models which describe the string propagations on Lorentzian and Euclidean $`AdS_3`$, respectively. We saw that in spite of the recent intensive studies there still remain open questions at the fundamental level. An emphasis was put on the importance of further investigating the fundamental problems such as the modular invariance, the closure of the OPE, the issue of the spectrum and the calculation of the correlation functions. With this in mind, we discussed some attempts at clarifying such problems. First, we discussed the modular invariance of the $`SL(2,R)`$ WZW model and showed that it gives important information of the spectrum. Next, we discussed the correlation functions of the $`H_3^+`$ WZW model using the full Lagrangian approach.
Although our attempts were quite incomplete, some of the problems seem still tractable. Thus further investigations will lead us to a deeper understanding of the string theory on $`AdS_3`$. We expect that such investigations also shed some light on the AdS/CFT correspondence.
Note added
Some of the questions raised in this talk have been discussed also in recent papers .
Acknowledgements
I would like to thank N. Ishibashi, A. Kato and K. Okuyama for the collaborations on the subjects discussed here. I would also like to thank I. Bars, J. de Boer, A. Giveon, H. Ishikawa, K. Ito, M. Kato, P.M. Petropoulos and S.-K. Yang for useful discussions and correspondences. Finally, I am grateful to the organizers of the workshop ‘Developments in Superstring and M-theory’ held at YITP, Kyoto, 27-29 October, 1999, for giving a chance to deepen my understanding on the string theory on $`AdS_3`$.
References
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# 1 Two Reggeon fusion diagram
The role of secondary Reggeons in central meson production
N.I.Kochelev<sup>a,</sup><sup>1</sup><sup>1</sup>1e-mail address: kochelev@thsun1.jinr.ru, T.Morii<sup>b,</sup><sup>2</sup><sup>2</sup>2e-mail address: morii@kobe-u.ac.jp, B.L.Reznik<sup>c,</sup><sup>3</sup><sup>3</sup>3e-mail address: reznik@dvgu.ru, A.V.Vinnikov<sup>a,c,</sup><sup>4</sup><sup>4</sup>4e-mail address: vinnikov@thsun1.jinr.ru
<sup>a</sup> Bogoliubov Laboratory of Theoretical Physics,
Joint Institute for Nuclear Research,
Dubna, Moscow region, 141980 Russia
<sup>b</sup> Faculty of Human Development, Division of Sciences for Natural Environment
and Graduate School of Science and Technology,
Kobe University, Nada, Kobe 657-8501, Japan
<sup>c</sup> Far Eastern State University, Sukhanova 8, GSP, Vladivostok, 690600 Russia
## Abstract
We estimate the contribution of $`f_2`$ trajectory exchange to the central $`\eta `$ and $`\eta ^{}`$ production. It is shown that secondary Reggeons may give a large contribution to processes of double diffractive meson production at high energies.
PACS number(s): 12.40.Nn, 13.60.Le, 12.39.Mk
The Regge theory provides a natural and economical description of hadron-hadron interactions at high energies and small momentum transfers . In this approach the interaction between colliding particles is described by the exchange of effective particles, i.e. Reggeons. At high energies, the pomeron with vacuum quantum numbers gives the dominant contribution to the hadron-hadron total cross sections. The Reggeons with quantum numbers different from the vacuum ones can also contribute to the total cross sections and their contribution is vanishing at very high energies.
An interest in double diffractive processes (DDP) is mediated by the idea that they can be a pure source of glueballs , . Intensive studies of these processes have been recently performed by WA102 collaboration at CERN. The mechanism of the central meson production in DDP at high energies is usually related to the double pomeron exchange (DPE) , , . This conclusion is based on the following observations:
i) if the two-pomeron fusion contributes dominantly to the central meson production, one can explain rather weak energy dependence of the total production cross sections;
ii) $`t`$\- and azimuthal dependences of differential cross sections for the most mesons are consistent with the two vector-fusion mechanism and with additional assumption that pomeron transforms as vector current;
iii) quantum numbers of these effective vector exchanges, e.g. flavour, P- and C-parities, are the same as for the pomeron.
However, these arguments fail if one looks into the details of the experimental data. For example, even in the simplest case of the light pseudoscalar meson production we find:
i) the observed cross section of $`\eta `$ production is larger than that of the $`\eta ^{}`$ production , while the DPE mechanism predicts the opposite: $`\sigma _\eta <<\sigma _\eta ^{}`$. This conclusion comes from the consideration of the flavour-singlet structure of the two-pomeron fusion which leads to the enhancement of flavour-singlet meson production .
ii) the cross section of $`\pi ^0`$ meson production does not show any angular dependence in the range $`0^{}<\phi <150^{}`$ and shows a peak (presumably a contribution from the diffractive $`\mathrm{\Delta }`$ and nucleon resonance production ) at $`\phi =180^{}`$, while the mechanism of two-vector fusion predicts the behaviour like $`\mathrm{sin}^2\phi `$ with a maximum at $`90^{}`$. Hence, the mechanism of $`\pi ^0`$ production is not consistent with the two-vector fusion at all.
The main goal of this letter is to underline the importance of secondary Reggeon exchanges for central meson production. As an example, we estimate the contribution of the $`f_2`$ exchange with $`P=C=+1`$ to the central production of $`\eta `$ and $`\eta ^{}`$ mesons for WA102 kinematics.
Let us consider a two-Reggeon fusion contribution to central $`\eta `$, $`\eta ^{}`$ productions, as is presented in Fig. 1, where the pomeron and $`f_2`$ Reggeon are taken into consideration.
The cross section of meson production in the reaction
$$p(p_1)+p(p_2)p(p_1^{})+p(p_2^{})+M(p_M)$$
is given by the formula
$$d\sigma =\frac{dPS^3}{4\sqrt{(p_1.p_2)^2m_p^4}}\underset{spin}{}|T|^2,$$
(1)
where $`dPS^3`$ is the 3-body phase space volume, $`m_p`$ is the proton mass, and $`T`$ is the matrix element for DDP reaction. $`_{spin}`$ stands for the summation and averaging over the spins in the final and initial proton states, respectively.
At high energies and small momentum transfers, the four-momenta of initial and final protons in the center-of-mass system are given as follows:
$`p_1`$ $``$ $`(P+m_p^2/2P,\stackrel{}{0},P),p_2(P+m_p^2/2P,\stackrel{}{0},P),`$
$`p_{1}^{}{}_{}{}^{}`$ $``$ $`(x_1P+(m_p^2+\stackrel{}{p}_{1T}^{}{}_{}{}^{2})/2x_1P,\stackrel{}{p}_{1T},x_1P),`$ (2)
$`p_{2}^{}{}_{}{}^{}`$ $``$ $`(x_2P+\stackrel{}{p}_{2T}^{}{}_{}{}^{2})/2x_2P,\stackrel{}{p}_{2T},x_2P),`$
where $`P=\sqrt{s}/2`$, $`s=(p_1+p_2)^2`$. Using the result of Ref. for the high energy phase space at small momentum transfers, we obtain
$$dPS^3=\frac{1}{2^8\pi ^4}dt_1dt_2dx_1dx_2d\mathrm{\Phi }\delta (s(1x_1)(1x_2)M^2)$$
(3)
for the phase space volume in the DDP reaction, where $`\mathrm{\Phi }`$ is azimuthal angle between final protons, $`t_{1,2}=(p_{1,2}p_{1,2}^{})^2`$ and $`M`$ is the meson mass. Kinematic limits for the phase space integration in (13) are determined by positive $`\stackrel{}{p}_{1,2T}^{}{}_{}{}^{2}`$ in (12) and the condition $`s_{1,2}(M+m_p)^2`$, where
$$s_{1,2}=s(1x_{1,2})+m_p^2+2t_{1,2}.$$
(4)
Let us consider typical values of $`s_{1,2}`$ for the diffractive process, which appeared to be very important to understand the reaction mechanism. In the diffractive region where $`t_{1,2}`$ are small $`s_{1,2}`$ at given $`s`$ can be functions of $`x_{1,2}`$ only. One can see from (3) that $`x_{1,2}`$ are not independent variables. At fixed $`x_F=x_2x_1`$
$$x_1=1\sqrt{\frac{x_F^2}{4}+\frac{M^2}{s}}\frac{x_F}{2},x_2=1\sqrt{\frac{x_F^2}{4}+\frac{M^2}{s}}+\frac{x_F}{2}.$$
(5)
Using (4) and (5) we obtain the dependence of $`s_{1,2}`$ on $`x_F`$ (see Fig. 2).
As one can see, despite the large value of the total invariant mass of the system, the value of $`s_2`$ is always rather small. Therefore, the large contribution from secondary Regge trajectories can be expected for any central meson production. Fixing $`x_F`$ and varying $`s`$ (see Fig. 3), we see that this conclusion remains valid even at very high energies. This is why even at LHC energy the contribution from secondary Reggeon exchanges can be not small, contrary to the expectation .
Let us estimate more accurately the contribution of both pomeron and secondary Reggeons to the central $`\eta ,\eta ^{}`$ production. The leading correction to the two-pomeron fusion contribution to this reaction comes from the additional pomeron-$`f_2`$ Reggeon fusion, where $`f_2`$ is the trajectory with P=C=+1 and positive signature.
The matrix element of DPE reaction, $`T`$, is given by the formula:
$$T=i36\beta _P^2\lambda _MA_{PP}^Mϵ_{\mu \nu \rho \sigma }p_1^\mu p_2^\nu p_1^\rho p_2^\sigma F_p(t_1)F_p(t_2)F_{PPM}(t_1,t_2),$$
(6)
where
$$F_p=\frac{4m_p^22.79t}{(4m_p^2t)(1t/0.71)^2},$$
(7)
and
$$F_{PPM}(t_1,t_2)=\frac{1}{(1t_1/8\pi ^2f_{PS}^2)(1t_2/8\pi ^2f_{PS}^2)}$$
(8)
are form factors in proton-pomeron and pomeron-pomeron-pseudoscalar vertices, respectively (see ), and
$$A_{PP}^M=\beta _P^4D_{PP}^M\left(\frac{s_1}{s_0}\right)^{\alpha _P(t_1)1}\left(\frac{s_2}{s_0}\right)^{\alpha _P(t_2)1}\mathrm{exp}\left(\frac{i\pi }{2}[\alpha _P(t_1)+\alpha _P(t_2)2]\right),$$
(9)
$`s_0=1`$ GeV<sup>2</sup>, $`\beta _P=1.8`$ GeV<sup>-1</sup>, $`\alpha _P(t)=1+ϵ+\alpha ^{}t`$ is the pomeron trajectory with $`ϵ0.08,\alpha ^{}0.25`$ GeV<sup>-2</sup> and
$$\lambda _M=\frac{18}{R_M\alpha _{em}}\sqrt{\frac{2\mathrm{\Gamma }_{\gamma \gamma }}{\pi M^3}}.$$
(10)
Here $`\mathrm{\Gamma }_{\eta \gamma \gamma }=0.46\times 10^6`$ GeV, $`\mathrm{\Gamma }_{\eta \gamma \gamma }=4.28\times 10^6`$ GeV and factors $`D_{PP}^M`$ and $`R^M`$ are related to the wave functions of $`\eta `$ and $`\eta `$
$$\eta =\mathrm{sin}\mathrm{\Theta }\eta _0+\mathrm{cos}\mathrm{\Theta }\eta _8,\eta ^{}=\mathrm{cos}\mathrm{\Theta }\eta _0+\mathrm{sin}\mathrm{\Theta }\eta _8,$$
(11)
$$D_{PP}^\eta =\mathrm{sin}\mathrm{\Theta },D_{PP}^\eta ^{}=\mathrm{cos}\mathrm{\Theta },$$
$$R_\eta =2\sqrt{2}\mathrm{cos}\mathrm{\Theta }\mathrm{sin}\mathrm{\Theta },R_\eta ^{}=2\sqrt{2}\mathrm{cos}\mathrm{\Theta }+\mathrm{sin}\mathrm{\Theta },$$
where $`\mathrm{\Theta }=19.5^{}`$ is the singlet-octet mixing angle.
Using (3), (5), (6) and the equations
$$\stackrel{}{p}_{1,2T}^{}{}_{}{}^{2}=x_{1,2}t_{1,2}(1x_1)^2m_p^2,$$
(12)
which follow from (2), the cross section is finally written as
$`{\displaystyle \frac{d\sigma }{dt_1dt_2dx_Fd\mathrm{\Phi }}}={\displaystyle \frac{3^4\lambda _M^2F_p^2(t_1)F_p^2(t_2)F_{PPM}^2(t_1,t_2)}{2^9\pi ^4\sqrt{x_F^2+4M^2/s}}}`$
$`\times (x_1t_1+(1x_1)^2m_p^2)(x_2t_2+(1x_2)^2m_p^2)|A_{PP}^M|^2\mathrm{sin}^2\mathrm{\Phi }.`$ (13)
The $`f_2`$ Reggeon gives an additional contribution to the total amplitude
$$A^M=A_{PP}^M+A_{Pf_2}^M+A_{f_2P}^M,$$
(14)
where
$`A_{Pf_2}^M=\beta _P^2\beta _{f_2}^2D_{Pf_2}^M\left({\displaystyle \frac{s_1}{s_0}}\right)^{\alpha _P(t_1)1}\left({\displaystyle \frac{s_2}{s_0}}\right)^{\alpha _{f_2}(t_2)1}\mathrm{exp}\left({\displaystyle \frac{i\pi }{2}}[\alpha _P(t_1)+\alpha _{f_2}(t_2)2]\right),`$
$`A_{f_2P}^M=\beta _P^2\beta _{f_2}^2D_{Pf_2}^M\left({\displaystyle \frac{s_1}{s_0}}\right)^{\alpha _{f_2}(t_1)1}\left({\displaystyle \frac{s_2}{s_0}}\right)^{\alpha _P(t_2)1}\mathrm{exp}\left({\displaystyle \frac{i\pi }{2}}[\alpha _{f_2}(t_1)+\alpha _P(t_2)2]\right).`$
The factors $`D_{ij}^M`$ can be obtained from the quark decomposition of $`\eta _1`$, $`\eta _8`$ and $`f_2`$ mesons,
$`\eta _1={\displaystyle \frac{1}{\sqrt{3}}}(u\overline{u}+d\overline{d}+s\overline{s}),`$
$`\eta _8={\displaystyle \frac{1}{\sqrt{6}}}(u\overline{u}+d\overline{d}2s\overline{s}),`$ (15)
$`f_2={\displaystyle \frac{1}{\sqrt{2}}}(u\overline{u}+d\overline{d}),`$
where we assume that $`f_2`$ is an ideal mixing of $`SU(3)`$ flavour octet and singlet.
We have
$$D_{Pf_2}^\eta =\sqrt{\frac{2}{3}}\mathrm{sin}\mathrm{\Theta }+\frac{1}{\sqrt{3}}\mathrm{cos}\mathrm{\Theta },D_{Pf_2}^\eta =\sqrt{\frac{2}{3}}\mathrm{cos}\mathrm{\Theta }+\frac{1}{\sqrt{3}}\mathrm{sin}\mathrm{\Theta }.$$
(16)
The parameters of $`f_2`$ trajectory have been taken from Donnachie-Landshoff fit
$$\beta _{f_2}=3.6\text{GeV}^1,ϵ_{f_2}=0.45,$$
(17)
with $`\alpha _{f_2}^{}0.9`$ GeV<sup>-2</sup>.
The final result for $`\eta `$ and $`\eta `$ production cross sections for WA102 kinematics ($`P=450`$ GeV/c, $`0<x_F<0.1`$) becomes to be equal to
$$\sigma (\eta )=450\text{nb},\sigma (\eta )=550\text{nb}.$$
(18)
which can be compared with the DPE contribution alone :
$$\sigma (\eta )=46\text{nb},\sigma (\eta )=370\text{nb},$$
(19)
and with the experimental data
$$\sigma (\eta )=1295\pm 120\text{nb},\sigma (\eta )=588\pm 60\text{nb}.$$
(20)
We see that $`f_2`$ Reggeon contribution to the cross section of $`\eta `$ meson production is very large. The admixture of flavour non-singlet $`f_2`$ exchange increases the cross section by an order of magnitude and also gives the large enhancement of $`\eta ^{}`$ production. Taking into account uncertainties of our model for possible values of quark-pomeron and quark-$`f_2`$ coupling constants as well as representation of form factors in effective vertices, we conclude that the sum of DPE and $`f_2`$ exchange can explain the observed value of total cross section of $`\eta `$ and $`\eta `$ central production. On the other hand, the DPE alone fails to explain the cross section value. We should stress that $`f_2`$ exchange does not spoil a good DPE description of the $`t`$ and azimuthal dependence of the differential cross sections of $`\eta `$ and $`\eta ^{}`$ production , which is just fixed by the shape of nucleon and pseudoscalar meson form factors and vector-like structure of quark-Reggeon vertex.
We have also analysed the DPE contribution to the cross section for different values of $`x_F`$ and $`s`$ (see Figs. 4 and 5). As is seen here, even for $`\eta `$ production, its contribution is not dominant at any value of $`x_F`$ and $`s`$. Increasing energy does not lead to the prevailing of DPE in $`\eta `$ and $`\eta `$ production. This is the reason why meson production within DDP does not seem to be a pure kinematical region for the dominance of DPE mechanism even at LHC energy .
Although our calculation has been performed only for the simplest case of pseudoscalar $`\eta `$ and $`\eta ^{}`$ meson production, we think that any DDP cannot leave the region where exchanges by secondary Regge trajectories are significant. This conclusion is based mainly on kinematical arguments; this is why we expect them to be correct for any central meson production too. We should also mention that due to large admixture of flavour singlet component, $`f_2`$ exchange can play an important role even in reactions of central glueball production.
Concluding, we have estimated the contribution of secondary Reggeon trajectories into central production of $`\eta `$ and $`\eta ^{}`$ mesons. The contribution is shown to be very large. Therefore, before to make some definite conclusions about properties of the pomeron from the analysis of central production data, one should carefully disentangle the secondary Reggeon contribution.
We are grateful to A. Dorokhov, S. Gerasimov, V. Romanovsky, N. Russakovich and A. Titov for useful discussions
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# 1 Introduction
## 1 Introduction
The potential of a static quark<sup>1</sup><sup>1</sup>1 A static quark (anti-quark) is an external source in the (complex conjugate of the) fundamental representation of the gauge group. anti-quark pair in a Yang-Mills gauge theory has been computed by Monte Carlo simulations on the lattice. The observable from which the static potential is extracted is the Wilson loop. The results for gauge group SU(2) \[?,?\] and SU(3) \[?,?\], close to the continuum limit, show a linearly rising confinement potential at large separations of the static charges. When the Yang-Mills gauge theories are coupled to matter fields in the fundamental representation of the gauge group, the static potential is expected to flatten at large distances: the ground state of the system is better interpreted in terms of two weakly interacting static-light mesons which are bound states of a static and a dynamical quark. The dynamical quarks are pair-created in the strong gauge field binding the static quarks. This phenomenon is called string breaking or screening of the static charges. The name “string” refers to the gauge field configuration which confines the static quarks and leads to the linear confinement in pure gauge theories.
In recent attempts in QCD with two flavors of dynamical quarks \[?,?,?,?,?,?,?\], the flattening in the static potential determined from the Wilson loops was not visible. The string breaking distance $`r_\mathrm{b}`$ around which the static potential should start flattening off, could nevertheless be estimated in the quenched approximation of QCD to be \[?,?\]
$`r_\mathrm{b}`$ $``$ $`2.7r_0,`$ (1)
where $`r_00.5\mathrm{fm}`$ is a scale conveniently defined from the force, $`F(r)`$, between static quarks\[?\]:
$`F(r_0)r_0^2`$ $`=`$ $`1.65r_0,`$ (2)
The QCD results so far show a linear rise of the static potential for distances beyond $`r_\mathrm{b}`$. It has been commonly argued that the problem is the poor overlap of the Wilson loops with the ground state of the system. The investigation of the static potential in models other than QCD is therefore relevant in order to understand its origin and identify the reason for the failure of the method used to extract it in full QCD.
First studies of string breaking were performed with a hopping-parameter expansion in SU(2) gauge theory with Wilson fermions \[?\]. In the Schwinger model ($`\mathrm{QED}_2`$), the exact solution for the static potential can be given in the limit of zero fermion mass \[?\]: $`V(r)=(e\sqrt{\pi }/2)\{1\mathrm{exp}(er/\sqrt{\pi })\}`$, where $`e`$ is the charge of the static sources. String breaking was established by numerical simulation in the Schwinger model \[?,?\]. Numerical evidence of the screening of the static potential was also found in the U(1) Higgs model (scalar QED) in two dimensions \[?\]. The flattening of the static potential at large distances is also expected in the confinement “phase” of the SU(2) Higgs model. Indeed, early simulations yielded some qualitative evidence for string breaking \[?,?\].
String breaking can also be studied in Yang-Mills theories using static sources in the adjoint representation of the gauge group. The gauge field itself is responsible for the screening of the sources and the formation of hadrons called “gluelumps”. An important numerical investigation concerning this screening has been carried out by C. Michael in \[?\], where it has been noted that string breaking can be a mixing phenomenon. The static potential is extracted from a matrix correlation in which two types of states enter, the adjoint string and the “two-gluelump”. However, due to large errors, no clear evidence for string breaking could be given. The first numerical evidence for string breaking in non-Abelian gauge theories with dynamical matter fields was given using the mixing method in the four-dimensional \[?\] and three-dimensional \[?\] SU(2) Higgs model by the computation of the potential between static quarks. Most recently, the extraction of the static adjoint potential in the three-dimensional \[?,?\] and four-dimensional \[?\] SU(2) Yang-Mills theory shows also evidence for string breaking.
In full QCD, string breaking has been seen at finite temperature \[?\], where the static potential can be extracted from Polyakov loop correlators. A recent investigation in zero-temperature QCD with two flavors of dynamical quarks \[?\] gives some indication that string breaking can be observed in the potential determined from a matrix correlation containing string-type and meson-type states. The main problems are the computational costs of the light quark propagators entering the matrix correlation: the maximal variance reduction method of \[?\] has been applied in \[?\].
In this article, we present results for the spectrum of static-light mesons and for the static potential in the confinement “phase” of the SU(2) Higgs model. The method of computation of the static potential is the same as in our first work \[?\] but it is explained here in much more detail. In Sect. 2, we present a detailed investigation of the spectrum of the static-light mesons. This is relevant not only for determining the asymptotic value of the static potential, but also for finding a suitable smearing procedure for the Higgs field to be used in the computation of the static potential. In Sect. 3, we show our results for the ground state and the first excited static potential obtained at $`\beta =2.4`$, with a spatial resolution two times better than at $`\beta =2.2`$ used in \[?\].
In order to study the interpretation of string breaking as a mixing phenomenon \[?\] we first properly define overlaps of the string and two-meson states with the ground and first excited energy eigenstates. Their dependence on the distance $`r`$ then establishes string breaking as a level crossing phenomenon. In Sect. 4, we study the scaling of the static potentials by comparing the results at $`\beta =2.4`$ with $`\beta =2.2`$ on a line of constant physics \[?\]. The dependence of the static potentials on the value of the Higgs quartic self-coupling is also investigated.
## 2 The spectrum of the static-light mesons
The investigation of the spectrum of static-light mesons is an important study to be done before the extraction of the static potential. The system composed by a static quark anti-quark pair is expected to be described at large separation of the static sources by two weakly interacting static-light mesons, which are bound states of a static quark and the dynamical matter field. Denoting by $`\mu `$ the mass of the lowest meson state, the static potential $`V_0(r)`$ is expected to approach the value
$`\underset{r\mathrm{}}{lim}V_0(r)=2\mu .`$ (3)
Therefore, the mass $`\mu `$ basically determines the string breaking distance $`r_\mathrm{b}`$ around which the potential starts flattening out.
The extraction of the static-light meson spectrum is representative for the variational method that we employ also for the extraction of the static potential. We constructed a large basis of operators that create one-meson “states” when applied to the vacuum. The variational approach chooses the best linear combinations of these “states” which approximate the energy eigenstates. Through the determination of the meson spectrum we can therefore gain information about a suitable way of constructing the meson-type states. Then, we use this information in Sect. 3 to construct a basis of states for the determination of the static potentials.
### 2.1 The matrix correlation
In \[?\], an Hamiltonian formalism for the SU(2) Higgs model is constructed. Along the lines of \[?\], a transfer matrix operator is defined and its strict positivity for $`\kappa >0`$ and $`\lambda >0`$ is proved<sup>2</sup><sup>2</sup>2 One can show that the partition function of the SU(2) Higgs model satisfies the property $`Z(\beta ,\kappa ,\lambda )=Z(\beta ,\kappa ,\lambda )`$. There is a mapping of the observables such that the expectation values at positive $`\kappa `$ are reproduced by expectation values at negative $`\kappa `$. This motivates the restriction of the parameter region to the values $`\kappa >0`$. For $`\lambda =0`$, strict positivity of the transfer matrix holds for $`0<\kappa <1/6`$. . This property is equivalent to the reality of the energy spectrum in any sector of the Hilbert space. Different charge sectors of the Hilbert space are defined by the transformation property of the states under gauge transformation. Through Gauss’ law, this gauge transformation property is related to the presence of static charges in some irreducible representation of the gauge group.
The static-light meson states belong to the charge sector with one static charge in the fundamental representation of the gauge group localised at a certain space position $`\stackrel{}{x}`$. The meson states $`|i,(i=1,2,3,\mathrm{}),`$ are described by (composite) fields $`O_i^\mathrm{M}(x)`$, which are constructed with field variables taken at equal time $`x_0`$ and transform under a gauge transformation $`\{\mathrm{\Lambda }(x)\mathrm{SU}(2)\}`$ according to
$`[O_i^{\mathrm{M},\mathrm{\Lambda }}(x)]_a`$ $`=`$ $`\mathrm{\Lambda }_{aa^{}}^{}(x)[O_i^\mathrm{M}(x)]_a^{},`$ (4)
where $`a,a^{}=1,2`$ are color indices. The obvious choice is to take $`O^\mathrm{M}(x)`$ to be the Higgs field $`\mathrm{\Phi }(x)`$. We are also going to consider non-local linear combinations which take into account contributions from Higgs fields at neighboring sites (smeared fields) and more general composite fields, with the intent to model the true wave function of the meson. From the basis of meson-type fields $`O_i^\mathrm{M}(x)`$ a matrix correlation
$`C_{ij}^\mathrm{M}(t)`$ $`=`$ $`[O_j^\mathrm{M}(x+t\widehat{0})^{}]_aU(x,x+t\widehat{0})_{ab}^{}[O_i^\mathrm{M}(x)]_b.`$ (5)
is constructed representing the transition amplitude over a time interval $`t`$ from the meson state $`i`$ to the meson state $`j`$. The static charge is represented by a straight time-like Wilson line $`U(x,x+t\widehat{0})^{}`$ connecting $`x`$ with $`x+t\widehat{0}`$. With the help of the reconstruction theorem proved in \[?\], it is possible to show that in the limit of an infinite physical time extension $`T`$ of the lattice<sup>3</sup><sup>3</sup>3 In practice, the limit $`T\mathrm{}`$ is reached when $`Tm_\mathrm{H}1`$, where $`m_\mathrm{H}`$ is the Higgs mass defined as the mass gap in the zero charge (gauge invariant) sector of the Hilbert space. the correlation matrix eq. (5) can be written like
$`C_{ij}^\mathrm{M}(t)`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}j|\alpha \alpha |i\mathrm{e}^{tW_\alpha },`$ (6)
where $`|\alpha ,(\alpha =0,1,2,\mathrm{})`$ are the orthonormal meson energy-eigenstates with energies<sup>4</sup><sup>4</sup>4 We normalise the vacuum energy to be 0. $`W_\alpha `$, $`W_\alpha <W_{\alpha +1}`$. The matrix correlation eq. (5) can be measured in a Monte Carlo simulation: now, we describe the variational method for extracting the meson energy spectrum from it.
### 2.2 Variational method
For matrices of the type in eq. (6) a general lemma for the extraction of the energies $`W_\alpha `$ has been proved in \[?\]. In this reference, a variational method is proposed, which is superior to a straightforward application of the lemma. It consists in solving the generalised eigenvalue problem:
$`{\displaystyle \underset{j}{}}C_{ij}(t)v_{\alpha ,j}(t,t_0)`$ $`=`$ $`\lambda _\alpha (t,t_0){\displaystyle \underset{j}{}}C_{ij}(t_0)v_{\alpha ,j}(t,t_0),\lambda _\alpha >\lambda _{\alpha +1},`$ (7)
where $`t_0`$ is fixed and small (in practice we use $`t_0=0`$). The generalised eigenvalues $`\lambda _\alpha (t,t_0)`$ are computed as the eigenvalues of $`\overline{C}=C(t_0)^{1/2}C(t)C(t_0)^{1/2}`$ and the vectors
$`\overline{v}_{\alpha ,i}={\displaystyle \underset{j}{}}[C(t_0)^{1/2}]_{ij}v_{\alpha ,j}(t,t_0)`$ with $`{\displaystyle \underset{i}{}}\overline{v}_{\alpha ,i}\overline{v}_{\alpha ^{},i}=\delta _{\alpha \alpha ^{}}`$ (8)
are the orthonormal eigenvectors of $`\overline{C}`$. The positivity of the transfer matrix ensures that $`C(t)`$ is positive definite for all $`t`$. In \[?\] it is proven that the energies $`W_\alpha `$ are given by the expressions
$`aW_\alpha `$ $`=`$ $`\mathrm{ln}(\lambda _\alpha (ta,t_0)/\lambda _\alpha (t,t_0))+\mathrm{O}\left(\mathrm{e}^{t\mathrm{\Delta }W_\alpha }\right),`$ (9)
where $`\mathrm{\Delta }W_\alpha =\underset{\beta \alpha }{\mathrm{min}}|W_\alpha W_\beta |`$. It is expected that, for a good basis of states, the coefficients of the higher exponential corrections in eq. (9) are suppressed so that the energies can be read off at moderately large values of $`t`$ from the right-hand side of eq. (9).
The variational method eq. (7) and eq. (9) is our standard method for extracting the energy spectrum. What we have stated here about this method is valid for any charge sector of the Hilbert space. One has to start from a basis $`|i`$ of states belonging to that charge sector. The matrix correlation $`C_{ij}(t)`$ corresponds to matrix elements $`j|\mathrm{TT}^n|i,nt/a,`$ of powers of the transfer matrix operator $`\mathrm{TT}`$ appropriate for the charge sector. In \[?\], this correspondence is derived in detail for the sector with a static quark anti-quark pair: the energy spectrum are the static potential and its excitations.
### 2.3 Smeared Higgs fields
We studied different bases of meson-type fields $`O_i^\mathrm{M}(x)`$ by measuring in Monte Carlo simulations the matrix correlation function $`C_{ij}(t)`$ defined in eq. (5) and computing from it the energy spectrum of the static-light mesons using the variational method described in Sect. 2.2. Our aim was to find the best field basis for describing the ground state of the static-light mesons. For these studies we simulated the SU(2) Higgs model on a $`20^4`$ lattice with parameters $`\beta =2.2`$, $`\kappa =0.274`$ and $`\lambda =0.5`$. This parameter point is in the confinement “phase” of the model and is the point that we used in our first work \[?\]. The measurement of the matrix correlation is improved by the use of the one-link integral method \[?\].
We first studied a basis containing the fundamental Higgs field $`\mathrm{\Phi }(x)`$ and smeared Higgs fields obtained by iterating the application of a smearing operator $`S_1`$ to the Higgs field. The smearing operator $`S_1`$ is defined as
$`S_1\mathrm{\Phi }(x)`$ $`=`$ $`\mathrm{\Phi }(x)+{\displaystyle \underset{\genfrac{}{}{0pt}{}{|xy|=a}{x_0=y_0}}{}}U(x,y)\mathrm{\Phi }(y),`$ (10)
where $`U(x,y)`$ is the link connecting $`y`$ with $`x`$. Iterating the smearing operator $`S_1`$ we obtain smeared Higgs fields $`\mathrm{\Phi }_1^{(m)}(x)=S_1^m\mathrm{\Phi }(x),(m=0,1,2,\mathrm{}),`$ with different smearing levels $`m`$ ($`m=0`$ corresponds to the fundamental Higgs field in the Lagrangian). We measured a matrix correlation function with a basis of smeared Higgs fields corresponding to smearing levels 0,1 and 2 of $`S_1`$. The result for the ground state extracted according to eq. (9) is shown in Fig. 1 (triangles). We were not able to reach a plateau for the ratio $`\mathrm{ln}(\lambda _\alpha (ta)/\lambda _\alpha (t))`$ within the range of $`t`$ considered (up to 8 in lattice unit).
Then, we investigated a larger basis of meson-type fields, defining in particular a smearing operator $`S_2`$ as
$`S_2\mathrm{\Phi }(x)`$ $`=`$ $`𝒫\{𝒫\mathrm{\Phi }(x)+𝒫{\displaystyle \underset{\genfrac{}{}{0pt}{}{|xy|=\sqrt{2}a}{x_0=y_0}}{}}\overline{U}(x,y)\mathrm{\Phi }(y)+`$ (11)
$`𝒫{\displaystyle \underset{\genfrac{}{}{0pt}{}{|xy|=\sqrt{3}a}{x_0=y_0}}{}}\overline{U}(x,y)\mathrm{\Phi }(y)\},`$
where $`𝒫\mathrm{\Phi }=\mathrm{\Phi }/\sqrt{\mathrm{\Phi }^{}\mathrm{\Phi }}`$ and $`\overline{U}(x,y)`$ represents the average over the shortest link connections between $`y`$ and $`x`$. Through iteration of $`S_2`$ we obtain the smeared Higgs fields $`\mathrm{\Phi }_2^{(m)}(x)=S_2^m\mathrm{\Phi }(x),(m=0,1,2,\mathrm{})`$. We considered the following basis of meson-type fields $`O_i^\mathrm{M}(x),i=1,2,\mathrm{},11`$:
$`O_1^\mathrm{M}(x)`$ $`=`$ $`𝒫\mathrm{\Phi }(x),`$ (12)
$`O_2^\mathrm{M}(x)`$ $`=`$ $`𝒫{\displaystyle \underset{\genfrac{}{}{0pt}{}{|xy|=a}{x_0=y_0}}{}}U(x,y)\mathrm{\Phi }(y),`$ (13)
$`O_3^\mathrm{M}(x)`$ $`=`$ $`𝒫{\displaystyle \underset{\genfrac{}{}{0pt}{}{|xy|=\sqrt{2}a}{x_0=y_0}}{}}\overline{U}(x,y)\mathrm{\Phi }(y),`$ (14)
$`O_4^\mathrm{M}(x)`$ $`=`$ $`𝒫{\displaystyle \underset{\genfrac{}{}{0pt}{}{|xy|=\sqrt{3}a}{x_0=y_0}}{}}\overline{U}(x,y)\mathrm{\Phi }(y),`$ (15)
$`O_i^\mathrm{M}(x)`$ $`=`$ $`\mathrm{\Phi }_2^{(i4)}(x),i=5,6,7,8,`$ (16)
$`O_9^\mathrm{M}(x)`$ $`=`$ $`\mathrm{\Phi }(x)\times {\displaystyle \frac{1}{6}}{\displaystyle \underset{k=1}{\overset{3}{}}}\{\mathrm{\Phi }^{}(xa\widehat{k})U(xa\widehat{k},k)\mathrm{\Phi }(x)+`$ (17)
$`\mathrm{\Phi }^{}(x)U(x,k)\mathrm{\Phi }(x+a\widehat{k})\},`$
$`O_{10}^\mathrm{M}(x)`$ $`=`$ $`\mathrm{\Phi }(x)\times {\displaystyle \frac{1}{12}}{\displaystyle \underset{1k<l3}{}}\{P_{kl}(x)+P_{kl}(xa\widehat{k})+`$ (18)
$`P_{kl}(xa\widehat{k}a\widehat{l})+P_{kl}(xa\widehat{l})\},`$
$`O_{11}^\mathrm{M}(x)`$ $`=`$ $`\mathrm{\Phi }(x)\times (\mathrm{\Phi }^{}(x)\mathrm{\Phi }(x)),`$ (19)
where in eq. (18) $`P_{kl}(x)=U(x,k)U(x+a\widehat{k},l)U^{}(x+a\widehat{l},k)U^{}(x,l)`$ and we use the same notation conventions as in \[?\]. In Fig. 2, the result for the extraction of the mass of a static-light meson using the fields $`O_i^\mathrm{M}(x),i=1,\mathrm{},11`$ is shown (triangles). Note the enlarged scale on the y-axis as compared to Fig. 1. We obtain a nice plateau already at moderately large values of $`t`$. The situation remains practically unchanged (also the statistical errors) if we remove from the basis all fields except the smeared fields obtained by iterations of the smearing operator $`S_2`$. This means that this smearing procedure contains all relevant features for describing the ground state which could be obtained by using the larger basis.
When the generalised eigenvalue problem eq. (7) is solved, the optimal linear combination of the basis fields $`O_i^\mathrm{M}(x)`$ describing the ground state can be expressed in terms of the components of the vector $`v_0`$ as $`_iv_{0,i}O_i^\mathrm{M}(x)`$. Therefore, we call $`v_0`$ the ground state wave function. An interesting fact we can learn from $`v_0`$, is that the field $`O_2^\mathrm{M}`$, with nearest neighbor contributions, has a very small coefficient $`v_{0,2}`$. This explains our original difficulties in extracting the meson ground state. In Fig. 1, a direct comparison of the smearing operators $`S_1`$ and $`S_2`$ shows clearly that the contributions from the excited states are much more suppressed when we use $`S_2`$ (circles).
### 2.4 The meson spectrum at $`\beta =\mathbf{2.4}`$
In Fig. 3, we show the results for the static-light meson spectrum that we obtained for the parameters $`\beta =2.4,\kappa =0.2759,\lambda =0.7`$ (in the confinement “phase”) on a $`32^4`$ lattice. More details about this simulation will be given in Sect. 3. For the measurement of the matrix correlation function we used a basis with the six fields
$`\mathrm{\Phi }_2^{(m)}(x),m=1,3,5,7,10,15,`$ (20)
obtained by iterating the smearing procedure $`S_2`$ in eq. (11). As we will see in Sect. 3, the lattice spacing at $`\beta =2.4`$ is reduced by almost a factor two with respect to the lattice spacing at $`\beta =2.2`$. Therefore, at $`\beta =2.4`$ smeared fields with high smearing levels $`m`$ are expected to play a more important role than at $`\beta =2.2`$. This expectation is confirmed by the simulation. In order to determine with confidence the static-light meson masses, we plot in Fig. 4 the logarithmic ratios on the right-hand side of eq. (9) as functions of the correction terms $`\mathrm{exp}(t\mathrm{\Delta }W)`$. This enables us to choose the best time $`t`$ for reading off the masses from the logarithmic ratios and to estimate the systematic errors associated with this choice. For the mass of the ground state, we must take the largest value $`t/a=9`$. For the mass of the first excited state, we can take $`t/a=8`$. In both cases, the systematic errors<sup>5</sup><sup>5</sup>5 The systematic errors for the masses are estimated from the difference between the mass read off at the chosen value of $`t`$ and the crossing point of the dotted lines in Fig. 4 with the y-axis ($`t=\mathrm{}`$). are of the same magnitude as the statistical errors. However, these errors are small. The results for the meson spectrum are<sup>6</sup><sup>6</sup>6 We use the notation $`\mu `$, $`\mu ^{}`$ and $`\mu ^{}`$ for $`W_0`$, $`W_1`$ and $`W_2`$ respectively.
$`a\mu =\mathrm{\hspace{0.33em}0.517}(2),`$ $`a\mu ^{}=\mathrm{\hspace{0.33em}0.88}(3),`$ $`a\mu ^{}=\mathrm{\hspace{0.33em}1.21}(9).`$ (21)
We note that the convergence of the right-hand side of eq. (9) is not so “critical” in the case of the static potentials considered in Sect. 3.
## 3 String breaking and mixing
As mentioned in the introduction, the basic point concerning the determination of the static potential has been first noted by C. Michael \[?\] in a study of the SU(2) static adjoint potential. The energy-eigenstates of the system composed of a static quark anti-quark pair and of light dynamical matter fields are well described by a superposition of string-type and meson-type states. Using a suitable basis of such states a matrix correlation can be constructed from which the static potential and its excitations are extracted for arbitrary separations of the static quarks.
### 3.1 The matrix correlation
The static potentials are defined as the energy levels in the charge sector of the Hilbert space with a static quark at space position $`\stackrel{}{x}`$ and a static anti-quark at space position $`\stackrel{}{x}_r=\stackrel{}{x}+r\widehat{k}`$. States living in this charge sector are described by fields $`O_{ab}(x,x_r)`$ with color indices $`a,b`$ and equal times $`x_0=y_0`$, which transform under gauge transformations like \[?,?\]
$`O_{ab}^\mathrm{\Lambda }(x,x_r)`$ $`=`$ $`\mathrm{\Lambda }_{aa^{}}^{}(x)O_{a^{}b^{}}(x,x_r)\mathrm{\Lambda }_{b^{}b}(x_r).`$ (22)
The simplest choice of such fields describing string-type states is $`U_{ab}(x,x_r)`$ and for the meson-type states $`\mathrm{\Phi }_a(x)\mathrm{\Phi }_b^{}(x_r)`$. By $`U(x,y)`$ we denote the product of gauge links along the straight line connecting $`y`$ with $`x`$. The basic matrix correlation for the extraction of the static potentials can be expressed in terms of the following transition matrix elements \[?\]
$`C_{\mathrm{WW}}(r,t)`$ $`=`$ $`\text{tr}[U(x,x_r)U(x_r,x_r+t\widehat{0})U^{}(x+t\widehat{0},x_r+t\widehat{0})U^{}(x,x+t\widehat{0})],`$ (23)
$`C_{\mathrm{WM}}(r,t)`$ $`=`$ $`\mathrm{\Phi }^{}(x+t\widehat{0})U^{}(x,x+t\widehat{0})U(x,x_r)U(x_r,x_r+t\widehat{0})\mathrm{\Phi }(x_r+t\widehat{0}),`$ (24)
$`C_{\mathrm{MM}}(r,t)`$ $`=`$ $`\mathrm{\Phi }^{}(x+t\widehat{0})U^{}(x,x+t\widehat{0})\mathrm{\Phi }(x)\mathrm{\Phi }^{}(x_r)U(x_r,x_r+t\widehat{0})\mathrm{\Phi }(x_r+t\widehat{0}).`$ (25)
The static quark (anti-quark) is represented by a straight time-like Wilson line $`U^{}(x,x+t\widehat{0})`$ ($`U(x_r,x_r+t\widehat{0})`$). The matrix $`C`$ is real, symmetric and positive. This simplest choice of the states does not however correspond to the physical picture that we have of the system.
The string-type states should reproduce the flux tube \[?,?,?,?,?,?\] of the gauge field binding the static quarks. Therefore, we smear the space-like links describing the string-type states using the APE smearing procedure of \[?\] with smearing strength set to the numerical value $`ϵ=1/4`$. In order to construct “physical” two-meson states, we determine the spectrum of the static-light mesons as described in Sect. 2. In the matrix correlation eq. (5) we use smeared Higgs fields $`O_i^\mathrm{M}(x)=\mathrm{\Phi }_2^{(n_i)}(x)`$ defined with the smearing operator of eq. (11). The numbers $`n_i(i=1,2,\mathrm{},N)`$ denote the smearing levels. The eigenvectors $`v_\alpha \mathrm{IR}^N(\alpha =0,1,2,\mathrm{})`$, obtained by solving the generalised eigenvalue problem eq. (7) for large $`t`$, are the wave functions describing approximately (because of the finite basis of fields and the finite time $`t`$) the true eigenstates of the Hamiltonian. We define the fields
$`\mathrm{\Psi }_\alpha (x)`$ $`=`$ $`{\displaystyle \underset{i=1}{\overset{N}{}}}v_{\alpha ,i}\mathrm{\Phi }_2^{(n_i)}(x)(\alpha =0,1,2,\mathrm{}),`$ (26)
corresponding to the approximate one-meson eigenstates. The fields we choose to describe two-meson states are defined as
$`\left[\mathrm{\Psi }_\alpha (x)\right]_a\left[\mathrm{\Psi }_\beta ^{}(x_r)\right]_b,\alpha ,\beta =\mathrm{\hspace{0.17em}0},1,2.`$ (27)
The values $`\alpha =0,1,2`$ refer to the ground, first and second excited one-meson state. The field basis in eq. (27) contains combinations with $`\alpha \beta `$ which are not symmetric under interchange of the positions $`x`$ and $`x_r`$ of the static charges. Because we expect the ground two-meson state to be symmetric, we project into the symmetric linear combinations of the fields in eq. (27) when we analyse the data of the simulations. The “mixed” states (for example of one meson in the ground state and one meson in the first excited state) can be important when looking at the asymptotic behavior (in $`r`$) of excited static potentials \[?\]. The one-meson states have a space extension due to the smearing of the Higgs field. For a high number of smearing iterations, there is effectively an “interaction” between the mesons in the two-meson states eq. (27) due to the overlap of the smeared Higgs fields
Summarising, we use the basis of states $`|i`$ described by the fields
$`[O_i(x,x_r)]_{ab}=\{\begin{array}{cc}U_{ab}^{(m_i)}(x,x_r)& i=1,2,\mathrm{},N_\mathrm{U}\hfill \\ \left[\mathrm{\Psi }_{\alpha _i}(x)\right]_a\left[\mathrm{\Psi }_{\beta _i}^{}(x_r)\right]_b& i=N_\mathrm{U}+1,\mathrm{},N_\mathrm{U}+9\hfill \end{array}`$ (30)
where $`U^{(m_i)}(x,x_r)`$ is the product of smeared gauge links (with smearing level $`m_i`$) along the straight line connecting $`x_r`$ with $`x`$ and the pairs of indices $`(\alpha _i=0,1,2;\beta _i=0,1,2)`$ label the 9 combinations of two-meson states. We construct the following matrix correlation
$`C_{ij}(t,r)=[O_i(x,x_r)]_{ab}U_{bc}(x_r,x_r+t\widehat{0})[O_j(x+t\widehat{0},x_r+t\widehat{0})]_{cd}^{}U_{da}^{}(x,x+t\widehat{0}).`$ (31)
We denote the static potentials by $`V_\alpha (r),\alpha =0,1,2,\mathrm{}`$. The corresponding eigenstates of the Hamiltonian are denoted by $`|\alpha `$. Taking the limit of infinite time extension of the lattice $`T\mathrm{}`$, we obtain from the transfer matrix formalism the following spectral representation of eq. (31) \[?\]
$`C_{ij}(t,r)`$ $`=`$ $`{\displaystyle \underset{\alpha }{}}j|\alpha \alpha |i\mathrm{e}^{tV_\alpha (r)}.`$ (32)
For fixed separation $`r`$, we extract from $`C(t,r)`$ the potentials $`V_\alpha (r)`$ using the variational method described in Sect. 2.2.
### 3.2 Results at $`\beta =\mathbf{2.4}`$
In our first study \[?\] we obtained the static potential from a simulation at $`\beta =2.2,\kappa =0.274,\lambda =0.5`$ on a $`20^4`$ lattice. We observed string breaking at a distance $`r_b/a\mathrm{\hspace{0.17em}5}`$. We decided then to study the system with a better lattice resolution at $`\beta =2.4`$. The results that we describe in the following are obtained on a $`32^4`$ lattice for the parameter set
$`\beta =\mathrm{\hspace{0.17em}2.4},\kappa =\mathrm{\hspace{0.17em}0.2759},\lambda =\mathrm{\hspace{0.17em}0.7}.`$ (33)
The field basis is constructed according to eq. (26) and eq. (30) from smeared gauge ($`N_\mathrm{U}=3`$) and Higgs fields, whose smearing parameters are summarised in Table 1. The parameters for the simulation were fixed after some trial runs. In the matrix correlation function $`C_{ij}(t,r)(i,j=1,2,\mathrm{},12)`$, the time-like links are replaced by their one-link integrals \[?\]. We collected a statistics of 800 measurements. Autocorrelations in the measurements are practically absent: the statistical errors, computed by a jackknife analysis remain constant when we group the measurements in bins of length 1,2 or 4.
#### 3.2.1 Renormalised static potentials
The static potentials $`V_\alpha (r)(\alpha =0,1,2,\mathrm{})`$ are extracted from the matrix correlation eq. (31) using the variational method described in Sect. 2.2. We rewrite eq. (9) as
$`aV_\alpha (r)`$ $`=`$ $`\mathrm{ln}(\lambda _\alpha (ta,t_0)/\lambda _\alpha (t,t_0))+\mathrm{O}\left(\mathrm{e}^{t\mathrm{\Delta }V_\alpha (r)}\right),`$ (34)
where $`\mathrm{\Delta }V_\alpha (r)=\underset{\beta \alpha }{\mathrm{min}}|V_\alpha (r)V_\beta (r)|`$ and the eigenvalues $`\lambda _\alpha (t,t_0)`$ are obtained by solving the generalised eigenvalue problem eq. (7) with the matrix correlation function at fixed $`r`$. We choose $`t_0=0`$. At all distances $`r`$ we can read off with confidence and very good statistical precision (per mille level) values for the static potential $`V_0(r)`$ at $`t=7a`$ which agree fully with $`t=6a`$. From the static potential along a lattice axis we determined the scale $`r_0`$ exactly as explained in \[?\]. The result is
$`r_0/a`$ $`=`$ $`5.29(6).`$ (35)
Comparing this number with the values of $`r_0/a`$ computed in quenched QCD \[?\], we see that our point eq. (33) in the SU(2) Higgs model corresponds in resolution to $`\beta 6`$ in the SU(3) Yang-Mills theory with Wilson action. With respect to our first work \[?\], the lattice resolution is almost a factor 2 better.
The static potentials $`V_\alpha (r)`$ as they are obtained from eq. (34) are not renormalised quantities because they contain self-energy contributions of the static quarks which diverge like $`1/a`$ in the continuum. We consider instead the differences $`V_\alpha (r)2\mu `$ which are free of divergences \[?\] and multiply them by $`r_0`$ to obtain dimensionless renormalised potentials. In Fig. 5, we represent the ground state and the first excited state static potentials. The ground state potential shows an approximate linear rise at intermediate distances: around separation
$`r_\mathrm{b}`$ $``$ $`1.9r_0`$ (36)
the potential flattens: the string breaks. As expected, for large distances the potential approaches the asymptotic value $`2\mu `$. The first excited potential comes very close to the ground state potential around $`r_\mathrm{b}`$ and rises linearly at larger distances. The scenario of string breaking as a level crossing phenomenon \[?\] is confirmed beautifully.
#### 3.2.2 Overlaps
A certain measure for the efficiency of a basis of fields eq. (30) used to extract the ground state potential is given by the overlap. Using the approximate ground state wave function $`v_0`$ obtained from the variational method we define the projected correlation function
$`\mathrm{\Omega }(t)={\displaystyle \underset{i,j}{}}v_{0,i}C_{ij}(t)v_{0,j}={\displaystyle \underset{\alpha }{}}(\omega _\alpha )^2\mathrm{e}^{tV_\alpha (r)},`$ (37)
with normalisation $`\mathrm{\Omega }(t_0=0)=1`$. The positive coefficients $`(\omega _\alpha )^2`$ may be interpreted as the square of the overlap of the true eigenstates of the Hamiltonian $`|\alpha `$ with the approximate ground state characterized by $`v_0`$. The “overlap” is an abbreviation commonly used to denote the ground state overlap, $`\omega _0`$. We determine $`\omega _0`$ straightforwardly from the correlation function $`\mathrm{\Omega }(t)`$ by noting that
$`\mathrm{ln}(\omega _0)^2`$ $`_\stackrel{}{t\mathrm{}}`$ $`{\displaystyle \frac{t+a}{a}}\mathrm{ln}\mathrm{\Omega }(t){\displaystyle \frac{t}{a}}\mathrm{ln}\mathrm{\Omega }(t+a).`$ (38)
We extract safe values for $`(\omega _0)^2`$ at $`t=7a`$, which agree fully with $`t=6a`$ and are shown in the left part of Fig. 6. Our basis of fields is big (and good) enough such that $`(\omega _0)^2`$ exceeds about 60% for all distances.
It is interesting to consider also the overlap for the (smeared) Wilson loops alone, i.e. we restrict the matrix correlation function to the $`3\times 3`$ sub-block associated with string-type states. Let us denote the corresponding projected correlation function by $`\mathrm{\Omega }_\mathrm{W}(t)`$ and the overlap by $`\omega _0^\mathrm{W}`$. The computation of $`\omega _0^\mathrm{W}`$ is more difficult because it turns out to be very small at large $`r`$. In the right part of Fig. 6, we present the results for two estimates of $`(\omega _0^\mathrm{W})^2`$. The triangles correspond to the estimates $`(\omega _0^{\mathrm{W},\mathrm{naive}})^2`$ obtained directly from eq. (38), with $`\mathrm{\Omega }(t)`$ replaced by $`\mathrm{\Omega }_\mathrm{W}(t)`$. The circles correspond to the more reliable estimate using the information from the full matrix correlation: the expression
$`(\omega _0^\mathrm{W})^2`$ $`_\stackrel{}{t\mathrm{}}`$ $`(\omega _0)^2{\displaystyle \frac{\mathrm{\Omega }_\mathrm{W}(t)}{\mathrm{\Omega }(t)}}`$ (39)
converges reasonably fast and $`(\omega _0^\mathrm{W})^2`$ can be estimated from the r.h.s. for large $`t`$ ($`t/a=79`$ in practice). Using eq. (39), we see that (smeared) Wilson loops alone have an overlap which drops at intermediate distances and they are clearly inadequate to extract the ground state at large $`r`$. On the contrary, $`(\omega _0^{\mathrm{W},\mathrm{naive}})^2`$ is above 50% at large distances: what is estimated here, is actually the coefficient $`(\omega _1^\mathrm{W})^2`$, i.e. the square of the overlap of the (smeared) Wilson loops with the first excited state (this statement is supported by direct calculation, see Sect. 3.2.3). The fact that $`\omega _1^\mathrm{W}`$ is so large might explain the problems encountered in QCD for observing string breaking from the analysis of a correlation matrix with Wilson loops only.
#### 3.2.3 Mixing
Finally, we want to show that string breaking is a mixing phenomenon involving string-type and meson-type states. This leads to the crossing of the energy levels seen above.
We consider the diagonal sub-blocks of the matrix correlation function eq. (31) corresponding to string-type states (fields $`i=1,2,3`$ in eq. (30)) and to meson-type states (fields $`i=4,5,\mathrm{},12`$ in eq. (30)) separately. We determine approximate ground state wave functions $`v_0^\mathrm{W}`$ in the subspace of the string-type states and $`v_0^\mathrm{M}`$ in the subspace of the meson-type states. With the help of these wave functions we construct a $`2\times 2`$ projected matrix correlation function
$`\mathrm{\Omega }_{kl}(t)={\displaystyle \underset{i,j}{}}v_{0,i}^kC_{ij}(t)v_{0,j}^l={\displaystyle \underset{\alpha }{}}\psi _l|\alpha \alpha |\psi _k\mathrm{e}^{tV_\alpha (r)}(k,l=\mathrm{W},\mathrm{M}),`$ (40)
where
$`|\psi _k={\displaystyle \underset{i}{}}v_{0,i}^k|i.`$ (41)
An inspection of
$`\mathrm{\Omega }_{\mathrm{WM}}(t_0=0)=\psi _\mathrm{M}|\psi _\mathrm{W}`$ (42)
shows that string-type and meson-type states are orthogonal only for large values of $`r`$ \[?\]. The coefficients
$`\omega _k(\alpha )`$ $``$ $`\alpha |\psi _k(k=\mathrm{W},\mathrm{M}),`$ (43)
in the expansion eq. (40), express the overlap of the string-type ($`k=\mathrm{W}`$) and meson-type ($`k=\mathrm{M}`$) states with the true eigenstates of the Hamiltonian. We can choose our phase conventions for the states such that the coefficients $`\omega _k(\alpha )`$ are real and in addition $`\omega _\mathrm{W}(0)>0`$ and $`\omega _\mathrm{W}(1)>0`$. We truncate the sum in eq. (40) after $`\alpha =1`$ and consider the diagonal matrix elements $`\mathrm{\Omega }_{kk}(t)`$ for two fixed times $`t=t_1`$ and $`t=t_2`$: inserting the known values for $`V_0(r)`$ and $`V_1(r)`$, we can solve for $`\omega _k^2(0)`$ and $`\omega _k^2(1)`$. The sign of the coefficients $`\omega _\mathrm{M}(0)`$ and $`\omega _\mathrm{M}(1)`$ is fixed by the off-diagonal matrix elements $`\mathrm{\Omega }_{\mathrm{WM}}(t_1)`$ and $`\mathrm{\Omega }_{\mathrm{WM}}(t_2)`$: we find that for all $`r`$, $`\omega _\mathrm{M}(0)>0`$ and $`\omega _\mathrm{M}(1)<0`$ (in our sign convention). The overlaps $`\omega _\mathrm{W}(0)`$ of the string-type states (circles) and $`\omega _\mathrm{M}(0)`$ of the meson-type states (triangles) with the ground state of the Hamiltonian are shown on the left of Fig. 7 and the corresponding overlaps $`\omega _\mathrm{W}(1)`$ and $`\omega _\mathrm{M}(1)`$ with the first excited eigenstate of the Hamiltonian on the right of Fig. 7. String-type states have a large overlap at short distances with the ground state and at large distances with the first excited state. Meson-type states have a large overlap at short distances with the first excited state and at large distances with the ground state.
In addition, we observe that the overlap of the meson-type states with the ground state is also large at very short distances. The explanation for this fact is that string-type and meson-type states have an overlap with each other at short distances. In the string breaking region around $`r/a=910`$, the overlaps of the string-type and meson-type states have similar magnitude, both when the ground state or the first excited state is considered. This fact is reflected in the crossing of the energy levels Fig. 5. Here, we would like to point out that the overlaps represented in Fig. 7 are not quantities which have a strict continuum limit. They depend on the $`\beta `$-value and the other parameters (e.g. of the smearing) that we consider. However, as long as one chooses a good basis (say with $`\omega _0>0.5`$) which can be separated into “string like” and “meson like”, the qualitative behavior in Fig. 7 is expected to persist also at smaller lattice spacings.
## 4 Scaling of the static potentials
In order to compare the renormalised static potentials at different values of the lattice spacing and estimate the size of scaling violations, a way of determining lines of constant physics (LCP) in the confinement “phase” of the SU(2) Higgs model is needed. This question has been addressed in \[?\], where a non-perturbative determination of the LCPs is described.
The bare parameters $`\kappa `$ and $`\lambda `$ are renormalised along a LCP by keeping two physical quantities $`F_1`$ and $`F_2`$ constant. A good choice is to take
$`F_1`$ $`=`$ $`r_0[2\mu V_0(r_0)]`$ (44)
and $`F_2`$ to be the generalised Binder cumulant $`c_3`$ defined in \[?\]. <sup>7</sup><sup>7</sup>7 Physically $`F_1`$ is a (non-perturbative) measure of the Higgs mass: the variation of $`F_1`$ with the bare parameters is dominantly caused by the variation of the meson mass $`\mu `$. This bound state mass is of course expected to depend strongly on the mass of its constituents. On the other hand, the interpretation of $`F_2`$ is less obvious. It was chosen to have a second renormalised quantity which is both sensitive to the bare coupling $`\lambda `$ and can be computed in the MC simulations\[?\].
Taking the parameter set (33) which we used in the simulation at $`\beta =2.4`$ we obtain the conditions
$`F_1=F_1^{}\mathrm{\hspace{0.17em}1.26}`$ and $`F_2=F_2^{}`$ (45)
defining a LCP (the numerical value $`F_2^{}`$ can be found in \[?\]). As a result of the non-perturbative matching, the parameter sets (33) and
$`\beta =2.2,`$ $`\kappa ^{}=\kappa (\lambda ^{}),`$ $`\lambda ^{}=0.96(10)`$ (46)
lie on a LCP. The value of $`\kappa `$ at $`\beta =2.2`$ is determined using the polynomial fit \[?\]
$`\kappa (\lambda )`$ $`=`$ $`0.3131+0.0564(\lambda 1)0.0286(\lambda 1)^2`$ (47)
$`+\mathrm{\hspace{0.17em}0.0198}(\lambda 1)^30.0246(\lambda 1)^4`$
which is obtained by requiring $`F_1=F_1^{}1.26`$ and correlates the uncertainties in $`\kappa `$ and $`\lambda `$ of the LCP. In Fig. 8, we compare the results for the renormalised ground state and first excited state static potentials that we obtain at $`\beta =2.4`$ and at $`\beta =2.2`$ along the LCP. The results for the potentials are compatible with scaling within minute errors under variation of the lattice spacing by almost a factor 2.
Another interesting issue is the $`\lambda `$-dependence of physical observables. Exploratory results \[?\] indicated that the physics of the SU(2) Higgs model in the confinement “phase” is weakly dependent on $`\lambda `$. This is certainly true near the continuum limit because it is well accepted that the scalar part of the SU(2) Higgs model is a trivial theory. Nevertheless, at finite value of the lattice spacing the model can be considered as an effective field theory with three independent renormalised couplings. In Fig. 9, we compare the static potentials for two different values $`\lambda =0.5`$ and $`\lambda =0.96`$ at $`\beta =2.2`$. The parameter $`\kappa `$ is determined from eq. (47). There is no significant difference between the different $`\lambda `$ values. This also means that the uncertainty in $`\lambda ^{}`$ in (46) is irrelevant for our scaling test Fig. 8.
## 5 Conclusions
We presented the results for the static potentials (ground state and first excited state) in the confinement “phase” of the SU(2) Higgs model. The string breaking or flattening of the ground state potential is clearly visible around separation $`r_b1.9r_0`$. The comparison with the first excited potential shows a nice crossing of the energy levels. The interpretation of string breaking as level crossing phenomenon between string-type and meson-type states is substantiated by the investigation of properly defined overlaps.
We also addressed the question of scaling violations in the measurements of the static potentials. They are shown to be tiny already at $`r_0/a>2.5`$. Moreover, the dependence of the static potentials on the Higgs quartic coupling $`\lambda `$ is very weak once the parameter $`\kappa `$ is determined by keeping the physical quantity $`F_1`$ constant. These results are a strong indication for a continuum-like behavior of the static potentials already at the relatively large values of the lattice spacing that we used.
The method for the determination of the static potential that we presented can be applied in QCD. Some steps in this direction have already been made\[?\]. The main problem for QCD is the statistical accuracy.
Acknowledgement. We thank the Konrad-Zuse-Zentrum für Informationstechnik Berlin (ZIB) for granting CPU-resources to this project.
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# 1 Introduction
## 1 Introduction
Anti-de Sitter (AdS) spacetimes naturally arise as the near-horizon geometries of non-dilatonic $`p`$-branes in supergravity theories. The metric for such a solution is usually the direct sum of AdS and an internal sphere. These geometries are of particular interest because of the conjecture that supergravity on such a background is dual to a conformal field theory on the boundary of the AdS . Examples include all the anti-de Sitter spacetimes AdS<sub>d</sub> with $`2d7`$, with the exception of $`d=6`$. The origin of AdS<sub>6</sub> is a little more involved, and it was first suggested in that it was related to the ten-dimensional massive type IIA theory. Recently, it was shown that the massive type IIA theory admits a warped-product solution of AdS<sub>6</sub> with $`S^4`$ , which turns out to be the near-horizon geometry of a semi-localised D4/D8 brane intersection . It is important that the warp factors depend only on the internal $`S^4`$ coordinates, since this implies that the reduced theory in $`D=6`$ has AdS spacetime as its vacuum solution. The consistent embedding of $`D=6`$, $`N=1`$ gauged supergravity in massive type IIA supergravity was obtained in . Ellipsoidal distributions of the D4/D8 system were also obtained, giving rise to AdS domain walls in $`D=6`$, supported by a scalar potential involving 3 scalars .
In fact, configurations with AdS in a warped spacetime are not rare occurrences. In , a semi-localised M5/M5 system was studied, and it was shown that the near-horizon geometry turns out to be a warped product of AdS<sub>5</sub> with an internal 6-space. This makes it possible to study AdS<sub>5</sub>/CFT<sub>4</sub> from the point of view of M-theory. In this paper, we shall consider AdS with a warped spacetime in a more general context and obtain such geometries for all the AdS<sub>d</sub>, as the near-horizon limits of semi-localised multiple intersections in both type IIA and type IIB theories.
The possibility of this construction is based on the following observations. As is well known, a non-dilatonic $`p`$-brane has the near-horizon geometry AdS$`{}_{d}{}^{}\times S^n`$. The internal $`n`$-sphere can be described geometrically as a foliation of $`S^p\times S^q`$ surfaces with $`n=p+q+1`$ (see appendix A), and so, in particular, if $`n4`$ the $`n`$-sphere can be viewed in terms of a foliation with $`S^3\times S^{n4}`$ surfaces, viz.
$$d\mathrm{\Omega }_n^2=d\alpha ^2+\mathrm{cos}^2\alpha d\mathrm{\Omega }_3^2+\mathrm{sin}^2\alpha d\mathrm{\Omega }_{n4}^2.$$
(1)
In appendix B, we show that when a non-dilatonic $`p`$-brane with an $`n`$-sphere in the transverse space intersects with a Kaluza-Klein monopole (a Taub-NUT with charge $`Q_\mathrm{N}`$) in a semi-localised manner, the net result turns out to be effectively a coordinate transformation of a solution with a distribution of pure $`p`$-branes with no NUT present. The round $`S^3`$ in (1) becomes the cyclic lens space $`S^3/Z_{Q_\mathrm{N}}`$ with metric
$$d\overline{\mathrm{\Omega }}_3^2=\frac{1}{4}d\mathrm{\Omega }_2^2+\frac{1}{4}(\frac{dy}{Q_\mathrm{N}}+\omega )^2,$$
(2)
where $`d\omega =\mathrm{\Omega }_2`$ is the volume form of the unit 2-sphere. This metric retains the same local structure as the standard round 3-sphere, and it has the same curvature tensor, but the $`y`$ coordinate on the $`U(1)`$ fibres is now identified with a period which is $`1/Q_\mathrm{N}`$ of the period for $`S^3`$ itself. We can now perform a dimensional reduction, or a T-duality transformation, on the fibre coordinate $`y`$, and thereby obtain AdS in a warped spacetime. The warp factor depends only on the internal “latitude” coordinate $`\alpha `$, but is independent of the lower-dimensional spacetime coordinates. In fact, the M5/M5 system with AdS<sub>5</sub> found in can be obtained in precisely such a manner from the D3-brane by using type IIA/IIB T-duality. Note that an isotropic $`p`$-brane can be viewed as carrying a single unit of NUT charge. Although this semi-localised way of introducing a Taub-NUT seems trivial, in that it amounts to a coordinate transformation, performing Kaluza-Klein reduction on the fibre coordinate does create a non-trivial intersecting component, since the Kaluza-Klein 2-form field strength now carries a non-trivial flux. This fact was used in to construct multi-charge $`p`$-branes starting from flat spacetime.
An analogous procedure can instead be applied to the anti-de Sitter spacetime, rather than the sphere, in the near-horizon limit AdS$`{}_{d}{}^{}\times S^n`$ of a non-dilatonic $`p`$-brane. As discussed in appendix A, AdS<sub>d</sub> can be described in terms of a foliation of AdS$`{}_{p}{}^{}\times S^q`$ surfaces with $`d=p+q+1`$ and so, in particular, for $`d4`$ it can be expressed as a foliation of AdS$`{}_{3}{}^{}\times S^{d3}`$:
$$ds_{\mathrm{AdS}_\mathrm{d}}^2=d\rho ^2+\mathrm{cosh}^2\rho ds_{\mathrm{AdS}_3}^2+\mathrm{sinh}^2\rho d\mathrm{\Omega }_{d4}^2.$$
(3)
In the presence of a pp-wave that is semi-localised on the world-volume of the $`p`$-brane, the AdS<sub>3</sub> turns out to have the form of a $`U(1)`$ bundle over AdS<sub>2</sub> ,
$$ds_{\mathrm{AdS}_3}^2=r^2W^1dt^2+\frac{dr^2}{r^2}+r^2W(dy+(W^11))dt)^2,$$
(4)
where $`W=1+Q_w/r^2`$, and $`Q_w`$ is the momentum carried by the pp-wave. This is precisely the structure of the extremal BTZ black hole . We can now perform a Kaluza-Klein reduction, or T-duality transformation, on the fibre coordinate $`y`$. In the near-horizon limit where the “1” in $`W`$ can be dropped, we obtain AdS<sub>2</sub> in a warped spacetime with a warp factor that depends only on the foliation coordinate, $`\rho `$.
A T-duality transformation on such a fibre coordinate of AdS<sub>3</sub> or $`S^3`$ has been called Hopf T-duality . It has the effect of (un)twisting the AdS<sub>3</sub> or $`S^3`$. The effect of this procedure on the six-dimensional dyonic string, whose near-horizon limit is AdS$`{}_{3}{}^{}\times S^3`$, was extensively studied in . In this paper, we apply the same technique to AdS<sub>3</sub> or $`S^3`$ geometries that are themselves factors in the foliation surfaces of certain larger-dimensional AdS spacetimes or spheres.
In section 2, we consider the semi-localised D3/NUT system and show that the effect of turning on the NUT charge $`Q_\mathrm{N}`$ in the intersection is merely to convert the internal 5-sphere, viewed as a foliation of $`S^1\times S^3`$, into a foliation of $`S^1\times (S^3/Z_{Q_\mathrm{N}})`$, where $`S^3/Z_{Q_\mathrm{N}}`$ is the cyclic lens space of order $`Q_\mathrm{N}`$. We can then perform a T-duality transformation on the Hopf fibre coordinate of the lens space and thereby obtain an AdS<sub>5</sub> in a warped spacetime as a solution in M theory, as the near-horizon geometry of a semi-localised M5/M5 system.
In section 3, we consider a semi-localised D3/pp-wave system, for which the AdS<sub>5</sub> becomes a foliation of a circle with the extremal BTZ black hole, which is locally AdS<sub>3</sub> and can be viewed as a $`U(1)`$ bundle over AdS<sub>2</sub>. We then perform a Hopf T-duality transformation on the fibre coordinate to obtain a solution with AdS<sub>2</sub> in a warped spacetime in M-theory, as the near-horizon geometry of a semi-localised M2/M2 system.
In sections 4 and 5, we apply the same analysis to the M2/NUT and M2/pp-wave systems, and the M5/NUT and M5/pp-wave systems, respectively; we obtain various configurations of AdS in warped spacetimes by performing Kaluza-Klein reductions and Hopf T-duality transformations on the fibre coordinates.
In section 6, we consider the D4/D8 system, which has the near-horizon geometry of a warped product of AdS<sub>6</sub> and $`S^4`$. We perform a Hopf T-duality transformation on the fibre coordinate of the foliating lens space of $`S^4`$, and thereby embed AdS<sub>6</sub> in a warped spacetime solution of type IIB theory.
We end with concluding remarks in section 7. In appendix A, we show how arbitrary-dimensional spheres and AdS spacetimes can be described in terms of foliations. In appendix B, we show that the solution describing the semi-local intersection of a non-dilatonic $`p`$-brane with a Kaluza-Klein monopole (Taub-NUT) is equivalent, after a coordinate transformation, to a solution purely composed of distributed $`p`$-branes, with no NUT.
## 2 D3/NUT systems and AdS<sub>5</sub> in M-theory from T-duality
AdS<sub>5</sub> spacetime arises naturally from type IIB theory as the near-horizon geometry of the D3-brane. Its origin in M-theory is more obscure. One way to embed the AdS<sub>5</sub> in M-theory is to note that $`S^5`$ can be viewed as a $`U(1)`$ bundle over $`CP^2`$, and hence we can perform a Hopf T-duality transformation on the $`U(1)`$ fibre coordinate. The resulting M-theory solution becomes AdS$`{}_{5}{}^{}\times CP^2\times T^2`$ . However, this solution is not supersymmetric at the level of supergravity, since $`CP^2`$ does not admit a spin structure. Charged spinors exist but, after making the T-duality transformation, the relevant electromagnetic field is described by the winding-mode vector and it is only in the full string theory that states charged with respect to this field arise. It was therefore argued in that the lack of supersymmetry (and indeed of any fermions at all) is a supergravity artifact and that, when the full string theory is considered, the geometry is supersymmetric. Such a phenomenon was referred as “supersymmetry without supersymmetry” in .
Recently, AdS<sub>5</sub> in warped eleven-dimensional spacetime was constructed in . It arises as the near-horizon limit of the semi-localised M5/M5 intersecting system. After performing a T-duality transformation, the warped spacetime of the near-horizon limit becomes AdS$`{}_{5}{}^{}\times (S^5/Z_{Q_\mathrm{N}})`$. In this section, we shall review this example in detail and show that the M5/M5 system originates from a semi-localised D3/NUT intersection in type IIB supergravity.
### 2.1 D3/NUT system
Any $`p`$-brane with a transverse space of sufficiently high dimension can intersect with a NUT. The D3/NUT solution of type IIB supergravity is given by
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{B}}^2`$ $`=`$ $`H^{1/2}(dt^2+dw_1^2+\mathrm{}+dw_3^2)+H^{1/2}(dx_1^2+dx_2^2`$ (5)
$`K(dz^2+z^2d\mathrm{\Omega }_2^2)+K^1(dy+Q_\mathrm{N}\omega )^2),`$
$`F_5`$ $`=`$ $`dtd^3wdH^1+(dtd^3wdH^1),`$
where $`z^2=z_1^2+z_2^2+z_3^2`$, and $`\omega `$ is a 1-form satisfying $`d\omega =\mathrm{\Omega }_2`$. The solution can be best illustrated by the following diagram:
$`t`$ $`w_1`$ $`w_2`$ $`w_3`$ $`x_1`$ $`x_2`$ $`z_1`$ $`z_2`$ $`z_3`$ $`y`$ D3 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ NUT $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $`K`$
Diagram 1. The D3/NUT brane intersection. Here $`\times `$ and $``$ denote the
worldvolume and transverse space coordinates respectively,
and $``$ denotes the fibre coordinate of the Taub-NUT.
The function $`K`$ is associated with the NUT component of the intersection; it is a harmonic function in the overall transverse Euclidean 3-space coordinatised by $`z_i`$. The function $`H`$ is associated with the D3-brane component. It satisfies the equation
$$_\stackrel{}{z}^2H+K_\stackrel{}{x}^2H=0.$$
(6)
Equations of this type were also studied in . In the absence of NUT charge, i.e. $`K=1`$, the function $`H`$ is harmonic in the the transverse 6-space of the D3-brane. When the NUT charge $`Q_\mathrm{N}`$ is non-zero, $`K`$ is instead given by
$$K=1+\frac{Q_\mathrm{N}}{z},$$
(7)
and the function $`H`$ cannot be solved analytically, but only in terms of a Fourier expansion in $`\stackrel{}{x}`$ coordinates. The usual way to solve for the solution is to consider the zero-modes in the Fourier expansion. In other words, one assumes that $`H`$ is independent of $`\stackrel{}{x}`$. The consequence of this assumption is that the resulting metric no longer has an AdS structure in its near-horizon region. In , it was observed that an explicit closed-form solution for $`H`$ can be obtained in the case where the “1” in function K is dropped. This solution is given by
$$K=\frac{Q_\mathrm{N}}{z},H=1+\underset{k}{}\frac{Q_k}{(|\stackrel{}{x}\stackrel{}{x}_{0k}|^2+4Q_\mathrm{N}z)^2}.$$
(8)
In this paper, we shall consider the case where the D3-brane is located at the origin of the $`\stackrel{}{x}`$ space and so we have
$$H=1+\frac{Q}{(x^2+4Q_\mathrm{N}z)^2},$$
(9)
where $`x^2=x^ix^i`$. Thus, the D3-brane is also localised in the space of the $`\stackrel{}{x}`$ as well. Let us now make a coordinate transformation
$$x_1=r\mathrm{cos}\alpha \mathrm{cos}\theta ,x_2=r\mathrm{cos}\alpha \mathrm{sin}\theta ,z=\frac{1}{4}Q_{\mathrm{N}}^{}{}_{}{}^{1}r^2\mathrm{sin}^2\alpha .$$
(10)
In terms of the new coordinates, the metric for the solution becomes
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{B}}^2`$ $`=`$ $`H^{1/2}(dt^2+dw_1^2+dw_2^2+dw_3^2)+H^{1/2}(dr^2+r^2dM_5^2),`$
$`H`$ $`=`$ $`1+{\displaystyle \frac{Q}{r^4}}.`$ (11)
where
$$dM_5^2=d\alpha ^2+c^2d\theta ^2+\frac{1}{4}s^2\left(d\mathrm{\Omega }_2^2+(\frac{dy}{Q_\mathrm{N}}+\omega )^2\right),$$
(12)
and $`s=\mathrm{sin}\alpha `$, $`c=\mathrm{cos}\alpha `$. Thus, we see that $`dM_5^2`$ describes a foliation of $`S^1`$ times the lens space $`S^3/Z_{Q_\mathrm{N}}`$. For a unit NUT charge, $`Q_\mathrm{N}=1`$, the metric $`dM_5^2`$ describes the round 5-sphere and the solution becomes an isotropic D3-brane. It is interesting to note that the regular D3-brane can be viewed as a semi-localised D3-brane intersecting with a NUT with unit charge.<sup>1</sup><sup>1</sup>1An analogous observation was also made in , where multi-charge solutions were obtained from flat space by making use of the fact that $`S^3`$ can be viewed as a $`U(1)`$ bundle over $`S^2`$. In other words, flat space can be viewed as a NUT, with unit charge, located on the $`U(1)`$ coordinate. In the near-horizon limit $`r0`$, where the constant 1 in the function $`H`$ can be dropped, the metric becomes AdS$`{}_{5}{}^{}\times M_5`$:
$$ds_{10\mathrm{I}\mathrm{I}\mathrm{B}}^2=Q^{1/2}r^2(dt^2+dw^idw^i)+Q^{1/2}\frac{dr^2}{r^2}+Q^{1/2}\left(d\alpha ^2+c^2d\theta ^2+\frac{1}{4}s^2(d\mathrm{\Omega }_2^2+(\frac{dy}{Q_\mathrm{N}}+\omega )^2)\right).$$
(13)
### 2.2 M5/M5 system and AdS<sub>5</sub> in M-theory
Since the near-horizon limit of a semi-localised D3-brane/NUT is a direct product of AdS<sub>5</sub> and an internal 5-sphere that is a foliation of a circle times a lens space, it follows that if we perform a T-duality transformation on the $`U(1)`$ fibre coordinate $`y`$, we shall obtain AdS<sub>5</sub> in a warped spacetime as a solution of the type IIA theory. The warp factor is associated with the scale factor $`s^2`$ of $`dy^2`$ in (13). This type of Hopf T-duality has the effect of untwisting a 3-sphere into $`S^2\times S^1`$ . If one performs the T-duality transformation on the original full solution (5), rather than concentrating on its near-horizon limit, then one obtains a semi-localised NS5/D4 system of the type IIA theory, which can be further lifted back to $`D=11`$ to become a semi-localised M5/M5 system, obtained in . In , the near-horizon structures of these semi-localised branes of M-theory were analysed, and AdS<sub>5</sub> was obtained as a warped spacetime solution. We refer the readers to Ref. and shall not discuss this solution further, but only mention that, from the above analysis, it can be obtained by implementing the T-duality transformation on the coordinate $`y`$ in (13).
## 3 D3/pp-wave system and extremal BTZ black hole
In this section, we study the semi-localised pp-wave intersecting with a D3-brane. The solution is given by
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{B}}^2`$ $`=`$ $`H^{1/2}\left(W^1dt^2+W(dy+(W^11)dt)^2+dx_1^2+dx_2^2\right)`$ (14)
$`+H^{1/2}(dz_1^2+\mathrm{}dz_6^2),`$
$`F_{\left(5\right)}`$ $`=`$ $`dtdydx_1dx_2dH^1+(dtdydx_1dx_2dH^1),`$
The solution can be illustrated by the following diagram
$`t`$ $`y`$ $`x_1`$ $`x_2`$ $`z_1`$ $`z_2`$ $`z_3`$ $`z_4`$ $`z_5`$ $`z_6`$ D3 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ wave $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`W`$
Diagram 2. The D3/pp-wave brane intersection. Here $``$ denotes the wave coordinate.
In the usual construction of such an intersection, the harmonic functions $`H`$ and $`W`$ depend only on the overall transverse space coordinates $`\stackrel{}{z}`$. The near-horizon limit of the solution then becomes K$`{}_{5}{}^{}\times S^6`$, where K<sub>5</sub> is the generalised Kaigorodov metric in $`D=5`$, and the geometry is dual to a conformal field theory in the infinite momentum frame . On the other hand, the semi-localised solution is given by
$$H=\frac{Q}{|\stackrel{}{z}|^4},W=1+Q_w(|\stackrel{}{x}|^2+\frac{Q}{|\stackrel{}{z}|^2}).$$
(15)
We now let
$$x_1=\frac{1}{r}\mathrm{cos}\alpha \mathrm{cos}\theta ,x_2=\frac{1}{r}\mathrm{cos}\alpha \mathrm{sin}\theta ,z_i=\frac{rQ^{1/2}}{\mathrm{sin}\alpha }\nu _i,$$
(16)
where $`\nu _i`$ coordinates, satisfying $`\nu _i\nu _i=1`$, define a 5-sphere with the unit sphere metric $`d\mathrm{\Omega }_5^2=d\nu _id\nu _i`$. Using these coordinates, the metric of the semi-localised D3/wave system becomes
$$ds_{10\mathrm{I}\mathrm{I}\mathrm{B}}^2=Q^{1/2}s^2(ds_{\mathrm{AdS}_3}^2+d\alpha ^2+c^2d\theta ^2+s^2d\mathrm{\Omega }_5^2),$$
(17)
where $`ds_{\mathrm{AdS}_3}^2`$ is given by
$`ds_{\mathrm{AdS}_3}^2`$ $`=`$ $`r^2W^1dt^2+r^2W(dy+(W^11)dt)^2+{\displaystyle \frac{dr^2}{r^2}},`$
$`W`$ $`=`$ $`1+{\displaystyle \frac{Q_w}{r^2}}.`$ (18)
Note that the above metric is exactly the extremal BTZ black hole , and hence it is locally AdS<sub>3</sub>. Thus we have demonstrated that the semi-localised D3/pp-wave system is in fact a warped product of AdS<sub>3</sub> (the extremal BTZ black hole) with a 7-sphere, where $`S^7`$ is described as a foliation of $`S^1\times S^5`$ surfaces.<sup>2</sup><sup>2</sup>2A D3-brane with an $`S^3\times \text{I}\mathrm{R}`$ worlvolume was obtained in . In that solution, which was rather different from ours, the dilaton was not constant. Note that the metric (17) can also be expressed as a direct product of AdS$`{}_{5}{}^{}\times S^5`$, with the AdS<sub>5</sub> metric written in the following form:
$$ds_5^2=s^2(ds_{\mathrm{AdS}_3}^2+d\alpha ^2+c^2d\theta ^2).$$
(19)
Making a coordinate transformation $`\mathrm{tan}(\alpha /2)=e^\rho `$, the metric becomes
$$ds_5^2=d\rho ^2+\mathrm{sinh}^2\rho d\theta ^2+\mathrm{cosh}^2\rho ds_{\mathrm{AdS}_3}^2,$$
(20)
which is precisely the AdS<sub>5</sub> metric written as a foliation of a circle times AdS<sub>3</sub> (see appendix A).
The extremal BTZ black hole occurs as the near-horizon geometry of the boosted dyonic string in six-dimensions, which can be viewed as an intersection of a string and a 5-brane in $`D=10`$. The boosted D1/D5 system was used to obtain the first stringy interpretation of the microscopic entropy of the Reissner-Nordström black hole in $`D=5`$. The boosted dyonic string has three parameters, namely the electric and magnetic charges $`Q_e`$, $`Q_m`$, and the boost momentum parameter $`Q_w`$. On the other hand, the extremal BTZ black hole itself has only two parameters: the cosmological constant, proportional to $`\sqrt{Q_eQ_m}`$, and the mass (which is equal to the angular momentum in the extremal limit), which is related to $`Q_w`$. (Analogous discussion applies to $`D=4`$ .) In our construction of the BTZ black hole in warped spacetime, the original configuration also has only two parameters, namely the D3-brane charge $`Q`$, related to the cosmological constant of the BTZ black hole, and the pp-wave charge, associated with the mass.
### 3.1 NS1/D2 and M2/M2 systems and AdS<sub>2</sub>
We can perform a T-duality transformation on the coordinate $`y`$ in the previous solution. The D3-brane is T-dual to the D2-brane, and the wave is T-dual to the NS-NS string. Thus the D3/pp-wave system of the type IIB theory becomes an NS1/D2 system in the type IIA theory, given by
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2`$ $`=`$ $`W^{1/4}H^{3/8}[(WH)^1dt^2+H^1(dx_1^2+dx_2^2)+W^1dy_1^2,`$
$`+dz_1^2+\mathrm{}dz_6^2],`$
$`e^\varphi `$ $`=`$ $`W^{1/2}H^{1/4},`$ (21)
$`F_{\left(4\right)}`$ $`=`$ $`dtdx_1dx_2dH^1,F_{\left(3\right)}=dtdy_1dW^1.`$
This solution can be represented diagrammatically as follows:
$`t`$ $`x_1`$ $`x_2`$ $`y_1`$ $`z_1`$ $`z_2`$ $`z_3`$ $`z_4`$ $`z_5`$ $`z_6`$ D2 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ NS1 $`\times `$ $``$ $``$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $`W`$
Diagram 3. The NS1/D2 brane intersection.
In the near-horizon limit where the 1 in $`W`$ is dropped, the metric of the NS1/D2 system (21), in terms of the new coordinates (16), becomes
$$ds_{10}^2=Q_w^{1/4}Q^{5/8}s^{5/2}\left(ds_{\mathrm{AdS}_2}^2+d\alpha ^2+c^2d\theta ^2+s^2d\mathrm{\Omega }_5^2+(Q_wQ)^1s^4dy_1^2\right),$$
(22)
where
$$ds_{\mathrm{AdS}_2}^2=\frac{r^4dt^2}{Q_w}+\frac{dr^2}{r^2}.$$
(23)
Thus we see that the near-horizon limit of the NS1/D2 system is a warped product of AdS<sub>2</sub> with a certain internal 8-space, which is a warped product of a 7-sphere with a circle.
We can further lift the solution back to $`D=11`$, where it becomes a semi-localised M2/M2 system,
$`ds_{11}^2`$ $`=`$ $`(WH)^{1/3}[(WH)^1dt^2+H^1(dx_1^2+dx_2^2)+W^1(dy_1^2+dy_2^2),`$
$`+dz_1^2+\mathrm{}+dz_6^2],`$
$`F_{\left(4\right)}`$ $`=`$ $`dtdx_1dx_2dH^1+dtdy_1dy_2dW^1.`$ (24)
The configuration for this solution can be summarised in the following diagram:
$`t`$ $`x_1`$ $`x_2`$ $`y_1`$ $`y_2`$ $`z_1`$ $`z_2`$ $`z_3`$ $`z_4`$ $`z_5`$ $`z_6`$ M2 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ M2 $`\times `$ $``$ $``$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $`W`$
Diagram 4. The M2-M2 brane intersection.
It is straightforward to verify that the near-horizon geometry of this system is a warped product of AdS<sub>2</sub> with a certain 9-space, namely
$$ds_{11}^2=Q_w^{1/3}Q^{2/3}s^{8/3}(ds_{\mathrm{AdS}_2}^2+d\alpha ^2+c^2d\theta ^2+s^2d\mathrm{\Omega }_5^2+(Q_wQ)^1s^4(dy_1^2+dy_2^2)),$$
(25)
where $`ds_{\mathrm{AdS}_2}^2`$ is an AdS<sub>2</sub> metric given by (23), and the internal 9-space is a warped product of a 7-sphere and a 2-torus.
### 3.2 Further possibilities
Note that in the above examples, we can replace the round sphere $`d\mathrm{\Omega }_5^2`$ by a lens space of the following form:
$$d\mathrm{\Omega }_5^2=d\stackrel{~}{\alpha }^2+\stackrel{~}{c}^2d\stackrel{~}{\theta }^2+\stackrel{~}{s}^2(d\stackrel{~}{\mathrm{\Omega }}_2^2+(\frac{d\stackrel{~}{y}}{\stackrel{~}{Q}_\mathrm{N}}+\stackrel{~}{\omega })^2),$$
(26)
where $`\stackrel{~}{c}\mathrm{cos}\stackrel{~}{\alpha }`$, $`\stackrel{~}{s}\mathrm{sin}\stackrel{~}{\alpha }`$ and $`d\stackrel{~}{\omega }=\stackrel{~}{\mathrm{\Omega }}_2`$. As we have discussed in appendix B, this can be viewed as an additional NUT with charge $`\stackrel{~}{Q}_\mathrm{N}`$ intersecting with the system. We can now perform a Kaluza-Klein reduction or T-duality transformation on the fibre coordinate $`\stackrel{~}{y}`$, leading to many further examples of warped products of AdS<sub>2</sub> or AdS<sub>3</sub> with certain internal spaces. The warp factors again depend only on the coordinates of the internal space. These geometries can be viewed as the near-horizon limits of three intersecting branes, with charges $`Q`$, $`Q_\mathrm{N}`$ and $`\stackrel{~}{Q}_\mathrm{N}`$. Of course, this system can equally well be obtained by replacing the horospherical AdS<sub>5</sub> in (13) with (19).
For example, let us consider the M2/M2 system with an additional NUT component. The solution of this semi-localised intersecting system is given by
$`ds_{11}^2`$ $`=`$ $`(WH)^{1/3}[(WH)^1dt^2+H^1(dx_1^2+dx_2^2)+W^1(dy_1^2+dy_2^2),`$
$`+K(dz^2+z^2d\mathrm{\Omega }_2^2)+K^1(dy+Q_\mathrm{N}\omega )^2+du_1^2+du_2^2],`$
$`F_{\left(4\right)}`$ $`=`$ $`dtdx_1dx_2dH^1+dtdy_1dy_2dW^1.`$ (27)
where the functions $`H`$, $`W`$ and $`K`$ are given by
$$H=\frac{Q}{(|\stackrel{}{u}|^2+4Q_\mathrm{N}z)^2},W=1+Q_w(|\stackrel{}{x}|^2+\frac{Q}{|\stackrel{}{u}|^2+4Q_\mathrm{N}z}),K=\frac{Q_\mathrm{N}}{z}.$$
(28)
We illustrate this solution in the following diagram:
$`t`$ $`x_1`$ $`x_2`$ $`y_1`$ $`y_2`$ $`z_1`$ $`z_2`$ $`z_3`$ $`y`$ $`u_1`$ $`u_2`$ M2 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ M2 $`\times `$ $``$ $``$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $`W`$ NUT $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $`\times `$ $`\times `$ $`K`$
Diagram 5. The M2/M2/NUT brane intersection.
The near-horizon structure of this solution is basically the same as that of the M2/M2 system with the round $`S^3`$ in the foliation replaced by the lens space $`S^3/Z_{Q_\mathrm{N}}`$. We can now perform Kaluza-Klein reduction on the fibre coordinate $`y`$ and the solution becomes the semi-localised D2/D2/D6 brane intersection, given by
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2`$ $`=`$ $`(WH)^{3/8}K^{1/8}[(WH)^1dt^2+H^1(dx_1^2+dx_2^2)+W^1(dy_1^2+dy_2^2),`$
$`+K(dz^2+z^2d\mathrm{\Omega }_2^2)+du_1^2+du_2^2],`$
$`F_{\left(4\right)}`$ $`=`$ $`dtdx_1dx_2dH^1+dtdy_1dy_2dW^1.`$ (29)
$`e^\varphi `$ $`=`$ $`(WH)^{1/4}K^{3/4},F_{\left(2\right)}=Q_\mathrm{N}\mathrm{\Omega }_2.`$ (30)
The solution can be illustrated by the following diagram:
$`t`$ $`x_1`$ $`x_2`$ $`y_1`$ $`y_2`$ $`z_1`$ $`z_2`$ $`z_3`$ $`u_1`$ $`u_2`$ D2 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ D2 $`\times `$ $``$ $``$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $`W`$ D6 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $`\times `$ $`\times `$ $`K`$
Diagram 6. The D2/D2/D6 brane intersection.
## 4 M2/NUT and M2/pp-wave systems
In this section, we apply an analogous analysis to the M2-brane. We show that the semi-localised M2-brane intersecting with a NUT is in fact an isotropic M2-brane with the internal 7-sphere itself being described as a foliation of a regular $`S^3`$ and lens space $`S^3/Z_{Q_\mathrm{N}}`$, where $`Q_\mathrm{N}`$ is the NUT charge. Reducing the system to $`D=10`$, we obtain a semi-localised D2/D6 system whose near-horizon geometry is a warped product of AdS<sub>4</sub> with an internal 6-space. We also show that a semi-localised pp-wave intersecting with the M2-brane is in fact a warped product of AdS<sub>3</sub> (the BTZ black hole) and an 8-space. The system can be reduced to $`D=10`$ to become a semi-localised D0/NS1 intersection.
### 4.1 M2-brane/NUT system
The solution for the intersection of an M2-brane and a NUT is given by
$`ds_{11}^2`$ $`=`$ $`H^{2/3}(dt^2+dw_1^2+dw_2^2)+H^{1/3}(dx_1^2+\mathrm{}+dx_4^2`$
$`+K(dz^2+z^2d\mathrm{\Omega }_2^2)+K^1(dy+Q_\mathrm{N}\omega )^2),`$
$`F_{\left(4\right)}`$ $`=`$ $`dtdw_1dw_2dH^1,`$ (31)
where $`z^2=z_1^2+z_2^2+z_3^2`$ and $`d\omega =\mathrm{\Omega }_2`$. The solution can be illustrated by the following diagram:
$`t`$ $`w_1`$ $`w_2`$ $`x_1`$ $`x_2`$ $`x_3`$ $`x_4`$ $`z_1`$ $`z_2`$ $`z_3`$ $`y`$ M2 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ NUT $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $`K`$
Diagram 7. The M2/NUT brane intersection.
If the function $`K`$ associated with the NUT components of the intersection takes the form $`K=Q_\mathrm{N}/z`$, then the function $`H`$ associated with the M2-brane component can be solved in the semi-localised form
$$H=1+\frac{Q}{(|\stackrel{}{x}|^2+4Q_\mathrm{N}z)^3}.$$
(32)
Thus, the solution is also localised on the space of the $`\stackrel{}{x}`$ coordinates. Let us now make a coordinate transformation
$$x_i=r\mathrm{cos}\alpha \mu _i,,z=\frac{1}{4}Q_{\mathrm{N}}^{}{}_{}{}^{1}r^2\mathrm{sin}^2\alpha ,$$
(33)
where $`\mu _i\mu _i=1`$, defining a 3-sphere, with the unit 3-sphere metric given by $`d\mathrm{\Omega }_3^2=d\mu _id\mu _i`$. In terms of the new coordinates, the metric for the solution becomes
$`ds_{11}^2`$ $`=`$ $`H^{2/3}(dt^2+dw_1^2+dw_2^2)+H^{1/3}(dr^2+r^2dM_7^2),`$
$`H`$ $`=`$ $`1+{\displaystyle \frac{Q}{r^6}},`$ (34)
where
$$dM_7^2=d\alpha ^2+c^2d\mathrm{\Omega }_3^2+\frac{1}{4}s^2\left(d\mathrm{\Omega }_2^2+(\frac{dy}{Q_\mathrm{N}}+\omega )^2\right).$$
(35)
Thus we see that $`dM_7^2`$ is a foliation of a regular 3-sphere, together with a lens space $`S^3/Z_{Q_\mathrm{N}}`$. When $`Q_\mathrm{N}=1`$ the metric $`dM_7^2`$ describes a round 7-sphere and the solution becomes an isotropic M2-brane. Interestingly, the regular M2-brane can be viewed as an intersecting semi-localised M2-brane with a NUT of unit charge. In the near-horizon limit $`r0`$, where the 1 in the function $`H`$ can be dropped, the metric becomes AdS$`{}_{4}{}^{}\times M_7`$.
### 4.2 D2-D6 system
In the M2-brane and NUT intersection (31), we can perform a Kaluza-Klein reduction on the $`y`$ coordinate. This gives rise to a semi-localised intersection of D2-branes and D6-branes:
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2`$ $`=`$ $`H^{5/8}K^{1/8}(dt^2+dw_1^2+dw_2^2)+H^{3/8}K^{1/8}(dx_1^2+\mathrm{}+dx_4^2)`$
$`H^{3/8}K^{7/8}(dz_1^2+dz_2^2+dz_3^2),`$
$`e^\varphi `$ $`=`$ $`H^{1/4}K^{3/4},`$ (36)
$`F_{\left(4\right)}`$ $`=`$ $`dtd^2wdH^1,F_2=e^{3/2\varphi }(dtd^2wd^4xdK^1).`$
The solution can be illustrated by the following diagram
$`t`$ $`w_1`$ $`w_2`$ $`x_1`$ $`x_2`$ $`x_3`$ $`x_4`$ $`z_1`$ $`z_2`$ $`z_3`$ D2 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ D6 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $`K`$
Diagram 8. The D2/D6 brane intersection.
Again, in the usual construction of a D2-D6 system, the harmonic functions $`H`$ and $`K`$ are taken to depend only on the overall transverse space coordinates $`\stackrel{}{z}`$. In the semi-localized construction, the function $`H`$ depends on $`\stackrel{}{x}`$ as well. In terms of the new coordinates defined in (33), the metric becomes
$$ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2=(\frac{rs}{2Q_\mathrm{N}})^{1/4}[H^{5/8}(dt^2+dw_1^2+dw_2^2)+H^{3/8}(dr^2+r^2(d\alpha ^2+c^2d\mathrm{\Omega }_3^2+\frac{1}{4}s^2d\mathrm{\Omega }_2^2)].$$
(37)
Thus, in the near-horizon limit where the 1 in $`H`$ can be dropped, the solution becomes a warped product of AdS<sub>4</sub> with an internal 6-space:
$$ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2=(2Q_\mathrm{N})^{1/4}Q^{3/8}s^{1/4}(ds_{\mathrm{AdS}_4}^2+d\alpha ^2+c^2d\mathrm{\Omega }_3^2+\frac{1}{4}s^2d\mathrm{\Omega }_2^2),$$
(38)
where $`ds_4^2`$ is the metric on AdS<sub>4</sub>, given by
$$ds_{\mathrm{AdS}_4}^2=\frac{r^4}{Q}(dt^2+dw_1^2+dw_2^2)+\frac{dr^2}{r^2}.$$
(39)
The internal 6-space is a warped product of a 4-sphere with a 2-sphere.
### 4.3 AdS<sub>4</sub> in type IIB from T-duality
In the above discussion, we found that our starting point is effectively to replace the round 7-sphere of the M2-brane by the foliation of a round 3-sphere together with a lens space $`S^3/Z_{Q_\mathrm{N}}`$. We can also replace the round 3-sphere by another lens space $`S^3/Z_{\stackrel{~}{Q}_\mathrm{N}}`$, given by
$$d\overline{\mathrm{\Omega }}_3^2=\frac{1}{4}\left(d\stackrel{~}{\mathrm{\Omega }}_2^2+(\frac{d\stackrel{~}{y}}{\stackrel{~}{Q}_\mathrm{N}}+\omega )^2\right).$$
(40)
As discussed in the appendix, the lens space arises from introducing a NUT around the fibre coordinate $`\stackrel{~}{y}`$, with NUT charge $`\stackrel{~}{Q}_\mathrm{N}`$. The system can then be viewed as the near-horizon limit of three intersecting branes, with charges $`Q`$, $`Q_\mathrm{N}`$ and $`\stackrel{~}{Q}_\mathrm{N}`$. For example, with this replacement the D2/D6 system becomes a D2/D6/NUT system. Performing a T-duality transformation on the fibre coordinate $`\stackrel{~}{y}`$, the $`S^3`$ untwists to become $`S^2\times S^1`$. The resulting type IIB metric is given by
$$ds_{10\mathrm{I}\mathrm{I}\mathrm{B}}^2=\left(\frac{Qsc}{4Q_\mathrm{N}\stackrel{~}{Q}_\mathrm{N}}\right)^{1/2}\left(ds_{\mathrm{AdS}_4}^2+d\alpha ^2+\frac{1}{4}c^2d\stackrel{~}{\mathrm{\Omega }}_2^2+\frac{1}{4}s^2d\mathrm{\Omega }_2^2+\frac{(4Q_\mathrm{N}\stackrel{~}{Q}_\mathrm{N})^2}{Qs^2c^2}d\stackrel{~}{y}^2\right).$$
(41)
This metric can be viewed as describing the near-horizon geometry of a semi-localised D3/D5/NS5 system in the type IIB theory. This metric (41) provides a background for consistent reduction of type IIB supergravity to give rise to four-dimensional gauged supergravity with AdS background.
In order to construct the semi-localised D3/D5/NS5 intersecting system in the type IIB theory, we start with the D2/D6/NUT system, given by
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2`$ $`=`$ $`H^{5/8}K^{1/8}(dt^2+dw_1^2+dw_2^2)+H^{3/8}K^{7/8}(dz_1^2+dz_2^2+dz_3^2)`$
$`+H^{3/8}K^{1/8}(\stackrel{~}{K}(dx^2+x^2d\stackrel{~}{\mathrm{\Omega }}_2^2)+\stackrel{~}{K}^1(dy+\stackrel{~}{Q}_\mathrm{N}\stackrel{~}{\omega })^2),`$
$`e^\varphi `$ $`=`$ $`H^{1/4}K^{3/4},`$ (42)
$`F_{\left(4\right)}`$ $`=`$ $`dtd^2wdH^1,F_2=e^{3/2\varphi }(dtd^2wd^4xdK^1).`$
where $`x^2=x_1^2+x_2^2+x_3^2`$ and the functions $`H`$, $`K`$ and $`\stackrel{~}{K}`$ are given by
$$H=1+\frac{Q}{(4\stackrel{~}{Q}_\mathrm{N}x+4Q_\mathrm{N}z)^3},K=\frac{Q_\mathrm{N}}{z},\stackrel{~}{K}=\frac{\stackrel{~}{Q}_\mathrm{N}}{x}.$$
(43)
It is instructive to illustrate the solution in the following diagram:
$`t`$ $`w_1`$ $`w_2`$ $`x_1`$ $`x_2`$ $`x_3`$ $`y`$ $`z_1`$ $`z_2`$ $`z_3`$ D2 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ D6 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $`K`$ NUT $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $`\times `$ $`\times `$ $`\times `$ $`\stackrel{~}{K}`$
Diagram 9. The D2/D6/NUT system
We can now perform the T-duality on the coordinate $`y`$, and obtain the semi-localised D3/D5/NS5 intersection of the type IIB theory, given by
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{B}}^2`$ $`=`$ $`H^{1/2}(K\stackrel{~}{K})^{1/4}[dt^2+dw_1^2+dw_2^2`$ (44)
$`H\stackrel{~}{K}(dx_1^2+dx_2^2+dx_3^2)+K\stackrel{~}{K}dy^2+HK(dz_1^2+dz_2^2+dz_3^2)].`$
It is straightforward to verify that the near-horizon structure of the above D3/D5/NS5 system is of the form (41). The solution can be illustrated by the following diagram:
$`t`$ $`w_1`$ $`w_2`$ $`x_1`$ $`x_2`$ $`x_3`$ $`y`$ $`z_1`$ $`z_2`$ $`z_3`$ D3 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $`\times `$ $``$ $``$ $``$ $`H`$ D5 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $`K`$ NS5 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $`\times `$ $`\times `$ $`\times `$ $`\stackrel{~}{K}`$
Diagram 10. The D3/D5/NS5 system
### 4.4 M2/pp-wave system
The M2/pp-wave solution is given by
$`ds_{11}^2`$ $`=`$ $`H^{2/3}(W^1dt+W(dy+(W^11)dt)^2+dx^2)+H^{1/3}(dz^2+z^2d\mathrm{\Omega }_7^2),`$
$`F_{\left(4\right)}`$ $`=`$ $`dtdydxdH^1.`$ (45)
The solution can be illustrated by the following diagram:
$`t`$ $`y_1`$ $`x_1`$ $`z_1`$ $`z_2`$ $`z_3`$ $`z_4`$ $`z_5`$ $`z_6`$ $`z_7`$ $`z_8`$ M2 $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ wave $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`W`$
Diagram 11. The M2/pp-wave brane intersection.
When both functions $`H`$ and $`W`$ are harmonic on the overall transverse space of the $`z^i`$ coordinates, the metric becomes a direct product of the Kaigorodov metric with a 7-sphere in the near-horizon limit. Here, we instead consider a semi-localised solution, with $`H`$ and $`K`$ given by
$$H=\frac{Q}{z^6},W=1+Q_w(x^2+\frac{Q/4}{z^4}).$$
(46)
Making the coordinate transformation
$$x=\frac{\mathrm{cos}\alpha }{r},z^2=\frac{rQ^{1/2}}{2\mathrm{sin}\alpha },$$
(47)
the metric becomes AdS$`{}_{4}{}^{}\times S^7`$, with
$$ds_{11}^2=\frac{Q^{1/3}}{4s^2}(ds_{\mathrm{AdS}_3}^2+d\alpha ^2)+Q^{1/3}d\mathrm{\Omega }_7^2.$$
(48)
Here $`ds_{\mathrm{AdS}_3}^2`$ is the metric of AdS<sub>3</sub> (the BTZ black hole), given by (18). Thus, we have demonstrated that the semi-localised M2/pp-wave system is a warped product of AdS<sub>3</sub> and an 8-space. Making the coordinate transformation $`\mathrm{tan}(\alpha /2)=e^\rho `$, the first part of (48) can be expressed as
$$ds_4^2=d\rho ^2+\mathrm{cosh}^2\rho ds_{\mathrm{AdS}_3}^2.$$
(49)
This is AdS<sub>4</sub> expressed as a foliation of AdS<sub>3</sub> (see appendix A).
### 4.5 The NS1/D0 system
Reducing the above solution on the coordinate $`y_1`$, it becomes an intersecting NS1/D0 system, with
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}`$ $`=`$ $`H^{3/4}W^{7/8}\left(dt^2+Wdx^2+WH(dz_1^2+\mathrm{}+dz_8^2)\right),`$
$`F_{\left(3\right)}`$ $`=`$ $`dtdxdH^1,F_{\left(2\right)}=dtdW^1,`$
$`e^\varphi `$ $`=`$ $`H^{1/2}W^{3/4}.`$ (50)
The metric of the near-horizon region describes a warped product of AdS<sub>2</sub> with an 8-space:
$$ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2=8^{3/4}Q^{3/8}Q_w^{1/8}s^{9/4}(ds_{\mathrm{AdS}_2}^2+d\alpha ^2+4s^2d\mathrm{\Omega }_7^2),$$
(51)
where $`ds_{\mathrm{AdS}_2}^2`$ is the metric of AdS<sub>2</sub>, given by (23). The NS1/D0 system can be illustrated by the following diagram:
$`t`$ $`x_1`$ $`z_1`$ $`z_2`$ $`z_3`$ $`z_4`$ $`z_5`$ $`z_6`$ $`z_7`$ $`z_8`$ NS1 $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`H`$ D0 $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`W`$
Diagram 12. The NS1/D0 brane intersection.
In the M2/pp-wave and NS1/D0 systems, the internal space has a round 7-sphere. We can replace it by foliating of two lens spaces $`S^3/Z_{Q_\mathrm{N}}`$ and $`S^3/Z_{\stackrel{~}{Q}_\mathrm{N}}`$. As discussed in the appendix B, this can be achieved by introducing two NUTs in the intersecting system. We can then perform Kaluza-Klein reductions or T-duality transformations on the two associated fibre coordinates of the lens spaces. The resulting configurations can then be viewed as the near-horizon geometries of four intersecting $`p`$-branes, with charges $`Q`$, $`Q_w`$, $`Q_\mathrm{N}`$ and $`\stackrel{~}{Q}_\mathrm{N}`$
## 5 M5/NUT and M5/pp-wave systems
### 5.1 M5/NUT and NS5/D6 systems
The solution of an M5-brane intersecting with a NUT is given by
$`ds_{11}^2`$ $`=`$ $`H^{1/3}(dt^2+dw_1^2+\mathrm{}+dw_5^2)+H^{2/3}(dx_1^2+K(dz^2+z^2d\mathrm{\Omega }_2^2)+K^1(dy+\omega )^2),`$
$`F_{\left(4\right)}`$ $`=`$ $`(dtd^5wdH^1).`$ (52)
The solution can be illustrated by the following diagram:
$`t`$ $`w_1`$ $`w_2`$ $`w_3`$ $`w_4`$ $`w_5`$ $`x_1`$ $`z_1`$ $`z_2`$ $`z_3`$ $`y`$ M5 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $`H`$ NUT $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $`K`$
Diagram 13. The M5/NUT brane intersection.
In the usual construction where the harmonic functions $`H`$ and $`K`$ depend only the $`z`$ coordinate, the metric does not have an AdS structure in the near-horizon region. Here, we instead consider a semi-localised solution, given by
$$H=1+\frac{Q}{(x^2+4Q_\mathrm{N}z)^{3/2}},K=\frac{Q_\mathrm{N}}{z}.$$
(53)
After an analogous coordinate transformation, we find that the metric can be expressed as
$`ds_{11}^2`$ $`=`$ $`H^{1/3}(dt^2+dw_idw_i)+H^{2/3}(dr^2+r^2dM_4^2),`$
$`dM_4^2`$ $`=`$ $`d\alpha ^2+\frac{1}{4}s^2(d\mathrm{\Omega }_2^2+({\displaystyle \frac{dy}{Q_\mathrm{N}}}+\omega )^2).`$ (54)
Thus, in the near-horizon limit, the metric is AdS$`{}_{7}{}^{}\times M_4`$, where $`M_4`$ is a foliation of a lens space $`S^3/Z_{Q_\mathrm{N}}`$.
We can dimensionally reduce the solution (52) on the fibre coordinate $`y`$. The resulting solution is the NS-NS 5-brane intersecting with a D6-brane:
$`t`$ $`w_1`$ $`w_2`$ $`w_3`$ $`w_4`$ $`w_5`$ $`x_1`$ $`z_1`$ $`z_2`$ $`z_3`$ NS5 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $`H`$ D6 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $`K`$
Diagram 14. The NS5/D6 brane intersection.
The solution is given by
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2`$ $`=`$ $`H^{1/4}K^{1/8}(dt^2+dw_idw_i)+H^{3/4}K^{1/8}dx^2+H^{3/4}K^{7/8}dz_idz_i,`$
$`e^\varphi `$ $`=`$ $`H^{1/2}K^{3/4},F_{\left(3\right)}=e^{\varphi /2}(dtd^5wdH^1),`$
$`F_{\left(2\right)}`$ $`=`$ $`e^{3\varphi /2}(dtd^5wdxdK^1),`$ (55)
In the near-horizon limit, the metric becomes a warped product of AdS<sub>7</sub> with a 3-space
$$ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2=\frac{Q^{3/4}}{(2Q_\mathrm{N})^{1/4}}s^{1/4}(\frac{r}{Q}(dt^2+dw_idw_i)+\frac{dr^2}{r^2}+d\alpha ^2+\frac{1}{4}s^2d\mathrm{\Omega }_2^2).$$
(56)
### 5.2 M5/pp-wave and D0/D4 system
The solution of an M5-brane with a pp-wave is given by
$`ds_{11}^2`$ $`=`$ $`H^{1/3}(W^1dt^2+W(dy_1+(W^11)dt)^2+dx_1^2+\mathrm{}+dx_4^2)`$
$`+H^{2/3}(dz_1^2+\mathrm{}+dz_5^2),`$
$`F_4`$ $`=`$ $`(dtdy_1d^4xdH^1).`$ (57)
The solution can be illustrated by the following diagram:
$`t`$ $`y_1`$ $`x_1`$ $`x_2`$ $`x_3`$ $`x_4`$ $`z_1`$ $`z_2`$ $`z_3`$ $`z_4`$ $`z_5`$ M5 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $`H`$ wave $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`W`$
Diagram 15. The M5/pp-wave brane intersection.
We shall consider semi-localised solutions, with the functions $`H`$ and $`W`$ given by
$$H=\frac{Q}{z^3},W=1+Q_w(x^2+\frac{4Q}{z}).$$
(58)
Using analogous coordinate transformations, we find that the metric of the semi-localised M5/pp-wave system becomes
$$ds_{11}^2=4Q^{2/3}s^2(ds_{\mathrm{AdS}_3}^2+d\alpha ^2+c^2d\mathrm{\Omega }_3^2)+Q^{2/3}d\mathrm{\Omega }_4^2,$$
(59)
where $`ds_{\mathrm{AdS}_3}^2`$, given by (18), is precisely the extremal BTZ black hole and hence is is locally AdS<sub>3</sub>. After making the coordinate transformation $`\mathrm{tan}(\alpha /2)=e^\rho `$, the first part of the metric (59) can be expressed as
$$ds_7^2=d\rho ^2+\mathrm{sinh}^2\rho d\mathrm{\Omega }_3^2+\mathrm{cosh}^2\rho ds_3^2.$$
(60)
This is AdS<sub>7</sub> written as a foliation of AdS<sub>3</sub> and $`S^3`$.
Performing a dimensional reduction of the solution (57) on the coordinate $`y_1`$, we obtain a D0/D4 intersecting system, given by
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2`$ $`=`$ $`H^{3/8}W^{7/8}(dt^2+Wdx_idx_i+HWdz_idz_i),`$
$`e^\varphi `$ $`=`$ $`H^{1/4}W^{3/4},F_{\left(2\right)}=dtdW^1,`$
$`F_4`$ $`=`$ $`e^{\varphi /2}(dtd^4xdH^1).`$ (61)
The near-horizon limit of the semi-localised D0/D4 system is a warped product of AdS<sub>2</sub> with an 8-space:
$$ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2=2^{9/4}Q^{3/4}Q_w^{1/8}s^{9/4}(ds_2^2+d\alpha ^2+c^2d\mathrm{\Omega }_3^2+\frac{1}{4}s^2d\mathrm{\Omega }_4^2),$$
(62)
where $`ds_2^2`$ is given by (23). We illustrate this intersecting system with the following diagram
$`t`$ $`x_1`$ $`x_2`$ $`x_3`$ $`x_4`$ $`z_1`$ $`z_2`$ $`z_3`$ $`z_4`$ $`z_5`$ D4 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $`H`$ D0 $`\times `$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $``$ $`W`$
Diagram 16. The D0/D4 brane intersection.
In this example in the internal space the round $`S^3`$ and $`S^4`$ can be replaced by a lens space $`S^3/Z_{Q_\mathrm{N}}`$ and the foliation of a lens space $`S^3/Z_{\stackrel{~}{Q}_\mathrm{N}}`$, respectively. We can then perform Kaluza-Klein reductions or T-duality transformations on the fibre coordinates of the lens spaces, leading to four-component intersections with charges $`Q`$, $`Q_w`$, $`Q_\mathrm{N}`$ and $`\stackrel{~}{Q}_\mathrm{N}`$.
## 6 AdS<sub>6</sub> in type IIB from T-duality
So far in this paper we have two examples of intersecting D$`p`$/D$`(p+4)`$ systems in the type IIA theory that give rise to warped products of AdS<sub>p+2</sub> with certain internal spaces, namely for $`p=0`$ and $`p=2`$. It was observed also that the D4/D8 system, arising from massive type IIA supergravity, gives rise to the warped product of AdS<sub>6</sub> with a 4-sphere in the near-horizon limit:
$$ds_{10\mathrm{I}\mathrm{I}\mathrm{A}}^2=s^{1/12}(ds_{\mathrm{AdS}_6}^2+g^2(d\alpha ^2+c^2d\mathrm{\Omega }_3^2)).$$
(63)
Note that the D4/D8 system is less trivial than the previous examples, in the sense that it cannot be mapped by T-duality to a non-dilatonic $`p`$-brane intersecting with a NUT or a wave.
We can now introduce a NUT in the intersecting system which has the effect, in the near-horizon limit, of replacing the round 3-sphere by a lens space, given in (2). We can then perform a Hopf T-duality transformation and obtain an embedding of AdS<sub>6</sub> in type IIB theory:
$$ds_{10}^2=c^{1/2}\left[ds_{\mathrm{AdS}_6}^2+g^2(d\alpha ^2+\frac{1}{4}c^2d\mathrm{\Omega }_2^2)+s^{2/3}c^2dy^2\right].$$
(64)
This solution can be viewed as the near-horizon geometry of an intersecting D5/D7/NS5 system. It provides a background for the exact embedding of six-dimensional gauged supergravity in type IIB theory.
The D5/D7/NS5 semi-localised solution can be obtained by performing the T-duality on the D4/D8/NUT system. The solution is given by
$`ds_{10\mathrm{I}\mathrm{I}\mathrm{B}}^2`$ $`=`$ $`(H_1K)^{1/4}(dt^2+dw_1^2+\mathrm{}+dw_4^2+H_1K(dx_1^2+dx_2^2+dx_3^2)`$ (65)
$`+H_2Kdy^2+H_1H_2dz^2).`$
The functions $`H_1`$, $`H_2`$ and $`K`$ are given by
$$H_1=1+\frac{Q_1}{(4Q_\mathrm{N}|\stackrel{}{x}|+\frac{4Q_2}{9}z^3)^{5/3}},H_2=Q_2z,K=\frac{Q_\mathrm{N}}{|\stackrel{}{x}|}.$$
(66)
It is straightforward to verify that the near-horizon structure of this system is of the form (64). The solution can be illustrated by the following:
$`t`$ $`w_1`$ $`w_2`$ $`w_3`$ $`w_4`$ $`x_1`$ $`x_2`$ $`x_3`$ $`y`$ $`z`$ D5 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $``$ $`H_1`$ D7 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $`H_2`$ NS5 $`\times `$ $`\times `$ $`\times `$ $`\times `$ $`\times `$ $``$ $``$ $``$ $``$ $`\times `$ $`K`$
Diagram 17. The D5/D7/NS5 brane intersection.
## 7 Conclusion
In this paper, we obtain various AdS spacetimes warped with certain internal spaces in eleven-dimensional and type IIA/IIB supergravities. These solutions arise as the near-horizon geometries of more general semi-localised multi-intersections of M-branes in $`D=11`$ or NS-NS branes or D-branes in $`D=10`$. We achieve this by noting that any bigger sphere (AdS spacetime) can be viewed as a foliation involving $`S^3`$ (AdS<sub>3</sub>). Then the $`S^3`$ (AdS<sub>3</sub>) can be replaced by a three-dimensional lens space (BTZ black hole), which arise naturally from the introduction of a NUT (pp-wave). We can then perform a Kaluza-Klein reduction or Hopf T-duality transformation on the fibre coordinate of the lens space (BTZ black hole).
It is important to note that the warp factor depends only on the internal foliation coordinate but not on the lower-dimensional spacetime coordinates. This implies the possibility of finding a larger class of consistent dimensional reduction of eleven-dimensional or type IIA/IIB supergravity on the internal space, giving rise to gauged supergravities in lower dimensions with AdS vacuum solutions. The first such example was obtained in . In this paper, we obtain further examples for possible consistent embeddings of lower-dimensional gauged supergravity in $`D=11`$ and $`D=10`$. For example, we obtain the vacuum solutions for the embedding of the six and four-dimensional gauged AdS supergravities in type IIB theory and for the embedding of the seven-dimensional gauged AdS supergravity in type IIA theory.
## Acknowledgement
C.N.P. would like to thank the Caltech-USC Center for Theoretical Physics, and Imperial College, London, for hospitality during the course of this work.
## Appendix A Spheres and AdS from foliations
There are two closely parallel constructions which arise in the various intersections involving NUTs and waves. The former involves a construction of the unit metric on the sphere $`S^{p+q+1}`$ as a foliation of $`S^p\times S^q`$ surfaces, while the latter involves an analogous construction of the unit metric on AdS<sub>p+q+1</sub>, as a foliation of AdS$`{}_{p}{}^{}\times S^q`$ surfaces.
Consider first the construction of the unit $`S^{p+q+1}`$ metric. We start from the unit metrics $`d\mathrm{\Omega }_p^2=dX^idX^i`$ and $`d\mathrm{\Omega }_q^2=dY^adY^a`$ on the spheres $`S^p`$ and $`S^q`$, defined as the surfaces
$$X^iX^i=1,Y^aY^a=1$$
(67)
in $`\text{I}\mathrm{R}^{p+1}`$ and $`\text{I}\mathrm{R}^{q+1}`$ respectively. We now introduce Cartesian coordinates $`Z^A=(Z^i,Z^a)`$ in $`\text{I}\mathrm{R}^{p+q+2}`$, defined by
$$Z^i=X^i\mathrm{cos}\alpha ,Z^a=Y^a\mathrm{sin}\alpha ,$$
(68)
and so $`Z^AZ^A=1`$, thus defining a unit sphere $`S^{p+q+1}`$ in $`\text{I}\mathrm{R}^{p+q+2}`$. Clearly (68) defines a complete parameterisation of points in $`\text{I}\mathrm{R}^{p+q+2}`$, with $`0\alpha \frac{1}{2}\pi `$, and so $`\alpha `$, together with the constrained coordinates $`x^i`$ and $`y^a`$ on the spheres $`S^p`$ and $`S^q`$, provide coordinates for the unit sphere $`S^{p+q+1}`$ with a manifest $`SO(p+q+2)`$ isometry group action on the $`Z^A`$ coordinates. The metric on $`S^{p+q+1}`$ is given by $`d\mathrm{\Omega }_{p+q+1}^2=dZ^AdZ^A`$, and so from the above definitions we obtain
$$d\mathrm{\Omega }_{p+q+1}^2=d\alpha ^2+\mathrm{cos}^2\alpha d\mathrm{\Omega }_p^2+\mathrm{sin}^2\alpha d\mathrm{\Omega }_q^2.$$
(69)
The foliating surfaces at a fixed value of the “latitude” coordinate $`\alpha `$ are $`S^p\times S^q`$, with radii $`\mathrm{cos}\alpha `$ and $`\mathrm{sin}\alpha `$ for the two factors. The construction is a generalisation of the Clifford Torus $`S^1\times S^1`$ foliating $`S^3`$.
In a similar manner, one can construct a metric $`d\omega _{p+q+1}^2`$ on the unit AdS<sub>p+q+1</sub> as follows. We start from a unit AdS<sub>p</sub>, with metric $`d\omega _p^2=dX^\mu dX^\nu \eta _{\mu \nu }`$, and a unit $`S^q`$ with metric $`d\mathrm{\Omega }_q^2=dY^adY^a`$, where the coordinates $`X^\mu `$ on $`\text{I}\mathrm{R}^{p+1}`$ satisfy the indefinite-signature condition
$$X^\mu X^\nu \eta _{\mu \nu }=1,\eta _{\mu \nu }=\text{diag}(1,1,1,1,\mathrm{},1),$$
(70)
while the coordinates $`Y^a`$ on $`\text{I}\mathrm{R}^{q+1}`$ satisfy $`Y^aY^a=1`$ as before. We now define coordinates $`Z^A=(Z^\mu ,Z^a)`$ by
$$Z^\mu =X^\mu \mathrm{cosh}\rho ,Z^a=Y^a\mathrm{sinh}\rho ,$$
(71)
which therefore satisfy
$$Z^AZ^B\eta _{AB}=1,\eta _{AB}=\text{diag}(1,1,1,1,\mathrm{},1).$$
(72)
The coordinates $`Z^A`$, subject to this constraint, therefore define AdS<sub>p+q+1</sub>, with a manifest $`SO(p+q1,2)`$ isometry. The metric $`d\omega _{p+q+1}^2=dZ^AdZ^B\eta _{AB}`$ is given by
$$d\omega _{p+q+1}^2=d\rho ^2+\mathrm{cosh}^2\rho d\omega _p^2+\mathrm{sinh}^2\rho d\mathrm{\Omega }_q^2.$$
(73)
## Appendix B NUTs without NUTs
In this appendix, we show explicitly that the semi-localised intersection of a $`p`$-brane with a Kaluza-Klein monopole (a NUT) can be recast, after appropriate coordinate transformations, as a restricted class of ordinary distributed $`p`$-branes. For definiteness, we take the case of a semi-localised intersection of the M2-brane with a NUT as an example. The analysis for the other cases is essentially identical.
The semi-localised solution obtained in is given by
$`ds_{11}^2`$ $`=`$ $`H^{2/3}dw^\mu dw_\mu +H^{1/3}[(dx_1^2+\mathrm{}+dx_4^2)`$
$`+K(dz_1^2+dz_2^2+dz_3^2)+K^1(dy+A_idz_i)^2],`$
$`K`$ $`=`$ $`{\displaystyle \frac{Q_\mathrm{N}}{|\stackrel{}{z}|}},A_idz_i=Q_\mathrm{N}\mathrm{cos}\theta d\phi ,`$ (74)
$`H`$ $`=`$ $`1+{\displaystyle \underset{k}{}}{\displaystyle \frac{Q_k}{\left(|\stackrel{}{x}\stackrel{}{x}_{0k}|^2+4Q_\mathrm{N}|\stackrel{}{z}|\right)^3}},`$
where $`Q_k`$ denotes the M2-brane charge located at $`\stackrel{}{x}_{0k}`$, $`Q_\mathrm{N}`$ is the NUT charge, and we take
$$(z_1,z_2,z_3)=\frac{R^2}{4Q_\mathrm{N}}(\mathrm{sin}\theta \mathrm{cos}\phi ,\mathrm{sin}\theta \mathrm{sin}\phi ,\mathrm{cos}\theta ).$$
(75)
It now follows that the part of the metric
$$d\overline{s}^2K(dz_1^2+dz_2^2+dz_3^2)+K^1(dy+A_idz_i)^2$$
(76)
is nothing but the locally-flat metric
$$d\overline{s}^2=dR^2+R^2d\overline{\mathrm{\Omega }}_3^2,$$
(77)
where
$$d\overline{\mathrm{\Omega }}_3^2\frac{1}{4}d\mathrm{\Omega }_2^2+\frac{1}{4}\left(\frac{dy}{Q_\mathrm{N}}+\mathrm{cos}\theta d\phi \right)^2$$
(78)
is the metric on the cyclic lens space $`S^3/Z_{Q_\mathrm{N}}`$. Locally, this is just the standard metric on the unit 3-sphere. Viewed as a $`U(1)`$ bundle over $`S^2`$ the coordinate $`y`$ on the $`U(1)`$ fibres is taken always to have the period $`4\pi `$. When $`Q_\mathrm{N}=1`$, the topology is therefore precisely $`S^3`$. However, if $`Q_\mathrm{N}`$ is a larger integer, the fibre coordinate has a period that is smaller by the fraction $`1/Q_\mathrm{N}`$ than the period that would be needed for $`S^3`$ itself, and consequently the topology is $`S^3/Z_{Q_\mathrm{N}}`$.
The solution (74) can therefore be recast as
$$ds_{11}^2=H_2^{2/3}dw^\mu dw_\mu +H_2^{1/3}(dx_1^2+\mathrm{}+dx_4^2+d\stackrel{~}{z}_1^2+\mathrm{}+d\stackrel{~}{z}_4^2),$$
(79)
with the harmonic function given by
$$H_2=1+\underset{k}{}\frac{Q_k}{\left(|\stackrel{}{x}\stackrel{}{x}_{0k}|^2+|\stackrel{}{\stackrel{~}{z}}|^2\right)^3}.$$
(80)
The coordinates $`\stackrel{~}{z}_i`$ live on $`\text{I}\mathrm{R}^4/Z_{Q_\mathrm{N}}`$, and are related to $`R`$ and the coordinates $`(\theta ,\phi ,y)`$ on the lens space $`S^3/Z_{Q_\mathrm{N}}`$ by
$$\stackrel{~}{z}_1+\mathrm{i}\stackrel{~}{z}_2=R\mathrm{sin}\frac{1}{2}\theta e^{\frac{\mathrm{i}}{2}(y/Q_\mathrm{N}+\phi )},\stackrel{~}{z}_3+\mathrm{i}\stackrel{~}{z}_4=R\mathrm{cos}\frac{1}{2}\theta e^{\frac{\mathrm{i}}{2}(y/Q_\mathrm{N}\phi )}.$$
(81)
In other words, if we make the following coordinate transformation from $`(z_1,z_2,z_3,y)`$ to $`(\stackrel{~}{z}_1,\stackrel{~}{z}_2,\stackrel{~}{z}_3,\stackrel{~}{z}_4)`$,
$`\stackrel{~}{z}_1+\mathrm{i}\stackrel{~}{z}_2`$ $`=`$ $`\left[{\displaystyle \frac{2Q_\mathrm{N}(r+z_3)(z_1+\mathrm{i}z_2)}{\sqrt{z_1^2+z_2^2}}}\right]^{1/2}e^{\frac{\mathrm{i}}{2Q_\mathrm{N}}y},`$
$`\stackrel{~}{z}_3+\mathrm{i}\stackrel{~}{z}_4`$ $`=`$ $`\left[{\displaystyle \frac{2Q_\mathrm{N}(rz_3)(z_1\mathrm{i}z_2)}{\sqrt{z_1^2+z_2^2}}}\right]^{1/2}e^{\frac{\mathrm{i}}{2Q_\mathrm{N}}y},`$ (82)
where $`r^2z_1^2+z_2^2+z_3^2`$, then the metric (76) is seen to be nothing but
$$d\overline{s}^2=d\stackrel{~}{z}_1^2+d\stackrel{~}{z}_2^2+d\stackrel{~}{z}_3^2+d\stackrel{~}{z}_4^2.$$
(83)
The semi-localised M2-brane/NUT intersection (74) can therefore be obtained by starting from a standard distribution of pure M2-branes (79), with charges spread over only four of the eight transverse directions as in (80). This is precisely equivalent to the semi-localised M2-brane/NUT intersection (74) with unit NUT charge, $`Q_\mathrm{N}=1`$. To obtain higher values of the NUT charge, one simply has to factor the $`\text{I}\mathrm{R}^4`$ space of the $`\stackrel{~}{z}_i`$ coordinates by $`Z_{Q_\mathrm{N}}`$, as defined above. Note that although this semi-localised way of introducing a NUT seems trivial, in that it amounts a coordinate transformation, performing Kaluza-Klein reduction on the fibre coordinate does create a non-trivial intersecting component, since the Kaluza-Klein 2-form field strength now carries a non-trivial flux.
The above discussion carries over, mutatis mutandis, to the cases of the semi-localised M5-brane/NUT and D3-brane/NUT.
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# ARPES study of Pb doped 𝐵𝑖₂𝑆𝑟₂𝐶𝑎𝐶𝑢₂𝑂₈ - a new Fermi surface picture
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## Abstract
High resolution angle resolved photoemission data from Pb doped $`Bi_2Sr_2CaCu_2O_8`$ (Bi2212) with suppressed superstructure is presented. Improved resolution and very high momentum space sampling at various photon energies reveal the presence of two Fermi surface pieces. One has the hole-like topology, while the other one has its van Hove singularity very close to $`(\pi ,0)`$, its topology at some photon energies resembles the electron-like piece. This result provides a unifying picture of the Fermi surface in the Bi2212 compound and reconciles the conflicting reports.
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The Fermi surface plays an important role in understanding the physics of any material. Among other things its shape and size determine the type and number of charge carriers in the material as well as the charge and spin dynamics. For example, in the context of the Fermi liquid approach to high temperature superconductors (HTSCs), Fermi surface topology is related to commensurate or incommensurate nature of the neutron data . Furthermore, detailed knowledge of the Fermi surface is essential to determine the superconducting gap size and symmetry in the superconducting state.
Angle resolved photoemission spectroscopy (ARPES) is a unique tool to probe the Fermi surface of the HTSCs. Over the last decade the HTSC system most extensively investigated by ARPES is Bi2212 . However, the existence of superstructure in the $`BiO`$ layer and shadow bands has made the Fermi surface determination in this compound complicated, especially around the $`M(\pi ,0)`$ point, where main bands, superstructure bands, and shadow bands cross the Fermi level. For several years there was a general agreement for a hole-like Fermi surface centered around $`Y(\pi ,\pi )`$, mainly based on ARPES experiments performed at $`22eV`$ photon energy . This hole-like Fermi surface picture was believed to apply over the entire doping range studied by ARPES (from underdoped samples with $`Tc15K`$ to overdoped samples with $`Tc68K`$). Later a vigorous discussion started, with experiments utilizing $`33eV`$ photon energy suggesting an electron-like Fermi surface centered around the $`\mathrm{\Gamma }`$ point . Other groups disputed electron-like Fermi surface reports, dismissing the observed results by invoking the interplay of matrix element effects with $`BiO`$ layer superstructure and shadowbands . Finally, recent reports on Bi2212 using 22 eV photons demonstrated the presence of the two Fermi surfaces in the material due to bonding and antibonding interaction of CuO planes . It is of great importance to reconcile all the reported results and to resolve the uncertainty in the Fermi surface of Bi2212 - the compound extensively studied and the source of many significant results.
Because the main discrepancy originates in the $`(\pi ,0)`$ region, where the superstructure effect of the $`BiO`$ layer is strongest, a definitive resolution of the Fermi surface issue can be found by studying the Pb-doped Bi2212 system. In this compound Pb is doped into the $`BiO`$ plane disrupting the $`BiO`$ plane modulation and removing the superstructure complication near the $`(\pi ,0)`$ region .
In this Letter we present results of Fermi surface mapping of Pb-doped Bi2212 with high energy resolution and very high $`k`$ space sampling. We used various photon energies and different methods to determine the Fermi surface. Our 22 eV data complement recent reports on the existence of two Fermi surface pieces. The photon energy dependence reveals that the relative intensity depends strongly on photon energy. While the bonding Fermi surface has a clear hole-like topology, the antibonding piece has its van Hove singularity very close to $`(\pi ,0)`$ and its Fermi surface is electron-like as seen at some photon energies. This result contradicts earlier reports of a single universal Fermi surface in this compound . On the other hand our data provides a unifying foundation for understanding the controversies about the Fermi surface of this important superconductor, as different reports stress different aspects of the global Fermi surface features.
ARPES data have been recorded at beamline $`\mathrm{10.0.1.1}`$ of the Advanced Light Source utilizing $`55,44,33,27`$ and $`22`$ $`eV`$ photon energy in $`410^{}{}_{}{}^{1}^1`$ $`Torr`$ vacuum. The sample was kept in the fixed position relative to the beam polarization, and the analyzer was rotated. The beam polarization was in the sample plane perpendicular to $`\mathrm{\Gamma }Y`$ direction, with beam nearly at grazing incidence with the sample surface. We used a Scienta SES 200 analyzer in the angle mode, where cuts parallel to $`\mathrm{\Gamma }Y`$ direction are carried out. The momentum resolution was $`\pm 0.06\AA ^1`$ in the scan direction and $`\pm 0.19\AA ^1`$ in the perpendicular direction for 55 eV photon energy and better for other energies, and the energy resolution was $`718`$ $`meV`$. An extensive and fine sampling mesh with more than 4000 EDCs for each photon energy was collected. The slightly overdoped Pb-doped Bi2212 ($`Tc=84K`$) and overdoped Pb-doped Bi2212 ($`Tc=70K`$) were grown using the floating-zone method. The single crystalline samples were oriented by using Laue diffraction ex situ and cleaved in situ in vacuum. The samples were measured at 100K (Fig. 1,2 and 4) and 20K (Fig. 3). The Fermi energy was obtained from the EDCs of polycrystalline $`Au`$.
In panel a) of Fig. 1 we show the map of spectral intensity at Fermi energy ($`E_F`$) obtained at 22 eV photon energy. The white arrow shows the polarization of radiation with respect to the crystal surface. To determine the spectral intensity map we divide each EDC by the integrated signal intensity from a 100 meV window above the $`E_F`$, which comes from higher order synchrotron light and is proportional to photon flux. The normalized EDCs represent electron spectral function weighted by the Fermi function and matrix element. Highest intensity points in the spectral intensity map at the Fermi energy give one method for determining Fermi surface. In panels 1b)-1f) we plot raw EDC data obtained along select cuts shown in panel a) by green arrows. Here Fermi surface crossing is defined as the location in the momentum space where the intensity of the spectral feature decreases drastically and the leading edge crosses the Fermi level. From both the near $`E_F`$ spectral weight image plot and from the EDCs one clearly sees two Fermi surfaces, as indicated by solid ovals in panels b to f. Concomitant presence of two Fermi surface pieces is consistent with recent data recorded at a similar photon energy
Fig. 2a) shows 22 eV data from another sample taken in a more extended k-space area. This data were taken at 20K in the superconducting state, where the spectral weight around M point is suppressed at the Fermi level due to the sc gap opening, so the map shown corresponds to 12 meV BE. This map can effectively be used to indicate the underlying Fermi surface. Although the maximum gap is larger than the energy window, finite resolution still reveals the underlying Fermi surface, and larger integration window does not change the picture. This data, as well as the normal state data in panel a) of Fig. 1, show striking resemblance to theoretical simulation by Bansil et al. for the same photon energy and polarization, as shown in panel b). The simulation uses first-principles one-step photoemission model calculation and comes up with two bands for two adjacent $`CuO_2`$ planes in a unit cell. These two bonding and antibonding bands give rise to two Fermi surfaces. The bonding piece is an outer hole-like piece, indicated by blue lines in the inset of panel a). On the other hand the antibonding piece, indicated by the red area in the inset of panel a), is hard to judge from these data, because the saddle point of the band is very close to $`E_F`$ at $`(\pi ,0)`$ . The simulation also indicates that the image plot is very similar whether this piece of the Fermi surface is hole-like or electron-like, i.e. whether the Van Hove saddle point is above or below $`E_F`$. The absence of superstructure complication in our data and the striking similarity between experiment and theory strongly suggest that there are indeed two pieces of the Fermi surface in Pb Bi2212.
The Fermi surface seen under other measurement conditions turns out to be very different. Fig. 3 shows data recorded at 55 eV under the same measurement geometry. Panel a) shows the map of spectral intensity at $`E_F`$ in the momentum space. In panels b), c) and d) we plot raw data obtained along the select cuts shown in panel a) by thick white lines. Cut b) is close to the nodal direction in the second Brillouin zone, while cuts c) and d) are cuts equidistant from $`M`$ point. The sampling density in the cuts is very high and is representative of the sampling density of the entire k-space studied. The high quality data clearly shows a quasiparticle dispersing towards the Fermi level, eventually crossing it and disappearing. While the intensity map in panel a) hardly shows any features in the first zone, panel d) clearly shows a well-resolved feature crossing the Fermi level, similar to panel c), with the overall intensity a factor of 10 lower than that in panel c). In fact, all features seen in the second zone are observed in the first zone as well. The Fermi surface shape emerging for this photon energy is electron-like. We have confirmed this result with data taken in all three Brillouin zones by using three complimentary methods: intensity map at $`E_F`$, the traditional method of tracking the EDCs, and the sharpest drop in $`n(\stackrel{}{k})`$. The contrast in data from Figures 1-3 immediately suggests that the 55 eV data picks out the inner piece of the Fermi surface.
To investigate the photon energy dependence of the FS further, we collected data at other photon energies and in Fig. 4 we plot the measured spectral intensity maps at $`E_f`$. In panel a) we plot our spectral intensity map collected at 27 eV. We see an electron-like FS in the first zone, and the spectral weight in the second zone is significantly suppressed. This complements the 55 eV data in panel d). In panel b) at 33 eV photon energy spectral intensity map shows strong suppression of the spectral weight at the M point and can be interpreted as either hole-like and electron-like. This data is quite different from earlier results recorded at another geometry . 44 eV data in panel c) looks very similar to 22 eV data in panel a) with bilayer split Fermi surfaces. Spectral weight map FS results are supported by individual EDC analysis.
There’s a very unusual photon energy related variation in the ARPES data. With the large unit cell size in the normal direction ($`c30.6\AA `$), one would expect periodicity with 1-2 eV steps in photon energy for the energy range studied. Surprisingly, this is not the case. Existence of bilayer split Fermi surface means strong interaction between the two $`CuO_2`$ planes in Bi2212 unit cell. It is reasonable to assume that the variation in the data with the photon energy is also driven by the separation of the layers $`d3.4\AA `$, with corresponding $`Kz=\frac{2\pi }{d}=1.85\AA `$. This would be more consistent with the observed variations with large photon energy intervals. Calculations similar to for a range of photon energies will be very useful to understand this phenomena.
Our results, in particular the variation of the Fermi surface picture with photon energies and the observation of the bilayer splitting, contradict earlier reports from Pb doped Bi2212. Previous data from the same material at different photon energies were interpreted as an evidence for the universal hole-like Fermi surface. We attribute this discrepancy mostly to poorer energy (70 meV compared to ours 7-18 meV) and momentum ($`0.094\AA ^1`$ compared to ours $`0.006\times 0.019\AA ^1`$) resolution used in that study for photon energies other then 22.4 eV. Attempts were also made to support a single universal hole-like Fermi surface picture by invoking matrix element arguments. While a band calculation did show a strong spectral intensity variation with photon energy , it is not applicable to Fermi surface determination at a particular photon energy. What really matters is the matrix element variation in a small region in the same Brillouin zone, which is actually small . The unambiguous evidence for two Fermi surface pieces show the matrix element argument for single universal hole-like Fermi surface to be misleading and incorrect.
The polarization setup used in our experiment is favorable for observing the bilayer splitting . Observation of the splitting of the bands due to the interaction between the layers was a long standing problem. Earlier photoemission results were interpreted as evidence for the absence of the bilayer splitting . Our data in Fig. 1, 2 and 4, collected with much better energy and momentum resolution, show the splitting to be present, even for the samples with Tc not far from optimal.
The picture emerging from the above is clear: there are two pieces of the Fermi surface, one of them clearly hole-like. The other piece is different, as it lies very close to the special point where Fermi surface changes from hole-like to electron-like with small change in chemical potential or $`\stackrel{}{K}_Z`$. This is the underlying reason for this piece to behave slightly differently at different photon energies. This picture provides a unifying foundation for all the controversial reports on the FS shape in Bi2212. Bi2212 Fermi surface has always been attributed to the $`CuO_2`$ plane, and because doping with Pb does not change the $`CuO_2`$ plane, the Fermi surface in Pb-doped and Pb-free compounds should be the same. Our finding indicates that the accepted picture of a single hole-like Fermi surface in Bi2212 for the entire doping range studied is incorrect. Instead, there exists another Fermi surface piece that is at the boundary between hole and electron character.
We would like to thank J. D. Denlinger for the help with data analysis software. The experiment was performed at the Advanced Light Source of Lawrence Berkeley National Laboratory, supported by DOE’ Office of Basic Energy Science, Division of Materials Science with contract DE-AC03-76SF00098. The Stanford work was supported by NSF grant through the Stanford MRSEC grant and NSF grant DMR-9705210. The SSRL’s work was also supported by the Office’s Division of Materials Science.
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# 1 Introduction
## 1 Introduction
Differential calculi can be constructed on spaces that are more general than differentiable manifolds. Indeed the algebraic construction of differential calculus in terms of Hopf structures allows to extend the usual differential geometric quantities (connection, curvature, metric, vielbein etc.) to a variety of interesting spaces that include quantum groups, noncommutative spacetimes (i.e. quantum cosets), and discrete spaces.
In this contribution we concentrate on the differential geometry of finite group “manifolds”. As we will discuss, these spaces can be visualized as collections of points, corresponding to the finite group elements, and connected by oriented links according to the particular differential calculus we build on them. Although functions $`fFun(G)`$ on finite groups $`G`$ commute, the calculi that are constructed on $`Fun(G)`$ by algebraic means are in general noncommutative, in the sense that differentials do not commute with functions, and the exterior product does not coincide with the usual antisymmetrization of the tensor product.
Among the physical motivations for finding differential calculi on finite groups we mention the possibility of using finite group spaces as internal spaces for Kaluza-Klein compactifications of Yang-Mills, (super)gravity or superstring theories ( for example Connes’ reconstruction of the standard model in terms of noncommutative geometry can be recovered as Kaluza-Klein compactification of Yang-Mills theory on an appropriate discrete internal space). Differential calculi on discrete spaces can be of use in the study of integrable models, see for ex. ref. . Finally gauge and gravity theories on finite group spaces may be used as lattice approximations. For example the action for pure Yang-Mills $`F{}_{}{}^{}F`$ considered on the finite group space $`Z^N\times Z^N\times Z^N\times Z^N`$, yields the usual Wilson action of lattice gauge theories, and $`N\mathrm{}`$ gives the continuum limit . New lattice theories can be found by choosing different finite groups.
A brief review of the differential calculus on finite groups is presented. Most of this material is not new, and draws on the treatment of ref.s , where the Hopf algebraic approach of Woronowicz for the construction of differential calculi is adapted to the setting of finite groups. Some developments on Lie derivative, diffeomorphisms and integration are new. The general theory is illustrated in the case of $`S_3`$.
## 2 Differential calculus on finite groups
Let $`G`$ be a finite group of order $`n`$ with generic element $`g`$ and unit $`e`$. Consider $`Fun(G)`$, the set of complex functions on $`G`$. An element $`f`$ of $`Fun(G)`$ is specified by its values $`f_gf(g)`$ on the group elements $`g`$, and can be written as
$$f=\underset{gG}{}f_gx^g,f_g𝑪$$
(2.1)
where the functions $`x^g`$ are defined by
$$x^g(g^{})=\delta _g^{}^g$$
(2.2)
Thus $`Fun(G)`$ is a n-dimensional vector space, and the $`n`$ functions $`x^g`$ provide a basis. $`Fun(G)`$ is also a commutative algebra, with the usual pointwise sum and product \[$`(f+h)(g)=f(g)+h(g)`$, $`(fh)(g)=f(g)h(g)`$, $`(\lambda f)(g)=\lambda f(g),f,hFun(G),\lambda 𝑪`$\] and unit $`I`$ defined by $`I(g)=1,gG`$. In particular:
$$x^gx^g^{}=\delta _{g,g^{}}x^g,\underset{gG}{}x^g=I$$
(2.3)
Consider now the left multiplication by $`g_1`$:
$$L_{g_1}g_2=g_1g_2,g_1,g_2G$$
(2.4)
This induces the left action (pullback) $`_{g_1}`$ on $`Fun(G)`$:
$$_{g_1}f(g_2)f(g_1g_2)|_{g_2},_{g_1}:Fun(G)Fun(G)$$
(2.5)
where $`f(g_1g_2)|_{g_2}`$ means $`f(g_1g_2)`$ seen as a function of $`g_2`$. Similarly we can define the right action on $`Fun(G)`$ as:
$$(_{g_1}f)(g_2)=f(g_2g_1)|_{g_2}$$
(2.6)
For the basis functions we find easily:
$$_{g_1}x^g=x^{g_1^1g},_{g_1}x^g=x^{gg_1^1}$$
(2.7)
Moreover:
$`_{g_1}_{g_2}=_{g_1g_2},_{g_1}_{g_2}=_{g_2g_1},`$ (2.8)
$`_{g_1}_{g_2}=_{g_2}_{g_1}`$ (2.9)
Bicovariant differential calculus
Differential calculi can be constructed on Hopf algebras $`A`$ by algebraic means, using the costructures of $`A`$ . In the case of finite groups $`G`$, differential calculi on $`A=Fun(G)`$ have been discussed in ref.s . Here we give the main results derived in , to which we refer for a more detailed treatment.
A first-order differential calculus on $`A`$ is defined by
i) a linear map $`d`$: $`A\mathrm{\Gamma }`$, satisfying the Leibniz rule
$$d(ab)=(da)b+a(db),a,bA;$$
(2.10)
The “space of 1-forms” $`\mathrm{\Gamma }`$ is an appropriate bimodule on $`A`$, which essentially means that its elements can be multiplied on the left and on the right by elements of $`A`$ \[more precisely $`A`$ is a left module if $`a,bA,\rho ,\rho ^{}\mathrm{\Gamma }`$ we have: $`a(\rho +\rho ^{})=a\rho +a\rho ^{},(a+b)\rho =a\rho +b\rho ,a(b\rho )=(ab)\rho ,I\rho =\rho `$. Similarly one defines a right module. A left and right module is a bimodule if $`a(\rho b)=(a\rho )b`$\]. From the Leibniz rule $`da=d(Ia)=(dI)a+Ida`$ we deduce $`dI=0`$.
ii) the possibility of expressing any $`\rho \mathrm{\Gamma }`$ as
$$\rho =\underset{k}{}a_kdb_k$$
(2.11)
for some $`a_k,b_k`$ belonging to $`A`$.
To build a first order differential calculus on $`Fun(G)`$ we need to extend the algebra $`A=Fun(G)`$ to a differential algebra of elements $`x^g,dx^g`$ (it is sufficient to consider the basis elements and their differentials). Note however that the $`dx^g`$ are not linearly independent. In fact from $`0=dI=d(_{gG}x^g)=_{gG}dx^g`$ we see that only $`n1`$ differentials are independent. Every element $`\rho =adb`$ of $`\mathrm{\Gamma }`$ can be expressed as a linear combination (with complex coefficients) of terms of the type $`x^gdx^g^{}`$. Moreover $`\rho b\mathrm{\Gamma }`$ (i.e. $`\mathrm{\Gamma }`$ is also a right module) since the Leibniz rule and the multiplication rule (2.3) yield the commutations:
$$dx^gx^g^{}=x^gdx^g^{}+\delta _g^{}^gdx^g$$
(2.12)
allowing to reorder functions to the left of differentials.
Partial derivatives
Consider the differential of a function $`fFun(g)`$:
$$df=\underset{gG}{}f_gdx^g=\underset{ge}{}f_gdx^g+f_edx^e=\underset{ge}{}(f_gf_e)dx^g\underset{ge}{}_gfdx^g$$
(2.13)
We have used $`dx^e=_{ge}dx^g`$ (from $`_{gG}dx^g=0`$). The partial derivatives of $`f`$ have been defined in analogy with the usual differential calculus, and are given by
$$_gf=f_gf_e=f(g)f(e)$$
(2.14)
Not unexpectedly, they take here the form of finite differences (discrete partial derivatives at the origin $`e`$).
Left and right covariance
A differential calculus is left or right covariant if the left or right action of $`G`$ ($`_g`$ or $`_g`$) commutes with the exterior derivative $`d`$. Requiring left and right covariance in fact defines the action of $`_g`$ and $`_g`$ on differentials: $`_gdbd(_gb),bFun(G)`$ and similarly for $`_gdb`$. More generally, on elements of $`\mathrm{\Gamma }`$ (one-forms) we define $`_g`$ as:
$$_g(adb)(_ga)_gdb=(_ga)d(_gb)$$
(2.15)
and similar for $`_g`$. Computing for example the left and right action on the differentials $`dx^g`$ yields:
$$_g(dx^{g_1})d(_gx^{g_1})=dx^{g^1g_1},_g(dx^{g_1})d(_gx^{g_1})=dx^{g_1g^1}$$
(2.16)
A differential calculus is called bicovariant if it is both left and right covariant.
Left invariant one forms
As in usual Lie group manifolds, we can introduce a basis in $`\mathrm{\Gamma }`$ of left-invariant one-forms $`\theta ^g`$:
$$\theta ^g\underset{hG}{}x^{hg}dx^h(=\underset{hG}{}x^hdx^{hg^1}),$$
(2.17)
It is immediate to check that $`_k\theta ^g=\theta ^g`$. The relations (2.17) can be inverted:
$$dx^h=\underset{gG}{}(x^{hg}x^h)\theta ^g$$
(2.18)
From $`0=dI=d_{gG}x^g=_{gG}dx^g=0`$ one finds:
$$\underset{gG}{}\theta ^g=\underset{gG}{}\underset{hG}{}x^hdx^{hg^1}=\underset{hG}{}x^h\underset{gG}{}dx^{hg^1}=0$$
(2.19)
Therefore we can take as basis of the cotangent space $`\mathrm{\Gamma }`$ the $`n1`$ linearly independent left-invariant one-forms $`\theta ^g`$ with $`ge`$ (but smaller sets of $`\theta ^g`$ can be consistently chosen as basis, see later).
The commutations between the basic 1-forms $`\theta ^g`$ and functions $`fFun(G)`$ are given by:
$$f\theta ^g=\theta ^g_gf$$
(2.20)
Thus functions do commute between themselves (i.e. $`Fun(G)`$ is a commutative algebra) but do not commute with the basis of one-forms $`\theta ^g`$. In this sense the differential geometry of $`Fun(G)`$ is noncommutative, the noncommutativity being milder than in the case of quantum groups $`Fun_q(G)`$(which are noncommutative algebras).
The right action of $`G`$ on the elements $`\theta ^g`$ is given by:
$$_h\theta ^g=\theta ^{ad(h)g},hG$$
(2.21)
where $`ad`$ is the adjoint action of $`G`$ on $`G`$, i.e. $`ad(h)ghgh^1`$. Then bicovariant calculi are in 1-1 correspondence with unions of conjugacy classes (different from $`\{e\}`$) : if $`\theta ^g`$ is set to zero, one must set to zero all the $`\theta ^{ad(h)g},hG`$ corresponding to the whole conjugation class of $`g`$.
We denote by $`G^{}`$ the subset corresponding to the union of conjugacy classes that characterizes the bicovariant calculus on $`G`$ ($`G^{}=\{gG|\theta ^g0\}`$). Unless otherwise indicated, repeated indices are summed on $`G^{}`$ in the following.
A bi-invariant (i.e. left and right invariant) one-form $`\mathrm{\Theta }`$ is obtained by summing on all $`\theta ^g`$ with $`ge`$:
$$\mathrm{\Theta }=\underset{ge}{}\theta ^g$$
(2.22)
Exterior product
For a bicovariant differential calculus on a Hopf algebra $`A`$ an exterior product, compatible with the left and right actions of $`G`$, can be defined by
$$\theta ^{g_1}\theta ^{g_2}=\theta ^{g_1}\theta ^{g_2}\theta ^{g_1^1g_2g_1}\theta ^{g_1}$$
(2.23)
where the tensor product between elements $`\rho ,\rho ^{}\mathrm{\Gamma }`$ is defined to have the properties $`\rho a\rho ^{}=\rho a\rho ^{}`$, $`a(\rho \rho ^{})=(a\rho )\rho ^{}`$ and $`(\rho \rho ^{})a=\rho (\rho ^{}a)`$.
Note that:
$$\theta ^g\theta ^g=0\text{(no sum on }g\text{)}$$
(2.24)
Left and right actions on $`\mathrm{\Gamma }\mathrm{\Gamma }`$ are simply defined by:
$$_h(\rho \rho ^{})=_h\rho _h\rho ^{},_h(\rho \rho ^{})=_h\rho _h\rho ^{}$$
(2.25)
(with the obvious generalization to $`\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma }`$) so that for example:
$$_h(\theta ^i\theta ^j)=\theta ^i\theta ^j,_h(\theta ^i\theta ^j)=\theta ^{ad(h)i}\theta ^{ad(h)j}$$
(2.26)
We can generalize the definition (2.28)to exterior products of $`n`$ one-forms:
$$\theta ^{i_1}\mathrm{}\theta ^{i_n}W_{j_1k_1}^{i_1i_2}W_{j_2k_2}^{k_1i_3}W_{j_3k_3}^{k_2i_4}\mathrm{}W_{j_{n1}j_n}^{k_{n2}i_n}\theta ^{j_1}\mathrm{}\theta ^{j_n}$$
(2.27)
where the matrix $`W`$ is defined by:
$$\theta ^i\theta ^jW_{kl}^{ij}\theta ^k\theta ^l=\theta ^i\theta ^j\mathrm{\Lambda }_{kl}^{ij}\theta ^k\theta ^l.$$
(2.28)
and $`\mathrm{\Lambda }_{kl}^{ij}`$ is the braiding matrix defined by (2.23). The space of $`n`$-forms $`\mathrm{\Gamma }^n`$ is therefore defined as in the usual case but with the new permutation operator $`\mathrm{\Lambda }`$, and can be shown to be a bicovariant bimodule, with left and right action defined as for $`\mathrm{\Gamma }\mathrm{}\mathrm{\Gamma }`$ with the tensor product replaced by the wedge product.
Exterior derivative
Having the exterior product we can define the exterior derivative
$$d:\mathrm{\Gamma }\mathrm{\Gamma }\mathrm{\Gamma }$$
(2.29)
$$d(a_kdb_k)=da_kdb_k,$$
(2.30)
which can easily be extended to $`\mathrm{\Gamma }^n`$ ($`d:\mathrm{\Gamma }^n\mathrm{\Gamma }^{(n+1)}`$), and has the following properties:
$$d(\rho \rho ^{})=d\rho \rho ^{}+(1)^k\rho d\rho ^{}$$
(2.31)
$$d(d\rho )=0$$
(2.32)
$$_g(d\rho )=d_g\rho $$
(2.33)
$$_g(d\rho )=d_g\rho $$
(2.34)
where $`\rho \mathrm{\Gamma }^k`$, $`\rho ^{}\mathrm{\Gamma }^n`$. The last two properties express the fact that $`d`$ commutes with the left and right action of $`G`$.
Tangent vectors
Using (2.18) to expand $`df`$ on the basis of the left-invariant one-forms $`\theta ^g`$ defines the (left-invariant) tangent vectors $`t_g`$:
$$df=\underset{gG}{}f_gdx^g=\underset{hG^{}}{}(_{h^1}ff)\theta ^h\underset{hG^{}}{}(t_hf)\theta ^h$$
(2.35)
so that the “flat” partial derivatives $`t_hf`$ are given by
$$t_hf=_{h^1}ff$$
(2.36)
The Leibniz rule for the flat partial derivatives $`t_g`$ reads:
$$t_g(ff^{})=(t_gf)_{g^1}f^{}+ft_gf^{}$$
(2.37)
In analogy with ordinary differential calculus, the operators $`t_g`$ appearing in (2.35) are called (left-invariant) tangent vectors, and in our case are given by
$$t_g=_{g^1}id$$
(2.38)
They satisfy the composition rule:
$$t_gt_g^{}=\underset{h}{}C_{g,g^{}}^ht_h$$
(2.39)
where the structure constants are:
$$C_{g,g^{}}^h=\delta _{g^{}g}^h\delta _g^h\delta _g^{}^h$$
(2.40)
and have the property:
$$C_{ad(h)g_2,ad(h)g_3}^{ad(h)g_1}=C_{g_2,g_3}^{g_1}$$
(2.41)
Note 2.1 : The exterior derivative on any $`fFun(G)`$ can be expressed as a commutator of $`f`$ with the bi-invariant one-form $`\mathrm{\Theta }`$:
$$df=[\mathrm{\Theta },f]$$
(2.42)
as one proves by using (2.20) and (2.35).
Note 2.2 : From the fusion rules (2.39) we deduce the “deformed Lie algebra” (cf. ref.s ):
$$t_{g_1}t_{g_2}\mathrm{\Lambda }_{g_1,g_2}^{g_3,g_4}t_{g_3}t_{g_4}=𝑪_{g_1,g_2}^ht_h$$
(2.43)
where the $`𝑪`$ structure constants are given by:
$$𝑪_{g_1,g_2}^gC_{g_1,g_2}^g\mathrm{\Lambda }_{g_1,g_2}^{g_3,g_4}C_{g_3,g_4}^g=C_{g_1,g_2}^gC_{g_2,g_2g_1g_2^1}^g=\delta _{g_1}^{ad(g_2^1)g}\delta _{g_1}^g$$
(2.44)
and besides property (2.41) they also satisfy:
$$𝑪_{g_1,g_2}^g=𝑪_{g,g_2^1}^{g_1}$$
(2.45)
Moreover the following identities hold:
i) deformed Jacobi identities:
$$𝑪_{h_1,g_1}^k𝑪_{k,g_2}^{h_2}\mathrm{\Lambda }_{g_1,g_2}^{g_3,g_4}𝑪_{h_1,g_3}^k𝑪_{k,g_4}^{h_2}=𝑪_{g_1,g_2}^k𝑪_{h_1,k}^{h_2}$$
(2.46)
ii) fusion identities:
$$𝑪_{h_1,g}^k𝑪_{k,g^{}}^{h_2}=C_{g,g^{}}^h𝑪_{h_1,h}^{h_2}$$
(2.47)
Thus the $`𝑪`$ structure constants are a representation (the adjoint representation) of the tangent vectors $`t`$.
Cartan-Maurer equations, connection and curvature
From the definition (2.17) and eq. (2.20) we deduce the Cartan-Maurer equations:
$$d\theta ^g+\underset{g_1,g_2}{}C_{g_1,g_2}^g\theta ^{g_1}\theta ^{g_2}=0$$
(2.48)
where the structure constants $`C_{g_1,g_2}^g`$ are those given in (2.40).
Parallel transport of the vielbein $`\theta ^g`$ can be defined as in ordinary Lie group manifolds:
$$\theta ^g=\omega _g^{}^g\theta ^g^{}$$
(2.49)
where $`\omega _{g_2}^{g_1}`$ is the connection one-form:
$$\omega _{g_2}^{g_1}=\mathrm{\Gamma }_{g_3,g_2}^{g_1}\theta ^{g_3}$$
(2.50)
Thus parallel transport is a map from $`\mathrm{\Gamma }`$ to $`\mathrm{\Gamma }\mathrm{\Gamma }`$; by definition it must satisfy:
$$(a\rho )=(da)\rho +a\rho ,aA,\rho \mathrm{\Gamma }$$
(2.51)
and it is a simple matter to verify that this relation is satisfied with the usual parallel transport of Riemannian manifolds. As for the exterior differential, $``$ can be extended to a map $`:\mathrm{\Gamma }^n\mathrm{\Gamma }\mathrm{\Gamma }^{(n+1)}\mathrm{\Gamma }`$ by defining:
$$(\phi \rho )=d\phi \rho +(1)^n\phi \rho $$
(2.52)
Requiring parallel transport to commute with the left and right action of $`G`$ means:
$`_h(\theta ^g)=(_h\theta ^g)=\theta ^g`$ (2.53)
$`_h(\theta ^g)=(_h\theta ^g)=\theta ^{ad(h)g}`$ (2.54)
Recalling that $`_h(a\rho )=(_ha)(_h\rho )`$ and $`_h(\rho \rho ^{})=(_h\rho )(_h\rho ^{}),aA,\rho ,\rho ^{}\mathrm{\Gamma }`$ (and similar for $`_h`$), and substituting (2.49) yields respectively:
$$\mathrm{\Gamma }_{g_3,g_2}^{g_1}𝑪$$
(2.55)
and
$$\mathrm{\Gamma }_{ad(h)g_3,ad(h)g_2}^{ad(h)g_1}=\mathrm{\Gamma }_{g_3,g_2}^{g_1}$$
(2.56)
Therefore the same situation arises as in the case of Lie groups, for which parallel transport on the group manifold commutes with left and right action iff the connection components are $`ad(G)`$ \- conserved constant tensors. As for Lie groups, condition (2.56) is satisfied if one takes $`\mathrm{\Gamma }`$ proportional to the structure constants. In our case, we can take any combination of the $`C`$ or $`𝑪`$ structure constants, since both are $`ad(G)`$ conserved constant tensors. As we see below, the $`C`$ constants can be used to define a torsionless connection, while the $`𝑪`$ constants define a parallelizing connection.
As usual, the curvature arises from $`^2`$:
$$^2\theta ^g=R_g^{}^g\theta ^g^{}$$
(2.57)
$$R_{g_2}^{g_1}d\omega _{g_2}^{g_1}+\omega _{g_3}^{g_1}\omega _{g_2}^{g_3}$$
(2.58)
The torsion $`R^g`$ is defined by:
$$R^{g_1}d\theta ^{g_1}+\omega _{g_2}^{g_1}\theta ^{g_2}$$
(2.59)
Using the expression of $`\omega `$ in terms of $`\mathrm{\Gamma }`$ and the Cartan-Maurer equations yields
$`R_{g_2}^{g_1}`$ $`=`$ $`(\mathrm{\Gamma }_{h,g_2}^{g_1}C_{g_3,g_4}^h+\mathrm{\Gamma }_{g_3,h}^{g_1}\mathrm{\Gamma }_{g_4,g_2}^h)\theta ^{g_3}\theta ^{g_4}=`$
$`=`$ $`(\mathrm{\Gamma }_{h,g_2}^{g_1}𝑪_{g_3,g_4}^h+\mathrm{\Gamma }_{g_3,h}^{g_1}\mathrm{\Gamma }_{g_4,g_2}^h\mathrm{\Gamma }_{g_4,h}^{g_1}\mathrm{\Gamma }_{g_4g_3g_4^1,g_2}^h)\theta ^{g_3}\theta ^{g_4}`$
$`R^{g_1}`$ $`=`$ $`(C_{g_2,g_3}^{g_1}+\mathrm{\Gamma }_{g_2,g_3}^{g_1})\theta ^{g_2}\theta ^{g_3}=`$ (2.61)
$`=`$ $`(𝑪_{g_2,g_3}^{g_1}+\mathrm{\Gamma }_{g_2,g_3}^{g_1}\mathrm{\Gamma }_{g_3,g_3g_2g_3^1}^{g_1})\theta ^{g_2}\theta ^{g_3}`$
Thus a connection satisfying:
$$\mathrm{\Gamma }_{g_2,g_3}^{g_1}\mathrm{\Gamma }_{g_3,g_3g_2g_3^1}^{g_1}=𝑪_{g_2,g_3}^{g_1}$$
(2.62)
corresponds to a vanishing torsion $`R^g=0`$ and could be referred to as a “Riemannian” connection.
On the other hand, the choice:
$$\mathrm{\Gamma }_{g_2,g_3}^{g_1}=𝑪_{g_3,g_2^1}^{g_1}$$
(2.63)
corresponds to a vanishing curvature $`R_g^{}^g=0`$, as can be checked by using the fusion equations (2.47) and property (2.45). Then (2.63) can be called the parallelizing connection: finite groups are parallelizable.
Tensor transformations
Under the familiar transformation of the connection 1-form:
$$(\omega _j^i)^{}=a_k^i\omega _l^k(a^1)_j^l+a_k^id(a^1)_j^k$$
(2.64)
the curvature 2-form transforms homogeneously:
$$(R_j^i)^{}=a_k^iR_l^k(a^1)_j^l$$
(2.65)
The transformation rule (2.64) can be seen as induced by the change of basis $`\theta ^i=a_j^i\theta ^j`$, with $`a_j^i`$ invertible $`x`$-dependent matrix (use eq. (2.51) with $`a\rho =a_j^i\theta ^j`$).
Metric
The metric tensor $`\gamma `$ can be defined as an element of $`\mathrm{\Gamma }\mathrm{\Gamma }`$:
$$\gamma =\gamma _{i,j}\theta ^i\theta ^j$$
(2.66)
Requiring it to be invariant under left and right action of $`G`$ means:
$$_h(\gamma )=\gamma =_h(\gamma )$$
(2.67)
or equivalently, by recalling $`_h(\theta ^i\theta ^j)=\theta ^i\theta ^j`$, $`_h(\theta ^i\theta ^j)=\theta ^{ad(h)i}\theta ^{ad(h)j}`$ :
$$\gamma _{i,j}𝑪,\gamma _{ad(h)i,ad(h)j}=\gamma _{i,j}$$
(2.68)
These properties are analogous to the ones satisfied by the Killing metric of Lie groups, which is indeed constant and invariant under the adjoint action of the Lie group.
On finite $`G`$ there are various choices of biinvariant metrics. One can simply take $`\gamma _{i,j}=\delta _{i,j}`$, or $`\gamma _{i,j}=𝑪_{l,i}^k𝑪_{k,j}^l`$.
For any biinvariant metric $`\gamma _{ij}`$ there are tensor transformations $`a_j^i`$ under which $`\gamma _{ij}`$ is invariant, i.e.:
$$a_h^{}^h\gamma _{h,k}a_k^{}^k=\gamma _{h^{},k^{}}a_h^{}^h\gamma _{h,k}=\gamma _{h^{},k^{}}(a^1)_k^k^{}$$
(2.69)
These transformations are simply given by the matrices that rotate the indices according to the adjoint action of $`G`$:
$$a_h^{}^h(g)=\delta _h^{}^{ad(\alpha (g))h}$$
(2.70)
where $`\alpha (g):GG`$ is an arbitrary mapping. Then these matrices are functions of $`G`$ via this mapping, and their action leaves $`\gamma `$ invariant because of the its biinvariance (2.68). Indeed substituting these matrices in (2.69) yields:
$$a_h^{}^h(g)\gamma _{h,k}a_k^{}^k(g)=\gamma _{ad([\alpha (g)]^1)h^{},ad([\alpha (g)]^1)k^{}}=\gamma _{h^{},k^{}}$$
(2.71)
proving the invariance of $`\gamma `$.
Consider now a contravariant vector $`\phi ^i`$ transforming as $`(\phi ^i)^{}=a_j^i(\phi ^j)`$. Then using (2.69) one can easily see that
$$(\phi ^k\gamma _{k,i})^{}=\phi ^k^{}\gamma _{k^{},i^{}}(a^1)_i^i^{}$$
(2.72)
i.e. the vector $`\phi _i\phi ^k\gamma _{k,i}`$ indeed transforms as a covariant vector.
Lie derivative and diffeomorphisms
The notion of diffeomorphisms, or general coordinate transformations, is fundamental in gravity theories. Is there such a notion in the setting of differential calculi on Hopf algebras ? The answer is affirmative, and has been discussed in detail in ref.s . As for differentiable manifolds, it relies on the existence of the Lie derivative.
Let us review the situation for Lie group manifolds. The Lie derivative $`l_{t_i}`$ along a left-invariant tangent vector $`t_i`$ is related to the infinitesimal right translations generated by $`t_i`$:
$$l_{t_i}\rho =\underset{\epsilon 0}{lim}\frac{1}{\epsilon }[_{\mathrm{exp}[\epsilon t_i]}\rho \rho ]$$
(2.73)
$`\rho `$ being an arbitrary tensor field. Introducing the coordinate dependence
$$l_{t_i}\rho (y)=\underset{\epsilon 0}{lim}\frac{1}{\epsilon }[\rho (y+\epsilon t_i)\rho (y)]$$
(2.74)
identifies the Lie derivative $`l_{t_i}`$ as a directional derivative along $`t_i`$. Note the difference in meaning of the symbol $`t_i`$ in the r.h.s. of these two equations: a group generator in the first, and the corresponding tangent vector in the second.
For finite groups the Lie derivative takes the form:
$$l_{t_g}\rho =[_{g^1}\rho \rho ]$$
(2.75)
so that the Lie derivative is simply given by
$$l_{t_g}=_{g^1}id=t_g$$
(2.76)
cf. the definition of $`t_g`$ in (2.38). For example
$$l_{t_g}(\theta ^{g_1}\theta ^{g_2})=\theta ^{ad(g^1)g_1}\theta ^{ad(g^1)g_2}\theta ^{g_1}\theta ^{g_2}$$
(2.77)
As in the case of differentiable manifolds, the Cartan formula for the Lie derivative acting on p-forms holds:
$$l_{t_g}=i_{t_g}d+di_{t_g}$$
(2.78)
see ref.s .
Exploiting this formula, diffeomorphisms (Lie derivatives) along generic tangent vectors $`V`$ can also be consistently defined via the operator:
$$l_V=i_Vd+di_V$$
(2.79)
This requires a suitable definition of the contraction operator $`i_V`$ along generic tangent vectors $`V`$, discussed in ref. .
We have then a way of defining “diffeomorphisms” along arbitrary (and x-dependent) tangent vectors for any tensor $`\rho `$:
$$\delta \rho =l_V\rho $$
(2.80)
and of testing the invariance of candidate lagrangians under the generalized Lie derivative.
Haar measure and integration
Since we want to be able to define actions (integrals on $`p`$-forms) we must now define integration of $`p`$-forms on finite groups.
Let us start with integration of functions $`f`$. We define the integral map $`h`$ as a linear functional $`h:Fun(G)𝑪`$ satisfying the left and right invariance conditions:
$$h(_gf)=0=h(_gf)$$
(2.81)
Then this map is uniquely determined (up to a normalization constant), and is simply given by the “sum over $`G`$” rule:
$$h(f)=\underset{gG}{}f(g)$$
(2.82)
Next we turn to define the integral of a p-form. Within the differential calculus we have a basis of left-invariant 1-forms, which may allow the definition of a biinvariant volume element. In general for a differential calculus with $`n`$ independent tangent vectors, there is an integer $`pn`$ such that the linear space of $`p`$-forms is 1-dimensional, and $`(p+1)`$\- forms vanish identically. We will see explicit examples in the next Section. This means that every product of $`p`$ basis one-forms $`\theta ^{g_1}\theta ^{g_2}\mathrm{}\theta ^{g_p}`$ is proportional to one of these products, that can be chosen to define the volume form $`vol`$:
$$\theta ^{g_1}\theta ^{g_2}\mathrm{}\theta ^{g_p}=ϵ^{g_1,g_2,\mathrm{}g_p}vol$$
(2.83)
where $`ϵ^{g_1,g_2,\mathrm{}g_p}`$ is the proportionality constant. Note that the volume $`p`$-form is obviously left invariant. We can prove that it is also right invariant with the following argument. Suppose that $`vol`$ be given by $`\theta ^{h_1}\theta ^{h_2}\mathrm{}\theta ^{h_p}`$ where $`h_1,h_2,\mathrm{}h_p`$ are given group element labels. Then the right action on $`vol`$ yields:
$$_g[\theta ^{h_1}\mathrm{}\theta ^{h_p}]=\theta ^{ad(g)h_1}\mathrm{}\theta ^{ad(g)h_p}=ϵ^{ad(g)h_1,\mathrm{}ad(g)h_p}vol$$
(2.84)
Recall now that the “epsilon tensor” $`ϵ`$ is necessarily made out of the $`W`$ tensors of eq. (2.28), defining the wedge product. These tensors are invariant under the adjoint action $`ad(g)`$, and so is the $`ϵ`$ tensor. Therefore $`ϵ^{ad(g)h_1,\mathrm{}ad(g)h_p}=ϵ^{h_1,\mathrm{}h_p}=1`$ and $`_gvol=vol`$. This will be verified in the examples of next Section.
Having identified the volume $`p`$-form it is natural to set
$$fvolh(f)=\underset{gG}{}f(g)$$
(2.85)
and define the integral on a $`p`$-form $`\rho `$ as:
$$\rho =\rho _{g_1,\mathrm{}g_p}\theta ^{g_1}\mathrm{}\theta ^{g_p}=\rho _{g_1,\mathrm{}g_p}ϵ^{g_1,\mathrm{}g_p}vol\underset{gG}{}\rho _{g_1,\mathrm{}g_p}(g)ϵ^{g_1,\mathrm{}g_p}$$
(2.86)
Due to the biinvariance of the volume form, the integral map $`:\mathrm{\Gamma }^p𝑪`$ satisfies the biinvariance conditions:
$$_gf=f=_gf$$
(2.87)
Moreover, under the assumption that the volume form belongs to a nontrivial cohomology class, that is $`d(vol)=0`$ but $`vold\rho `$, the important property holds:
$$𝑑f=0$$
(2.88)
with $`f`$ any $`(p1)`$-form: $`f=f_{g_2,\mathrm{}g_p}\theta ^{g_2}\mathrm{}\theta ^{g_p}`$. This property, which allows integration by parts, has a simple proof. Rewrite $`𝑑f`$ as:
$$𝑑f=(df_{g_2,\mathrm{}g_p})\theta ^{g_2}\mathrm{}\theta ^{g_p}+f_{g_2,\mathrm{}g_p}d(\theta ^{g_2}\mathrm{}\theta ^{g_p})$$
(2.89)
Under the cohomology assumption the second term in the r.h.s. vanishes, since $`d(\theta ^{g_2}\mathrm{}\theta ^{g_p})=0`$ (otherwise, being a $`p`$-form, it should be proportional to $`vol`$, and this would contradict the assumption $`vold\rho `$). Using now (2.35) and (2.85):
$`{\displaystyle 𝑑f}={\displaystyle (t_{g_1}f_{g_2,\mathrm{}g_p})\theta ^{g_1}}\theta ^{g_2}\mathrm{}\theta ^{g_p}={\displaystyle [_{g_1^1}f_{g_2,\mathrm{}g_p}f_{g_2,\mathrm{}g_p}]ϵ^{g_1,\mathrm{}g_p}vol}=`$
$`=ϵ^{g_1,\mathrm{}g_p}{\displaystyle \underset{gG}{}}[_{g_1^1}f_{g_2,\mathrm{}g_p}(g)f_{g_2,\mathrm{}g_p}(g)]=0`$ (2.90)
Q.E.D.
## 3 Bicovariant calculus on $`S_3`$
In this Section we illustrate the general theory on the particular example of the permutation group $`S_3`$.
Elements: $`a=(12)`$, $`b=(23)`$, $`c=(13)`$, $`ab=(132)`$, $`ba=(123)`$, $`e`$.
Nontrivial conjugation classes: $`I=[a,b,c]`$, $`II=[ab,ba]`$.
There are 3 bicovariant calculi $`BC_I`$, $`BC_{II}`$, $`BC_{I+II}`$ corresponding to the possible unions of the conjugation classes . They have respectively dimension 3, 2 and 5. We examine here the $`BC_I`$ and $`BC_{II}`$ calculi.
$`BC_I`$ differential calculus
Basis of the 3-dimensional vector space of one-forms:
$$\theta ^a,\theta ^b,\theta ^c$$
(3.1)
Basis of the 4-dimensional vector space of two-forms:
$$\theta ^a\theta ^b,\theta ^b\theta ^c,\theta ^a\theta ^c,\theta ^c\theta ^b$$
(3.2)
Every wedge product of two $`\theta `$ can be expressed as linear combination of the basis elements:
$$\theta ^b\theta ^a=\theta ^a\theta ^c\theta ^c\theta ^b,\theta ^c\theta ^a=\theta ^a\theta ^b\theta ^b\theta ^c$$
(3.3)
Basis of the 3-dimensional vector space of three-forms:
$$\theta ^a\theta ^b\theta ^c,\theta ^a\theta ^c\theta ^b,\theta ^b\theta ^a\theta ^c$$
(3.4)
and we have:
$`\theta ^c\theta ^b\theta ^a=\theta ^c\theta ^a\theta ^c=\theta ^a\theta ^c\theta ^a=\theta ^a\theta ^b\theta ^c`$
$`\theta ^b\theta ^c\theta ^a=\theta ^b\theta ^a\theta ^b=\theta ^a\theta ^b\theta ^a=\theta ^a\theta ^c\theta ^b`$
$`\theta ^c\theta ^a\theta ^b=\theta ^c\theta ^b\theta ^c=\theta ^b\theta ^c\theta ^b=\theta ^b\theta ^a\theta ^c`$ (3.5)
Basis of the 1-dimensional vector space of four-forms:
$$vol=\theta ^a\theta ^b\theta ^a\theta ^c$$
(3.6)
and we have:
$$\theta ^{g_1}\theta ^{g_2}\theta ^{g_3}\theta ^{g_4}=ϵ^{g_1,g_2,g_3,g_4}vol$$
(3.7)
where the nonvanishing components of the $`ϵ`$ tensor are:
$`ϵ_{abac}=ϵ_{acab}=ϵ_{cbca}=ϵ_{cacb}=ϵ_{babc}=ϵ_{bcba}=1`$ (3.8)
$`ϵ_{baca}=ϵ_{caba}=ϵ_{abcb}=ϵ_{cbab}=ϵ_{acbc}=ϵ_{bcac}=1`$ (3.9)
Cartan-Maurer equations:
$`d\theta ^a+\theta ^b\theta ^c+\theta ^c\theta ^b=0`$
$`d\theta ^b+\theta ^a\theta ^c+\theta ^c\theta ^a=0`$
$`d\theta ^c+\theta ^a\theta ^b+\theta ^b\theta ^a=0`$ (3.10)
The exterior derivative on any three-form of the type $`\theta \theta \theta `$ vanishes, as one can easily check by using the Cartan-Maurer equations and the equalities between exterior products given above. Then, as shown in the previous Section, integration of a total differential vanishes on the “group manifold” of $`S_3`$ corresponding to the $`BC_I`$ bicovariant calculus. This “group manifold” has three independent directions, associated to the cotangent basis $`\theta ^a,\theta ^b,\theta ^c`$. Note however that the volume element is of order four in the left-invariant one-forms $`\theta `$.
$`BC_{II}`$ differential calculus
Basis of the 2-dimensional vector space of one-forms:
$$\theta ^{ab},\theta ^{ba}$$
(3.11)
Basis of the 1-dimensional vector space of two-forms:
$$vol=\theta ^{ab}\theta ^{ba}=\theta ^{ba}\theta ^{ab}$$
(3.12)
so that:
$$\theta ^{g_1}\theta ^{g_2}=ϵ^{g_1,g_2}vol$$
(3.13)
where the $`ϵ`$ tensor is the usual 2-dimensional Levi-Civita tensor.
Cartan-Maurer equations:
$$d\theta ^{ab}=0,d\theta ^{ba}=0$$
(3.14)
Thus the exterior derivative on any one-form $`\theta ^g`$ vanishes and integration of a total differential vanishes on the group manifold of $`S_3`$ corresponding to the $`BC_{II}`$ bicovariant calculus. This group manifold has two independent directions, associated to the cotangent basis $`\theta ^{ab},\theta ^{ba}`$.
Visualization of the $`S_3`$ group “manifold”
We can draw a picture of the group manifold of $`S_3`$. It is made out of 6 points, corresponding to the group elements and identified with the functions $`x^e,x^a,x^b,x^c,x^{ab},x^{ba}`$.
$`BC_I`$ \- calculus:
From each of the six points $`x^g`$ one can move in three directions, associated to the tangent vectors $`t_a,t_b,t_c`$, reaching three other points whose “coordinates” are
$$_ax^g=x^{ga},_bx^g=x^{gb},_cx^g=x^{gc}$$
(3.15)
The 6 points and the “moves” along the 3 directions are illustrated in the Fig. 1. The links are not oriented since the three group elements $`a,b,c`$ coincide with their inverses.
$`BC_{II}`$ \- calculus:
From each of the six points $`x^g`$ one can move in two directions, associated to the tangent vectors $`t_{ab},t_{ba}`$, reaching two other points whose “coordinates” are
$$_{ab}x^g=x^{gba},_{ba}x^g=x^{gab}$$
(3.16)
The 6 points and the “moves” along the 3 directions are illustrated in Fig. 1. The arrow convention on a link labeled (in italic) by a group element $`h`$ is as follows: one moves in the direction of the arrow via the action of $`_h`$ on $`x^g`$. (In this case $`h=ab`$). To move in the opposite direction just take the inverse of $`h`$.
The pictures in Fig. 1 characterize the bicovariant calculi $`BC_I`$ and $`BC_{II}`$ on $`S_3`$, and were drawn in ref. as examples of digraphs, used to characterize different calculi on sets. Here we emphasize their geometrical meaning as finite group “manifolds”.
Fig. 1 : $`S_3`$ group manifold, and moves of the points under the group action
Acknowledgements
It is a pleasure to thank the organizers of the Corfu Summer Institute on Elementary Particle Physics for their invitation to discuss physics in such a beautiful and relaxed atmosphere.
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# Affine Lie algebras and tame quivers
## 0. Introduction
#### 0.0.1.
There is a remarkable connection between the theory of representations of quivers and the structure theory of Lie algebras. The first manifestation of this connection was discovered by P. Gabriel \[Gab72\]. Let $`Q`$ be a quiver obtained by orienting edges of the Dynkin graph corresponding to a simple simply laced Lie algebra $`𝔤`$. Gabriel proved that the set $`𝒯`$ of isomorphism classes of indecomposable complex representations of $`Q`$ is in one-to-one correspondence with the set $`R_+`$ of positive roots of $`𝔤`$.
Gabriel’s result was soon extended by J. Bernstein, I. Gelfand, and V. Ponomarev \[BGP73\], who introduced reflection functors $`𝒮_i`$ corresponding to Coxeter generators of the Weyl group of $`𝔤`$. The reflection functor $`𝒮_i`$ acts from the category $`(Q)`$ of complex representations of $`Q`$ to the category $`(Q^{})`$, where $`Q^{}`$ differs from $`Q`$ only by orientation. Using the reflection functors Bernstein, Gelfand and Ponomarev were able to give another proof of the Gabriel theorem.
#### 0.0.2.
The set $`R_+`$ of positive roots corresponds to a basis of a maximal nilpotent subalgebra $`𝔫`$ of $`𝔤`$, and one might guess that there exists an intrinsic Lie bracket on the $``$-linear span of the set $`𝒯`$, such that the resultant Lie algebra is isomorphic to $`𝔫`$. However this Lie bracket was introduced only 18 years later by C. M. Ringel \[Rin90b, Rin90c\]. Ringel actually considers representations of $`Q`$ over finite fields rather than over complex numbers, so we use a variant of his definition due to A. Schofield \[Sch91\] and G. Lusztig \[Lus91b\].
Let $`[𝐏_\alpha ]`$, $`[𝐏_\beta ]𝒯`$ be isomorphism classes of indecomposable representations of $`Q`$. Then their Lie bracket is defined as follows
(0.0.2.1)
$$[[𝐏_\alpha ],[𝐏_\beta ]]=\underset{[𝐏_\gamma ]𝒯}{}(\chi (N_{𝐏_\alpha ,𝐏_\beta ;𝐏_\gamma })\chi (N_{𝐏_\beta ,𝐏_\alpha ;𝐏_\gamma }))[𝐏_\gamma ],$$
where $`\chi (N_{𝐏_\alpha ,𝐏_\beta ;𝐏_\gamma })`$ is the Euler characteristic with compact support of the algebraically constructible set $`N_{𝐏_\alpha ,𝐏_\beta ;𝐏_\gamma }`$ of all subrepresentations $`𝐕𝐏_\gamma `$ such that $`𝐕`$ is isomorphic to $`𝐏_\alpha `$, and $`𝐏_\gamma /𝐕`$ is isomorphic to $`𝐏_\beta `$.
Thus we obtain a complex Lie algebra denoted by $`𝔫^{}`$ with a distinguished basis $`\{[𝐏]\}_{[𝐏]𝒯}`$. In \[Rin90c\] Ringel proved that $`𝔫^{}`$ and $`𝔫`$ are isomorphic as $`R_+`$-graded complex Lie algebras, and, moreover, he was able to find the structure constants in $`𝔫^{}`$. Namely, let $`𝐏_\alpha `$ be the unique up to an isomorphism indecomposable representation of $`Q`$ corresponding to a root $`\alpha R_+`$. Then given $`\alpha `$, $`\beta R_+`$ one has
(0.0.2.2)
$$[[𝐏_\alpha ],[𝐏_\beta ]]=\{\begin{array}{cc}ϵ(\alpha ,\beta )[𝐏_{\alpha +\beta }]\hfill & \text{if }\alpha +\beta R_+\text{ ,}\hfill \\ 0\hfill & \text{if }\alpha +\beta R_+\text{ },\hfill \end{array}$$
where $`ϵ`$ is a bimultiplicative two-cocycle on the root lattice of $`𝔤`$, uniquely defined by its values on pairs of simple roots: $`ϵ(i,i)=1`$ for any simple root $`i`$, $`ϵ(i,j)=1`$ if there is an arrow $`ij`$ in $`Q`$, $`ϵ(i,j)=1`$ otherwise. The choice of bimultiplicative cocycle $`ϵ`$ corresponds to a choice of orientation of edges which converts a Dynkin graph into quiver.
It is interesting to note that the choice of Chevalley basis in $`𝔫`$ and, more generally, in $`𝔤`$, such that structure constants are given by the cocycle $`ϵ`$, was first introduced in relation to a vertex operator construction of the affine Lie algebra $`\widehat{𝔤}`$ associated to $`𝔤`$ (see \[FK80, Seg81\]).
Let us also remark that there is an equivalent definition of the cocycle $`ϵ`$ in terms of the category $`(Q)`$. Namely,
(0.0.2.3)
$$ϵ(\alpha ,\beta )=(1)^{(dim_{}\mathrm{Hom}_{(Q)}(𝐏_\alpha ,𝐏_\beta )dim_{}\mathrm{Ext}_{(Q)}^1(𝐏_\alpha ,𝐏_\beta ))}.$$
Ringel’s proof of (0.0.2.2) is based on a case-by-case study of all possible varieties $`N_{𝐏_\alpha ,𝐏_\beta ;𝐏_\gamma }`$ and is rather long. In this paper we give a new and short proof of the Ringel theorem using reflection functors of Bernstein, Gelfand, and Ponomarev. Instead of studying varieties $`N_{𝐏_\alpha ,𝐏_\beta ;𝐏_\gamma }`$ and calculating their Euler characteristics for a particular quiver $`Q`$ we consider all quivers with the same underlying Dynkin graph and use functorial properties of the Lie algebra $`𝔫^{}`$.
More explicitly, given a quiver $`Q`$ we define a complex Lie algebra $`𝔫^ϵ`$. As a linear space $`𝔫^ϵ`$ has a basis $`\{\stackrel{~}{e}_\alpha \}_{\alpha R_+}`$ corresponding to isomorphism classes of indecomposable representations of $`Q`$. The Lie bracket is defined as follows
$$[\stackrel{~}{e}_\alpha ,\stackrel{~}{e}_\beta ]=\{\begin{array}{cc}ϵ(\alpha ,\beta )\stackrel{~}{e}_{\alpha +\beta }\hfill & \text{if }\alpha +\beta R_+\text{ ,}\hfill \\ 0\hfill & \text{if }\alpha +\beta R_+\text{ }.\hfill \end{array}$$
Because of (0.0.2.3) $`𝔫^ϵ`$ is functorial with respect to reflection functors. The same is true for $`𝔫^{}`$ by definition.
It turns out that there are enough reflection functors to ensure that a Lie algebra with a basis parameterized by isomorphism classes of indecomposable representations of $`Q`$ and with functorial structure constants is unique. More precisely, we introduce a homomorphism $`\mathrm{\Xi }:𝔫^ϵ𝔫^{}`$ given on generators by $`\mathrm{\Xi }(\stackrel{~}{e}_i)=[𝐏_i]`$ for any simple root $`i`$. Then we use reflection functors to prove that $`\mathrm{\Xi }(\stackrel{~}{e}_\alpha )=[𝐏_\alpha ]`$ for any $`\alpha R_+`$, which is equivalent to the Ringel theorem.
#### 0.0.3.
Shortly after appearance of Gabriel’s paper L. A. Nazarova \[Naz73\] and independently P. Donovan and M. R. Freislich \[DF73\] classified indecomposable representations of quivers associated to simply laced extended Dynkin graphs, which we call quivers of affine type (they are also called tame quivers).
For quivers of affine type the bijection between isomorphism classes of indecomposable representations and positive roots of the corresponding affine Lie algebra $`𝔤`$ holds only for positive real roots $`\alpha R_+^{\mathrm{re}}`$. In the case of a positive imaginary root $`\alpha R_+^{\mathrm{im}}`$ there exists an uncountable family $`𝒯_\alpha `$ of non-isomorphic indecomposable representations corresponding to $`\alpha `$. According to Dlab and Ringel \[DR76\] the set $`𝒯_\alpha `$ does not depend on $`\alpha R_+^{\mathrm{im}}`$ and admits a surjection
$$\mu :𝒯_\alpha ^1$$
such that $`\mu `$ is injective for all points of $`^1`$ except for $`L3`$ exceptional points for which fibers are finite sets.
Since $`𝒯_\alpha `$ is an infinite set for $`\alpha R_+^{\mathrm{im}}`$ one has to adjust the definition of $`𝔫^{}`$. One way to do it, proposed by Ringel \[Rin90a, Rin93\], consists of replacing representations of $`Q`$ by formal composition series. We adopt another approach due to A. Schofield \[Sch91\] and G. Lusztig \[Lus91a, Lus91b\]. Instead of formal linear combinations of elements of the set $`𝒯`$ of indecomposable representations we consider complex valued constructible (in some sense) functions on $`𝒯`$. One can generalize the definition of Lie bracket (0.0.2.1). Finally we define $`𝔫^{}`$ to be the Lie subalgebra of the Lie algebra of constructible functions on $`𝒯`$ generated by characteristic functions of simple representations.
The main result of the present paper is a generalization of the Ringel theorem described in 0.0.2 to affine case. We show, in particular, that the Lie algebra $`𝔫^{}`$ is isomorphic to the positive part $`𝔫`$ of the affine Lie algebra $`𝔤`$ corresponding to the quiver $`Q`$. The Lie algebra $`𝔫`$ is graded by the set $`R_+`$ of positive roots and is the precise analogue of a maximal nilpotent subalgebra of a finite dimensional simple Lie algebra.
Our results imply that the Lie algebra $`𝔫^{}`$ contains the characteristic function of indecomposable representation $`𝐏_\alpha `$ corresponding to any positive real root $`\alpha `$, and we can find the structure constants for the Lie bracket of two such characteristic functions. Namely, if $`\alpha `$, $`\beta `$, $`\alpha +\beta R_+^{\mathrm{re}}`$ then
$$[[𝐏_\alpha ],[𝐏_\beta ]]=ϵ^{}(\alpha ,\beta )[𝐏_{\alpha +\beta }].$$
Here abusing notation we denote by $`[𝐏_\alpha ]`$ the characteristic function of $`[𝐏_\alpha ]𝒯`$. The structure constants $`ϵ^{}`$ are related to the bimultiplicative two-cocycle $`ϵ`$ defined in 0.0.2 via a twist by coboundary
$$ϵ^{}(\alpha ,\beta )=ϵ(\alpha ,\beta )\xi (\alpha +\beta )\xi ^1(\alpha )\xi ^1(\beta ),$$
where $`\xi (\alpha )=(1)^{1+dim_{}\mathrm{Hom}_{(Q)}(𝐏_\alpha ,𝐏_\alpha )}`$. Note that in the case of a quiver with underlying Dynkin graph $`\mathrm{Hom}_{(Q)}(𝐏_\alpha ,𝐏_\alpha )=`$ for any indecomposable $`𝐏_\alpha `$ and, therefore, the twist does not affect Ringel’s structure constants.
The case of an imaginary root $`\alpha R_+^{\mathrm{im}}`$ is more involved. We give a description of the root space $`𝔫_\alpha ^{}`$ in terms of the surjection $`\mu :𝒯_\alpha ^1`$. Namely, $`𝔫_\alpha ^{}`$ consists of all functions $`f`$ on $`𝒯_\alpha `$ such that
(0.0.3.1)
$$\underset{[𝐏]\mu ^1(z)}{}f([𝐏])\text{ does not depend on }z^1\text{ }.$$
The following remarkable identity guarantees that $`𝔫_\alpha ^{}`$ has the correct dimension
$$\underset{z^1}{}(N_z1)=\mathrm{rank}𝔤2,$$
where $`N_z`$ is the order of the fiber $`\mu ^1(z)`$.
Considering constructible functions on $`𝒯`$ with values in $``$ satisfying condition (0.0.3.1) on $`𝒯_\alpha `$ for any $`\alpha R_+^{\mathrm{im}}`$ one obtains a lattice $`\mathrm{}^{}𝔫^{}`$. Our results imply that $`\mathrm{}^{}`$ is closed with respect to the Lie bracket in $`𝔫^{}`$ (i.e. it is an order). We also prove that $`\mathrm{}^{}`$ is characterized by the property that it is the minimal order in $`𝔫^{}`$ containing simple root generators. In particular, the lattice $`\mathrm{}_\alpha ^{}`$ for any imaginary root $`\alpha `$ is isomorphic to the root lattice of a simple Lie algebra $`𝔤_0`$, such that $`𝔤=\widehat{𝔤}_0`$.
Our strategy in affine case is similar to the one we use in the proof of the Ringel theorem. We define some ad hoc Lie algebra $`𝔫^ϵ`$ using the cocycle $`ϵ`$, and a homomorphism $`\mathrm{\Xi }:𝔫^ϵ𝔫^{}`$. Then we use functorial properties of $`𝔫^ϵ`$ and $`𝔫^{}`$ to describe $`\mathrm{\Xi }`$. In the affine case, however, the set of reflection functors is not enough to fix $`\mathrm{\Xi }`$, and we employ additional functors $`𝒞_z`$ introduced by Dlab and Ringel. A functor $`𝒞_z`$ corresponds to a point $`z^1`$ and is a full, faithful, exact functor from $`(C_{N_z})`$ to $`(Q)`$, where $`C_{N_z}`$ is the cyclic quiver with $`N_z`$ vertices. Both $`𝔫^{}`$ and $`𝔫^ϵ`$ are functorial with respect to functors $`𝒞_z`$ by construction. Using the reflection functors and the functors $`𝒞_z`$ we are able to describe the map $`\mathrm{\Xi }`$ for any real root space and a codimension one subspace in any imaginary root space. The proof is completed with analysis of the positive part of an affine Lie algebra $`\widehat{sl}_2`$ embedded into $`𝔤`$.
#### 0.0.4.
Our description of the isomorphism of the Lie algebra $`𝔫^{}`$ constructed via representation theory of a quiver of affine type with the positive part $`𝔫`$ of the corresponding affine Lie algebra reveals a certain fine structure of affine Lie algebras and their root systems. It turns out that for any orientation of edges of the extended Dynkin graph one can canonically associate $`L3`$ affine Lie algebras of type $`A_n^{(1)}`$ embedded into $`𝔤`$. In $`D_n^{(1)}`$ and $`E_n^{(1)}`$ cases these subalgebras correspond precisely to the connected components of the non-extended Dynkin graph with the branching vertex removed. In particular, it follows that $`L=3`$ in $`DE`$ case.
If orientation of $`Q`$ is such that each vertex is either a sink or a source, then one can interpret representations of $`Q`$ in terms of a finite subgroup $`\mathrm{\Gamma }`$ of $`SL(2,)`$ associated to $`Q`$ via McKay correspondence \[Lus92\]. In this case the type $`A_n^{(1)}`$ subalgebras arise from maximal cyclic subgroups of $`\mathrm{\Gamma }`$ and their positive roots correspond to representations of $`\mathrm{\Gamma }`$ induced from the cyclic subgroups. Thus our construction of the positive part $`𝔫`$ of an affine Lie algebra $`𝔤`$ can be viewed as a far-reaching expansion of the McKay correspondence between group $`\mathrm{\Gamma }SL(2,)`$ and extended Dynkin graph of $`𝔤`$.
#### 0.0.5.
We conclude Introduction with remarks concerning some generalizations and interpretations of our results.
First let us note that using species instead of quivers one can extend results of the paper to all simple and affine types of Lie algebras. Though we consider only simply laced case in order to preserve clarity of the exposition, the statements and the proofs can be generalized almost word-by-word to species.
Second we remark that Ringel’s construction of the nilpotent Lie algebra $`𝔫`$ was extended by L. Peng and J. Xiao \[PX97\] to the whole simple Lie algebra $`𝔤`$ via the root category $`(Q)=D^b(Q)/T^2`$. Our proof of the Ringel theorem can be extended to this setting. Moreover it would become more transparent as the reflection functors act more naturally in the root category than in abelian category $`(Q)`$. It would be very interesting to extend the construction of $`𝔫^{}`$ to the root category of a quiver of affine type.
Let us finally mention that results and techniques used in this paper could be stated in the language of algebraic stacks. For example, description (0.0.3.1) of imaginary root subspaces of $`𝔫^{}`$ indicates that the right setting for $`𝔫^{}`$ is that of cohomology of the set $`𝒯`$ considered as a stack. The details will be provided elsewhere.
#### Acknowledgements
The authors are grateful to Mikhail Khovanov for his careful reading of the manuscript and many valuable comments. This research was supported in part by NSF grant DMS-9700765.
## 1. Cartan datum, Lie algebras, and quivers
Throughout the paper $``$ and $``$ denote, respectively, integer numbers and complex numbers; $`_+=\{k,k0\}`$.
### 1.1. Cartan datum and Dynkin graph
#### 1.1.1.
A *Cartan datum* is a pair $`(I,<,>)`$, consisting of a finite set $`I`$ and a bilinear form $`<,>`$ on the free abelian group $`[I]`$, with values in $``$. The bilinear form should satisfy the following conditions:
$`<i,i>`$ $`=2\text{ for all }iI,`$
$`<i,j>`$ $`0\text{ for all }ij,`$
$`<i,j>`$ $`=0\text{ if }<j,i>=0.`$
The matrix $`a_{ij}=<i,j>`$ is called the Cartan matrix of the Cartan datum $`(I,<,>)`$.
#### 1.1.2.
A Cartan datum is said to be *irreducible* if the corresponding Cartan matrix cannot be made block-diagonal by simultaneous permutations of rows and columns.
A Cartan datum is said to be of *finite type* if the corresponding Cartan matrix is positive definite.
A Cartan datum is said to be of *affine type* if the corresponding Cartan matrix is irreducible, positive semi-definite, but not positive definite.
A Cartan datum is said to be *symmetric* if $`a_{ij}=a_{ji}`$ for all $`i`$, $`j`$.
A Cartan datum is said to be *simply laced* if $`a_{ij}\{0,1\}`$ for all $`ij`$.
If a Cartan datum is simply laced then it is symmetric.
#### 1.1.3.
From now on all the Cartan data is assumed to be symmetric. See, however, 0.0.5.
#### 1.1.4.
For a given finite set $`S`$ we denote by $`𝒫_2(S)`$ the set of all two-element subsets of $`S`$.
By definition, a finite graph is a triple $`(I,E,\mathrm{Ends})`$, consisting of two finite sets $`I`$ (vertices) and $`E`$ (edges), and a map $`\mathrm{Ends}:E𝒫_2(I)`$.
To a symmetric Cartan datum $`(I,<,>)`$ we associate a *Dynkin graph* $`(I,E,\mathrm{Ends})`$ as follows: the set of vertices coincides with $`I`$, two vertices $`i`$ and $`j`$ are joined by $`a_{ij}`$ edges, there are no edges joining a vertex with itself.
A Dynkin graph is called irreducible (resp. of finite type, of affine type) if the corresponding Cartan datum is irreducible (resp. of finite type, of affine type).
#### 1.1.5.
The lists of all irreducible Dynkin graphs of finite and affine types are contained in Figures 1 and 2 respectively (cf. \[Bou68\]).
The integers on vertices of affine Dynkin graphs on Figure 2 are components of the first imaginary root $`\delta `$ (see 5.1.1).
Note that all symmetric Cartan data of finite or affine type are simply laced except for $`A_1^{(1)}`$.
#### 1.1.6.
If one removes a vertex together with adjacent edges from a Dynkin graph of affine type one obtains a Dynkin graph of finite type (not necessarily irreducible).
### 1.2. A subalgebra of a Kac-Moody algebra
#### 1.2.1.
Let $`𝔫`$ be the quotient of the free Lie algebra over $``$ with generators $`\{e_i\}_{iI}`$ by the ideal generated by the following relations
(1.2.1.1)
$$(\mathrm{ad}(e_i))^{1a_{ij}}e_j=0\text{ if }ij\text{ },$$
where $`\mathrm{ad}(x)y=[x,y]`$.
Relations 1.2.1.1 are called the Serre relations.
The Lie algebra $`𝔫`$ is $`_+[I]`$-graded with $`\mathrm{deg}e_i=i`$.
### 1.3. Quiver and the Euler cocycle
#### 1.3.1.
A *quiver* $`Q`$ is an oriented graph, that is a quadruple $`(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$, consisting of two finite sets $`I`$ (vertices) and $`\mathrm{\Omega }`$ (oriented edges), and two maps $`\mathrm{In}`$ and $`\mathrm{Out}`$ from $`\mathrm{\Omega }`$ to $`I`$. The underlying non-oriented graph is given by $`(I,\mathrm{\Omega },\{\mathrm{In},\mathrm{Out}\})`$.
A quiver is said to be of finite (resp. affine) type if the underlying non-oriented graph is a Dynkin graph of finite (resp. affine) type.
#### 1.3.2.
For a quiver $`Q`$ of finite or affine type we denote by $`𝔫(Q)`$ the Lie algebra over $``$ associated with the underlying Dynkin graph of $`Q`$ as in 1.2.1.
#### 1.3.3.
The abelian group $`[I]`$ is called the *root lattice*.
#### 1.3.4.
Given a quiver $`Q`$ we introduce a biadditive form $`e:[I]\times [I]`$ given by the following values on generators:
$$e(i,j)=\delta _{ij}\text{(number of }h\mathrm{\Omega }\text{ such that }i=\mathrm{Out}(h)\text{ and }j=\mathrm{In}(h)\text{),}$$
where $`\delta _{ij}`$ is the Kronecker symbol. The form $`e`$ is called Euler form (see 1.4.6 for a justification of this name).
If $`Q`$ is of finite or affine type then
(1.3.4.1)
$$e(\alpha ,\beta )+e(\beta ,\alpha )=<\alpha ,\beta >.$$
Having the Euler form $`e`$ we define *Euler cocycle* $`ϵ`$ to be a bimultiplicative function from $`[I]\times [I]`$ to $`\{\pm 1\}`$ given by
$$ϵ(\alpha ,\beta )=(1)^{e(\alpha ,\beta )}.$$
If we want to specify the quiver $`Q`$ used in the definition of $`ϵ`$ we write $`ϵ_Q`$.
Euler cocycle $`ϵ`$ is the main building block of the structure constants of $`𝔫`$ in the natural basis associated to the quiver $`Q`$. The precise construction of the natural basis will be carried out in Sections 4 and 5 in the finite and the affine cases respectively.
### 1.4. Category of representations of a quiver
#### 1.4.1.
Let $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$ be a quiver.
A (finite-dimensional) *representation* of $`Q`$ (over $``$) is the following data:
* a finite dimensional $`I`$-graded vector space $`V=_{iI}V_i`$ over $``$,
* a collection of linear maps $`x=_{h\mathrm{\Omega }}x_h_{h\mathrm{\Omega }}\mathrm{Hom}_{}(V_{\mathrm{Out}(h)},V_{\mathrm{In}(h)})`$.
A morphism from a representation $`(V,x)`$ to another representation $`(V^{},x^{})`$ is an $`I`$-graded $``$-linear map $`\varphi :VV^{}`$, such that $`x_h^{}\varphi _{\mathrm{Out}(h)}=\varphi _{\mathrm{In}(h)}x_h`$ for each $`h\mathrm{\Omega }`$. The composition of morphisms is the composition of linear maps.
We denote the category of representations of the quiver by $`\overline{}(Q)`$. It is an abelian category with respect to the natural additive structure on morphisms.
#### 1.4.2.
Let $`(Q)`$ be a full subcategory of $`\overline{}(Q)`$ with objects being $`𝐌=(V,x)\mathrm{Ob}(\overline{}(Q))`$ for which there exists $`N(𝐌)`$ such that $`x_{h_N}\mathrm{}x_{h_1}=0`$ for any sequence $`h_1,\mathrm{},h_N\mathrm{\Omega }`$ with $`\mathrm{In}(h_i)=\mathrm{Out}(h_{i+1})`$. We call objects of $`(Q)`$ *nilpotent* representations of $`Q`$.
The subcategory $`(Q)`$ is closed with respect to extensions. If $`Q`$ does not contain an oriented cycle then $`(Q)=\overline{}(Q)`$.
#### 1.4.3.
*Dimension* of a representation $`𝐌=(V,x)`$ of $`Q`$ is an element of $`[I]`$ given by the graded dimension of $`V`$:
$$dim_{(Q)}𝐌=\underset{iI}{}dim_{}(V_i)i.$$
#### 1.4.4.
It is easy to see that isomorphism classes of simple objects in $`(Q)`$ are in one-to-one correspondence with vertices of $`Q`$ (we recall that $`\mathrm{Ob}((Q))`$ consists of nilpotent representations). The simple object corresponding to a vertex $`iI`$ is given (up to an isomorphism) by
$$𝐏_i=(V,x),\text{ where }\{\begin{array}{cc}V_j=\{0\},ji,\hfill & \\ V_i=,\hfill & \\ x=0.\hfill & \end{array}$$
Note that $`dim_{(Q)}𝐏_i=iI[I]`$.
#### 1.4.5.
Let $`𝐌\mathrm{Ob}((Q))`$. It follows from 1.4.4 that $`(dim_{(Q)}𝐌)_i`$ is equal to the number of factors isomorphic to $`𝐏_i`$ in Jordan-Hölder series of $`𝐌`$.
#### 1.4.6.
The following proposition explains the homological meaning of the Euler form $`e`$.
###### Proposition.
1. The category $`(Q)`$ is hereditary, i.e. $`\mathrm{Ext}_{(Q)}^n(𝐌,𝐍)=0`$ for any $`𝐌,𝐍\mathrm{Ob}((Q))`$ and any $`n2`$.
2. The following equality holds for any $`𝐌,𝐍\mathrm{Ob}((Q))`$.
(1.4.6.1)
$$\begin{array}{c}e(dim_{(Q)}𝐌,dim_{(Q)}𝐍)=\hfill \\ \hfill =dim_{}\mathrm{Hom}_{(Q)}(𝐌,𝐍)dim_{}\mathrm{Ext}_{(Q)}^1(𝐌,𝐍)=\\ \hfill =\underset{k=0}{\overset{\mathrm{}}{}}(1)^kdim_{}\mathrm{Ext}_{(Q)}^k(𝐌,𝐍)\end{array}$$
###### Proof.
The proof is standard. Let us consider the following exact sequence of functors from $`((Q))^{}\times (Q)`$ to the category of $``$-linear spaces.
(1.4.6.3)
$$\begin{array}{c}0\stackrel{}{}\mathrm{Hom}_{(Q)}(𝐌,𝐍)\stackrel{𝜋}{}_{iI}\mathrm{Hom}_{}(V_i,W_i)\stackrel{𝜌}{}\hfill \\ \hfill \stackrel{𝜌}{}_{h\mathrm{\Omega }}\mathrm{Hom}_{}(V_{\mathrm{Out}(h)},W_{\mathrm{In}(h)})\stackrel{𝜎}{}\mathrm{Ext}_{(Q)}^1(𝐌,𝐍)\stackrel{}{}0.\end{array}$$
Here $`𝐌=(V,x)`$, $`𝐍=(W,y)\mathrm{Ob}((Q))`$, and $`\pi `$, $`\rho `$, $`\sigma `$ are natural transformations given as follows
$$\begin{array}{cc}\hfill \pi (\varphi )& =\varphi _i,\hfill \\ \hfill \rho (\varphi _i)& =(\varphi _{\mathrm{In}(h)}x_hy_h\varphi _{\mathrm{Out}(h)}),\hfill \\ \hfill \sigma (\psi _h)& =[0𝐍𝐄𝐌0],\hfill \end{array}$$
where $`𝐄=(_{iI}(V_iW_i),_{h\mathrm{\Omega }}\left(\begin{array}{cc}x_h& 0\\ \psi _h& y_h\end{array}\right))`$. The two middle terms in the exact sequence 1.4.6.3 are exact functors (both with respect to the first and to the second arguments). Thus $`\mathrm{Ext}_{(Q)}^1`$ is right exact, which proves 1.4.6.1. The equality 1.4.6.2 follows from evaluation of dimensions in 1.4.6.3. ∎
#### 1.4.7.
Let $`𝒳(Q)`$ be the set of isomorphism classes of objects of $`(Q)`$. We denote the isomorphism class of $`𝐌\mathrm{Ob}((Q))`$ by $`[𝐌]`$.
We also use the following notation:
$`𝒳_\alpha (Q)=`$ $`\{[𝐌]𝒳(Q)|dim_{(Q)}𝐌=\alpha \},`$
$`𝒯(Q)=`$ $`\{[𝐌]𝒳(Q)|\text{ }𝐌\text{ is indecomposable }\},`$
$`𝒯_\alpha (Q)=`$ $`\{[𝐌]𝒯(Q)|dim_{(Q)}𝐌=\alpha \}.`$
So $`𝒯(Q)`$ is the set of isomorphism classes of indecomposable objects. For example, $`𝒳_i(Q)=𝒯_i(Q)=\{[𝐏_i]\}`$.
#### 1.4.8.
The set $`𝒳_\alpha (Q)`$ has a natural structure of the orbit space for an action of an algebraic group on an affine algebraic variety over $``$.
Let $`\alpha =_{iI}\alpha _ii_+[I]`$. Consider an affine space over $``$
$$𝐄_\alpha =\underset{h\mathrm{\Omega }}{}\mathrm{Hom}_{}(^{\alpha _{\mathrm{Out}(h)}},^{\alpha _{\mathrm{In}(h)}}).$$
Let $`𝐆_\alpha =_{iI}GL(\alpha _i,)`$. The group $`𝐆_\alpha `$ acts on $`𝐄_\alpha `$ by $`x_h^g=g_{\mathrm{In}(h)}x_hg_{\mathrm{Out}(h)}^1`$, $`g𝐆_\alpha `$.
When we want to specify the quiver $`Q`$ used in the definition of $`𝐄_\alpha `$ and $`𝐆_\alpha `$ we write $`𝐄_\alpha (Q)`$ and $`𝐆_\alpha (Q)`$.
There is a natural bijection between the set of isomorphism classes of objects in $`\overline{}(Q)`$ of dimension $`\alpha `$ and the set of orbits of $`𝐆_\alpha `$ in $`𝐄_\alpha `$. Namely, given a point $`x𝐄_\alpha `$, $`(^\alpha ,x)`$ is a representative of the isomorphism class of objects in $`\overline{}(Q)`$, corresponding to this point. This class depends only on the orbit to which the point belongs. And vice versa, in each isomorphism class there are objects of the form $`(^\alpha ,x)`$.
#### 1.4.9.
Let $`c=(h_1,\mathrm{},h_n)`$ be an oriented cycle, that is $`n`$-tuple of elements of $`\mathrm{\Omega }`$, such that $`\mathrm{In}(h_i)=\mathrm{Out}(h_{i+1})`$ for $`1in1`$ and $`\mathrm{In}(h_n)=\mathrm{Out}(h_1)`$. Given $`𝐌=(V,x)\mathrm{Ob}(\overline{}(Q))`$, we call $`x_{h_n}\mathrm{}x_{h_1}\mathrm{Hom}_{}(V_{\mathrm{Out}(h_1)},V_{\mathrm{Out}(h_1)})`$ the holonomy of $`x`$ around the cycle $`c`$.
An object $`𝐌=(V,x)\mathrm{Ob}(\overline{}(Q))`$ is nilpotent if and only if the trace of the holonomy of $`x`$ around every cycle $`c`$ is equal to zero. Thus the set $`𝐄_\alpha ^{nil}=\{x𝐄_\alpha |(^\alpha ,x)(Q)\}`$ is an affine subvariety in $`𝐄_\alpha `$. One can consider functions on $`𝒳_\alpha (Q)`$ as $`𝐆_\alpha `$-invariant functions on $`𝐄_\alpha ^{nil}`$.
## 2. Convolution algebra
### 2.1. Algebraically constructible functions
#### 2.1.1.
Let $`X`$ be an algebraic variety over $``$.
A *constructible* set in $`X`$ is a set obtained from subvarieties in $`X`$ by finitely many standard set theoretic operations.
A function on $`X`$ is called *constructible* if $`X`$ has a finite partition into constructible sets such that the function is constant on each of them.
We denote by $`M(X)`$ the set of all constructible functions on $`X`$ with values in $``$. The set $`M(X)`$ is naturally a $``$-linear space.
Let $`G`$ be an algebraic group acting on $`X`$. We denote by $`M_G(X)`$ a subspace of $`M(X)`$, consisting of all $`G`$-invariant functions.
#### 2.1.2.
Let $`f:XY`$ be a morphism of algebraic varieties. Then $`f^{}`$ denotes a $``$-linear map from $`M(Y)`$ to $`M(X)`$ defined as follows:
$$(f^{}(\varphi ))(x)=\varphi (f(x)).$$
Let $`f:XY`$ be a proper morphism of algebraic varieties. Then $`f_{}`$ denotes a $``$-linear map from $`M(X)`$ to $`M(Y)`$ defined as follows \[Mac74\]:
$$(f_{}(\varphi ))(y)=\underset{a}{}a\chi (f^1(y)\varphi ^1(a)),$$
where $`\chi `$ denotes the Euler characteristic with compact support.
Note that if $`f`$ is an equivariant morphism of $`G`$-varieties, then the restriction of $`f^{}`$ (resp. $`f_{}`$) gives a $``$-linear map from $`M_G(Y)`$ to $`M_G(X)`$ (resp. from $`M_G(X)`$ to $`M_G(Y)`$).
Let $`G`$ and $`H`$ be algebraic groups, $`X`$ and $`Y`$ be algebraic varieties with $`H`$ actions, and $`f:XY`$ be a locally trivial $`H`$-equivariant principal $`G`$-bundle. Then $`f_{\mathrm{}}`$ is a $``$-linear map from $`M_{G\times H}(X)`$ to $`M_H(Y)`$ defined as follows:
$$(f_{\mathrm{}}(\varphi ))(y)=\varphi (f^1(y)).$$
The map $`f_{\mathrm{}}(\varphi )`$ is well defined because $`\varphi M_G(X)`$.
###### Proposition.
The maps $`f^{}`$, $`f_{}`$, and $`f_{\mathrm{}}`$ have the following properties:
1. $`f^{}`$, $`f_{}`$, and $`f_{\mathrm{}}`$ are functorial:
$$(fg)^{}=g^{}f^{},$$
$$(fg)_{}=f_{}g_{},$$
$$(fg)_{\mathrm{}}=f_{\mathrm{}}g_{\mathrm{}},$$
when the right-hand side if defined,
2. if $`f:XY`$ is a locally trivial $`H`$-equivariant principal $`G`$-bundle then $`f_{\mathrm{}}`$ and $`f^{}`$ are inverse to each other and give an isomorphism between $`M_{G\times H}(X)`$ and $`M_H(Y)`$.
###### Proof.
All the statements except the functoriality of $`f_{}`$ are obvious. The functoriality of $`f_{}`$ follows from properties of the Euler characteristic (namely, additivity with respect to algebraic stratifications, and multiplicativity for fiber bundles). ∎
### 2.2. Convolution algebra and Lie subalgebras
#### 2.2.1.
We introduce the following $``$-linear spaces:
$$_\alpha (Q)=M_{𝐆_\alpha }(𝐄_\alpha ^{nil}),$$
$$(Q)=\underset{\alpha _+[I]}{}_\alpha (Q)$$
Due to 1.4.8 we can consider elements of $`(Q)`$ as functions on $`𝒳(Q)`$ with values in $``$. Note however, that not every function on $`𝒳(Q)`$ belongs to $`(Q)`$.
In particular, given a representation $`𝐌`$ of $`Q`$ we sometimes say ”the characteristic function of $`[M]`$” instead of ”the characteristic function of $`𝐆_\alpha `$-orbit in $`𝐄_\alpha `$, corresponding to $`[M]`$”. In either case, one should make sure that the function in question is constructible. It is always true, for example, if the number of $`𝐆_\alpha `$-orbits is finite, or if the function is in the image of a Hall map (see 3.1).
#### 2.2.2.
For example
$$_i(Q)=E_i\text{ for }iI[I],$$
where $`E_i`$ is the function equal to $`1`$ on $`𝐄_i=𝐄_i^{nil}`$ (which is a point).
#### 2.2.3.
We endow $`(Q)`$ with a bilinear product $``$, graded by $`_+[I]`$,
(2.2.3.1)
$$_\alpha (Q)_\beta (Q)_{\alpha +\beta }(Q).$$
By linearity it is enough to define $``$-product for a pair of functions $`f_\alpha _\alpha (Q)`$ and $`f_\beta _\beta (Q)`$. Moreover $`f_\alpha f_\beta `$ should belong to $`_{\alpha +\beta }(Q)`$ for the grading property (2.2.3.1) to hold.
The product $`f_\alpha f_\beta `$ is defined as follows \[Lus91a, Lus91b\]. Consider a diagram of varieties:
(2.2.3.2)
$$𝐄_\alpha ^{nil}\times 𝐄_\beta ^{nil}\stackrel{p_1}{}𝐄^{}\stackrel{p_2}{}𝐄^{\prime \prime }\stackrel{p_3}{}𝐄_{\alpha +\beta }^{nil},$$
where the notation is as follows:
$`𝐄_\alpha ^{nil}`$, $`𝐄_\beta ^{nil}`$, $`𝐄_{\alpha +\beta }^{nil}`$ are defined in 1.4.8,
$`𝐄^{\prime \prime }`$ is the variety of all pairs $`(x,W)`$, consisting of $`x𝐄_{\alpha +\beta }^{nil}`$ and an $`x`$-stable $`I`$-graded subspace of $`^{\alpha +\beta }`$ such that $`dimW=\alpha `$,
$`𝐄^{}`$ is the variety of all quadruples $`(x,W,R^{},R^{\prime \prime })`$, where $`(x,W)𝐄^{\prime \prime }`$, $`R^{}`$ is an isomorphism $`^\alpha \stackrel{~}{}W`$, $`R^{\prime \prime }`$ is an isomorphism $`^\beta \stackrel{~}{}^{\alpha +\beta }/W`$.
$`p_2(x,W,R^{},R^{\prime \prime })=(x,W)`$,
$`p_3(x,W)=x`$,
$`p_1(x,W,R^{},R^{\prime \prime })=(x^{},x^{\prime \prime })`$, where $`x_hR_{\mathrm{Out}(h)}^{}=R_{\mathrm{In}(h)}^{}x_h^{}`$, and $`x_hR_{\mathrm{Out}(h)}^{\prime \prime }=R_{\mathrm{In}(h)}^{\prime \prime }x_h^{\prime \prime }`$ for all $`hH`$.
Note that $`p_2`$ is a principal $`𝐆_\alpha \times 𝐆_\beta `$ fibration, and $`p_3`$ is proper.
Given $`f_\alpha _\alpha (Q)`$ and $`f_\beta _\beta (Q)`$, let $`g(x_1,x_2)=f_\alpha (x_1)f_\beta (x_2)`$ be an algebraically constructible function on $`𝐄_\alpha ^{nil}\times 𝐄_\beta ^{nil}`$. By definition
$$f_\alpha f_\beta =(p_3)_{}(p_2)_{\mathrm{}}(p_1)^{}(g).$$
###### Theorem.
The space $`(Q)`$ equipped with the $``$-product is a $`_+[I]`$-graded associative algebra over $``$.
###### Proof.
Associativity follows from functorial properties Proposition of the maps $`p^{},p_{},p_{\mathrm{}}`$. ∎
#### 2.2.4.
The algebra $`(Q)`$ with the $``$-product defined above is called the *Hall algebra*. In our definition the $``$-product is the opposite of the one given in \[Lus91b\] (which goes back to Ringel \[Rin90b\]).
#### 2.2.5.
Let $`Q`$ be a quiver of finite or affine type. Consider $`(Q)`$ as a $`_+[I]`$-graded Lie algebra over $``$, with the following Lie bracket:
(2.2.5.1)
$$[f,g]=fggf.$$
We denote by $`𝔫^{}(Q)`$ the Lie subalgebra of $`(Q)`$ generated by $`\{E_i\}_{iI}`$ (see 2.2.2).
###### 2.2.6 Proposition.
$`𝔫^{}(Q)`$ has the following properties:
1. $`𝔫^{}(Q)`$ is a $`_+[I]`$-graded Lie algebra: $`𝔫^{}(Q)=_{\alpha _+[I]}𝔫_\alpha ^{}(Q)`$,
2. $`dim_{}𝔫_\alpha ^{}(Q)<\mathrm{}`$ for any $`\alpha _+[I]`$,
3. $`𝔫_i^{}(Q)=E_i`$ for $`iI_+[I]`$, where $`E_i`$ is defined in 2.2.2,
4. the generators $`E_i`$ satisfy the following relations:
$$(\mathrm{ad}(E_i))^{1a_{ij}}E_j=0\text{ for }ij\text{ },$$
where $`\mathrm{ad}(X)Y=[X,Y]`$, and $`a_{ij}=<i,j>`$.
###### Proof.
2.2.6.4 is a simple calculation using the definition of the $``$-product, all the rest is obvious. ∎
#### 2.2.7.
Given $`𝐌_1`$, $`𝐌_2`$, $`𝐌_3\mathrm{Ob}((Q))`$ we denote by $`n_{𝐌_1,𝐌_2;𝐌_3}^Q`$ the Euler characteristic with compact support of the variety $`N_{𝐌_1,𝐌_2;𝐌_\mathrm{𝟑}}^Q`$ of all subobjects $`𝐕`$ of $`𝐌_3`$ such that $`[𝐕]=[𝐌_1]`$ and $`[𝐌_3/𝐕]=[𝐌_2]`$ (we consider $`N_{𝐌_1,𝐌_2;𝐌_\mathrm{𝟑}}^Q`$ as a constructible subvariety in a product of Grassmannians):
$$n_{𝐌_1,𝐌_2;𝐌_3}^Q=\chi (N_{𝐌_1,𝐌_2;𝐌_3}^Q).$$
We denote by $`^{ind}(Q)`$ the subspace of $`(Q)`$ consisting of all $`f(Q)`$ such that $`f([𝐌])=0`$ if $`𝐌\mathrm{Ob}((Q))`$ is decomposable.
###### 2.2.8 Proposition.
1. Let $`𝐌^{}`$ and $`𝐌^{\prime \prime }`$ be two indecomposable objects of $`(Q)`$. If $`n_{𝐌^{},𝐌^{\prime \prime };𝐌}^Q0`$ then either $`𝐌`$ is indecomposable or $`𝐌=𝐌^{}𝐌^{\prime \prime }`$.
2. Let $`𝐌^{}`$ and $`𝐌^{\prime \prime }`$ be two indecomposable objects of $`(Q)`$. Then
$$n_{𝐌^{},𝐌^{\prime \prime };𝐌^{}𝐌^{\prime \prime }}^Q=n_{𝐌^{\prime \prime },𝐌^{};𝐌^{}𝐌^{\prime \prime }}^Q=\{\begin{array}{cc}1\hfill & \text{ if }𝐌^{}\text{ is not isomorphic to }𝐌^{\prime \prime }\text{ },\hfill \\ 2\hfill & \text{ if }𝐌^{}\text{ is isomorphic to }𝐌^{\prime \prime }\text{ }.\hfill \end{array}$$
3. The subspace $`^{ind}(Q)(Q)`$ is closed with respect to Lie bracket 2.2.5.1.
4. $`𝔫^{}(Q)^{ind}(Q)`$, that is if $`f𝔫^{}(Q)`$ then $`f([𝐌])=0`$ for any decomposable $`𝐌(Q)`$.
###### Proof.
2.2.8.1 and 2.2.8.2 have been proven by Ch. Riedtmann \[Rie94\].
2.2.8.3 follows from 2.2.8.1 and 2.2.8.2.
2.2.8.4 follows from 2.2.8.3 and from the fact that $`E_i^{ind}(Q)`$ for any $`iI`$. ∎
#### 2.2.9.
Because of 2.2.6.4 we have a surjective $`[I]`$-graded homomorphism of Lie algebras $`\mathrm{\Xi }^{}`$ from the Lie algebra $`𝔫(Q)`$ introduced in 1.3.2 to $`𝔫^{}(Q)`$, defined as follows:
$$\mathrm{\Xi }^{}:𝔫(Q)𝔫^{}(Q),$$
$$\mathrm{\Xi }^{}(e_i)=E_i,iI.$$
Our goal in this paper is to investigate the homomorphism $`\mathrm{\Xi }^{}`$. For this study we utilize functorial properties of the $``$-product.
## 3. Functors
### 3.1. Hall functors and Hall maps
#### 3.1.1.
Let $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$, $`Q^{}=(I^{},\mathrm{\Omega }^{},\mathrm{In}^{},\mathrm{Out}^{})`$ be quivers, $`:(Q)(Q^{})`$ be a full, faithful, exact functor, such that $`\mathrm{im}`$ is *épaisse* (closed with respect to extensions).
We denote by $`dim`$ an additive map from $`[I]`$ to $`[I^{}]`$ such that
$$dim(dim_{(Q)}𝐌)=dim_{(Q^{})}((𝐌))$$
for any $`𝐌(Q)`$. The map $`dim`$ is well-defined due to 1.4.5.
#### 3.1.2.
Let $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$, $`Q^{}=(I^{},\mathrm{\Omega }^{},\mathrm{In}^{},\mathrm{Out}^{})`$ be quivers, $`:(Q)(Q^{})`$ be a full, faithful, exact functor with *épaisse* image. We call $``$ a *Hall functor* if there exist two sets of linear maps
$$\{\varphi _i_{h^{}\mathrm{\Omega }^{}}\mathrm{Hom}_{}(^{(dim(i))_{\mathrm{Out}(h^{})}},^{(dim(i))_{\mathrm{In}(h^{})}})\}_{iI},$$
$$\{\psi _h_{h^{}\mathrm{\Omega }^{}}\mathrm{Hom}_{}(^{(dim(\mathrm{Out}(h)))_{\mathrm{Out}(h^{})}},^{(dim(\mathrm{In}(h)))_{\mathrm{In}(h^{})}})\}_{h\mathrm{\Omega }},$$
such that $``$ is given by the following action on objects of $`(Q)`$
$$((V,x))=(W,z),\text{ where }$$
(3.1.2.1)
$$W=_{iI}(V_i_{}^{dim(i)}),$$
$$z=(_{iI}(\mathrm{Id}_{V_i}\varphi _i))(_{h\mathrm{\Omega }}(x_h\psi _h)),$$
and the natural action on morphisms.
We borrow all functors we use from the theory of representations of quivers, however we want them to be given by specific formulas (3.1.2.1). The following proposition asserts that this condition is not too restrictive.
###### Proposition.
Let $`Q`$ and $`Q^{}`$ be quivers, $`:(Q)(Q^{})`$ be a full, faithful, exact functor with *épaisse* image. Then there exist a Hall functor $`:(Q)(Q^{})`$ and an automorphism $`𝒢`$ of the category $`(Q)`$ such that $``$ is naturally equivalent to $`𝒢`$.
###### Proof.
Given $`h\mathrm{\Omega }`$ let $`𝐏_h=(U,t)\mathrm{Ob}((Q))`$, where
$`U_i`$ $`=\text{ if }i=\mathrm{Out}(h)\text{ or }i=\mathrm{In}(h),`$
$`U_i`$ $`=\{0\}\text{ otherwise },`$
$`t_h`$ $`=\mathrm{Id}_{},`$
$`t_s`$ $`=0\text{ if }sh.`$
We fix an object of the form $`(^{dim(i)},\varphi _i)`$ (resp. $`(^{dim(\mathrm{Out}(h))}^{dim(\mathrm{In}(h))}`$, $`\left(\begin{array}{cc}\varphi _{\mathrm{Out}\left(h\right)}& 0\\ \psi _h& \varphi _{\mathrm{In}\left(h\right)}\end{array}\right))`$) in the isomorphism class of $`(𝐏_i)`$ for each $`iI`$ (resp. $`(𝐏_h)`$ for each $`h\mathrm{\Omega }`$). Using $`\{\varphi _i\}_{iI}`$ and $`\{\psi _h\}_{h\mathrm{\Omega }}`$ we construct a functor $`:(Q)(Q^{})`$. Namely the action of $``$ on objects is given by (3.1.2.1), and the action on morphisms is the natural one. It is easy to see that $``$ is exact, $`(𝐏_i)`$ is isomorphic to $`(𝐏_i)`$ for any $`iI`$, and the natural map $`\mathrm{Ext}_{(Q)}^1(𝐏_i,𝐏_j)\mathrm{Ext}_{(Q^{})}^1((𝐏_i),(𝐏_j))`$ is a bijection for any $`i`$, $`jI`$. Now one can deduce the statement of the proposition by induction on length of Jordan-Hölder series of an object of $`(Q)`$. ∎
#### 3.1.3.
Let $`:(Q)(Q^{})`$ be a Hall functor. As $``$ is full and faithful it induces an injective map from the set $`𝒳(Q)`$ of isomorphism classes of representations of $`Q`$ to the set $`𝒳(Q^{})`$ of isomorphism classes of representations of $`Q^{}`$ (or from the set of $`𝐆_\alpha (Q)`$-orbits in $`𝐄_\alpha ^{nil}(Q)`$ to the set of $`𝐆_{dim(\alpha )}(Q^{})`$-orbits in $`𝐄_{dim(\alpha )}^{nil}(Q^{})`$). Therefore one can define a push-forward map $`𝔣`$ from $`(Q)`$ to the set of $``$-valued functions on $`𝒳(Q^{})`$ as follows
$$(𝔣(g))([(𝐌)])=g([𝐌]),$$
$$(𝔣(g))([𝐗])=0\text{ if }\mathrm{}𝐍\mathrm{Ob}((Q))\text{ such that }𝐗=(𝐍)\text{ }.$$
We call $`𝔣`$ the *Hall map* associated with the Hall functor $``$.
###### Proposition.
Let $`:(Q)(Q^{})`$ be a Hall functor, and $`𝔣`$ be the corresponding Hall map. Then
1. $`\mathrm{im}𝔣(Q^{})`$, that is $`𝔣`$ maps constructible $`𝐆_\alpha (Q)`$-invariant functions on the variety $`𝐄_\alpha ^{nil}(Q)`$ to constructible $`𝐆_{dim(\alpha )}(Q^{})`$-invariant functions on the variety $`𝐄_{dim(\alpha )}^{nil}(Q^{})`$,
2. $`𝔣:(Q)(Q^{})`$ is a homomorphism of algebras:
$$𝔣(g)𝔣(h)=𝔣(gh)\text{ for any }g\text{}h(Q).$$
###### Proof.
Follows from definitions. Note that being a Hall functor $``$ induces algebraic maps between varieties $`𝐄_\alpha ^{nil}`$, $`𝐄^{}`$, $`𝐄^{\prime \prime }`$ (used in (2.2.3)) for the quiver $`Q`$ and the corresponding varieties for the quiver $`Q^{}`$. ∎
#### 3.1.4.
Let $`:(Q)(Q^{})`$ be a Hall functor, in particular, $``$ is full and faithful. Then $`(𝐌)`$ is indecomposable if and only if $`𝐌`$ is indecomposable. It follows that $``$ induces an injective map from the set $`𝒯(Q)`$ of isomorphism classes of indecomposable objects of $`(Q)`$ to the set $`𝒯(Q^{})`$ of isomorphism classes of indecomposable objects of $`(Q^{})`$.
#### 3.1.5.
The following proposition describes functorial properties of the Euler cocycle $`ϵ`$ and the bilinear form $`<,>`$ with respect to a Hall functor.
###### Proposition.
Let $`:(Q)(Q^{})`$ be a Hall functor. Then
$$\begin{array}{cc}\hfill ϵ_Q(dim_{(Q)}𝐌_1,dim_{(Q)}𝐌_2)& =ϵ_Q^{}(dim_{(Q^{})}(𝐌_1),dim_{(Q^{})}(𝐌_2)),\hfill \\ \hfill <dim_{(Q)}𝐌_1,dim_{(Q)}𝐌_2>_Q& =<dim_{(Q^{})}(𝐌_1),dim_{(Q^{})}(𝐌_2)>_Q^{}.\hfill \end{array}$$
for any $`𝐌_1`$, $`𝐌_2\mathrm{Ob}((Q))`$.
###### Proof.
The proposition follows from 1.4.6 and from the fact that $``$ is full, faithful, exact, and has *épaisse* image. ∎
#### 3.1.6.
The simplest example of a Hall functor is an embedding functor.
We call a subquiver $`Q^{}=(I^{},\mathrm{\Omega }^{},\mathrm{In}^{},\mathrm{Out}^{})`$ of a quiver $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$ *full* if any $`h\mathrm{\Omega }`$ such that $`\mathrm{In}(h)I^{}`$ and $`\mathrm{Out}(h)I^{}`$ belongs to $`\mathrm{\Omega }^{}`$.
Let $`Q^{}`$ be a full subquiver of $`Q`$. Then the natural embedding $`_{Q^{}Q}:(Q^{})(Q)`$ clearly satisfies all the axioms of a Hall functor, and therefore induces an injective homomorphism $`𝔦_{Q^{}Q}:(Q^{})(Q)`$. Moreover,
$$𝔦_{Q^{}Q}(E_k)=E_k,$$
where $`k`$ is a vertex of $`Q^{}`$. It follows that $`𝔦_{Q^{}Q}`$ induces an injective homomorphism from $`𝔫^{}(Q^{})`$ to $`𝔫^{}(Q)`$.
In the next section we consider a more intricate example of an (almost) Hall functor – the reflection functor, and in Section 5 we have still more examples.
### 3.2. Reflection functors
#### 3.2.1.
Let $`(I,<,>)`$ be a Cartan datum. We denote by $`\sigma _i`$ an endomorphism of $`[I]`$ given by the reflection with respect to a vertex $`iI`$:
$$\sigma _i(\alpha )=\alpha <i,\alpha >i.$$
It is easy to see that $`\sigma _i`$ preserves the inner product $`<,>`$, and that $`\sigma _i\sigma _i=id`$. The subgroup of the group of endomorphisms of $`[I]`$ generated by $`\{\sigma _i\}_{iI}`$ is called the *Weyl group*. We denote the Weyl group by $`W`$.
#### 3.2.2.
Let $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$ be a quiver with the underlying Cartan datum $`(I,<,>)`$.
A vertex $`iI`$ is called a *source* if there is no $`h\mathrm{\Omega }`$ such that $`\mathrm{In}(h)=i`$.
A vertex $`iI`$ is called a *sink* if there is no $`h\mathrm{\Omega }`$ such that $`\mathrm{Out}(h)=i`$.
A vertex is called *admissible* if it is either a source or a sink.
A reflection $`\sigma _i`$ is called *admissible* if $`i`$ is an admissible vertex.
We denote by $`\sigma _iQ=(I,\mathrm{\Omega },\mathrm{In}^{},\mathrm{Out}^{})`$ a quiver with the same sets of vertices and edges and with the maps $`\mathrm{In}^{}`$ and $`\mathrm{Out}^{}`$ defined as follows:
$$\mathrm{In}^{}(h)=\mathrm{In}(h),\mathrm{Out}^{}(h)=\mathrm{Out}(h)\text{ if }\mathrm{In}(h)i\text{ and }\mathrm{Out}(h)i,$$
$$\mathrm{In}^{}(h)=\mathrm{Out}(h),\mathrm{Out}^{}(h)=\mathrm{In}(h)\text{ if }\mathrm{In}(h)=i\text{ or }\mathrm{Out}(h)=i.$$
In other words, $`\sigma _iQ`$ is obtained from $`Q`$ by reversing all the arrows coming to $`i`$ and going out of $`i`$.
#### 3.2.3.
We call a sequence of vertices $`i_1,\mathrm{},i_KI`$ *admissible* if for any $`l\{1,\mathrm{},K\}`$ the vertex $`i_l`$ is a sink for the quiver $`\sigma _{i_{l+1}}\mathrm{}\sigma _{i_K}Q`$.
We call $`c=\sigma _{i_1}\mathrm{}\sigma _{i_{|I|}}W`$ a *Coxeter element* if $`i_1,\mathrm{},i_{|I|}`$ is an admissible sequence of vertices in which each vertex is present exactly once (in other words it is an ordering of vertices). It is easy to see that if the quiver has no oriented cycles then there exists a unique Coxeter element.
It is clear that $`cQ=Q`$.
#### 3.2.4.
The following proposition asserts that the Euler cocycle $`ϵ`$ introduced in 1.3.4 is functorial with respect to admissible reflections.
###### Proposition.
Let $`i`$ be an admissible vertex for $`Q`$. Then
$$ϵ_{\sigma _iQ}(\sigma _i\alpha ,\sigma _i\beta )=ϵ_Q(\alpha ,\beta ),$$
for any $`\alpha `$, $`\beta [I]`$.
###### Proof.
An easy calculation using the definition of $`ϵ_Q`$. See also Remark 3.2.7. ∎
#### 3.2.5.
For an admissible $`iI`$ we denote by $`{}_{}{}^{i}(Q)`$ the full subcategory of $`(Q)`$ defined as follows. Let $`i`$ be a source (resp. sink) and $`(V,x)`$ be an object of $`(Q)`$. Then it is an object of $`{}_{}{}^{i}(Q)`$ if $`_{h|\mathrm{Out}(h)=i}x_h:V_i_{h|\mathrm{Out}(h)=i}V_{\mathrm{In}(h)}`$ is injective (resp. $`_{h|\mathrm{In}(h)=i}x_h:_{h|\mathrm{In}(h)=i}V_{\mathrm{Out}(h)}V_i`$ is surjective). Let us note that the subcategory $`{}_{}{}^{i}(Q)`$ is not abelian in general.
###### Proposition.
1. The subcategory $`{}_{}{}^{i}(Q)`$ is *épaisse*.
2. $`\mathrm{Ob}({}_{}{}^{i}(Q))`$ contains all indecomposable objects of $`(Q)`$ except for the simple object $`𝐏_i`$.
###### Proof.
Follows from the definition. ∎
#### 3.2.6.
Let $`i`$ be an admissible vertex. Following \[BGP73\] we define *reflection functor* $`𝒮_i(Q):(Q)(\sigma _iQ)`$.
If $`i`$ is a sink then the action of the functor $`𝒮_i(Q)`$ on objects is defined by $`𝒮_i(Q)((V,x))=(V^{},x^{})`$, where
$$V_k^{}=V_k\text{ if }ki,$$
$$V_i^{}=\mathrm{Ker}(\underset{h|\mathrm{In}(h)=i}{}x_h:\underset{h|\mathrm{In}(h)=i}{}V_{\mathrm{Out}(h)}V_i),$$
$$x_h^{}=x_h\text{ if }\mathrm{In}(h)i,$$
$$x_h^{}:V_i^{}V_{\mathrm{Out}(h)}\text{ is the inclusion composed with the projection if }\mathrm{In}(h)=i.$$
If $`i`$ is a source then the action of the functor $`𝒮_i(Q)`$ on objects is defined by $`𝒮_i(Q)((V,x))=(V^{},x^{})`$, where
$$V_k^{}=V_k\text{ if }ki,$$
$$V_i^{}=\mathrm{Coker}(\underset{h|\mathrm{Out}(h)=i}{}x_h:V_i\underset{h|\mathrm{Out}(h)=i}{}V_{\mathrm{In}(h)}),$$
$$x_h^{}=x_h\text{ if }\mathrm{Out}(h)i,$$
$$x_h^{}:V_{\mathrm{In}(h)}V_i^{}\text{ is the inclusion composed with the projection if }\mathrm{Out}(h)=i.$$
The action of $`𝒮_i(Q)`$ on morphisms is the natural one.
###### Proposition.
The functor $`𝒮_i(Q)`$ has the following properties:
1. The image of $`𝒮_i(Q)`$ coincides with the subcategory $`{}_{}{}^{i}(\sigma _iQ)`$.
2. The restriction of the functor $`𝒮_i(Q)`$ defines an equivalence of categories: $`𝒮_i(Q):{}_{}{}^{i}(Q)\stackrel{~}{}{}_{}{}^{i}(\sigma _iQ)`$, the inverse functor is $`𝒮_i(\sigma _iQ)`$.
3. An object $`𝐌\mathrm{Ob}({}_{}{}^{i}(Q))`$ is indecomposable if and only if $`𝒮_i(Q)𝐌`$ is indecomposable.
4. Let $`𝐌\mathrm{Ob}({}_{}{}^{i}(Q))`$, $`dim_{(Q)}𝐌=\alpha `$. Then $`dim_{(\sigma _iQ)}𝒮_i(𝐌)=\sigma _i(\alpha )`$.
5. If
$$0𝐌_1𝐌_2𝐌_30$$
is an exact sequence of objects of $`{}_{}{}^{i}(Q)`$ then
$$0𝒮_i(Q)(𝐌_1)𝒮_i(Q)(𝐌_2)𝒮_i(Q)(𝐌_3)0$$
is exact in $`{}_{}{}^{i}(\sigma _iQ)`$.
###### Proof.
3.2.6.1 and 3.2.6.4 follow from the definition of $`𝒮_i(Q)`$,
3.2.6.2 Let $`i`$ be a sink. Then there is the following split exact sequence
$$0𝒮_i(\sigma _iQ)𝒮_i(Q)(𝐌)𝐌𝐋_i0,$$
for any $`𝐌(Q)`$. Here $`𝐋_i=(V,x)`$, where
$$V_k=0\text{ if }ki,$$
$$V_i=\mathrm{Coker}(\underset{h|\mathrm{In}(h)=i}{}x_h:\underset{h|\mathrm{In}(h)=i}{}V_{\mathrm{Out}(h)}V_i),$$
$$x=0.$$
If $`𝐌{}_{}{}^{i}(Q)`$ then $`𝐋_i=0`$.
The proof for $`i`$ being a source is analogous.
3.2.6.3 follows from 3.2.6.2.
3.2.6.5 follows from the Snake Lemma.
#### 3.2.7.
*Remark.* Roughly speaking Proposition Proposition means that the restriction of the functor $`𝒮_i`$ to the subcategory $`{}_{}{}^{i}(Q)`$ has associated dimension mapping $`dim𝒮_i=\sigma _i`$ and has properties of a Hall functor. In particular together with 3.1.5 it explains 3.2.4. However one should be careful, as the subcategory $`{}_{}{}^{i}(Q)`$ is not abelian. It would be more accurate to define the reflection functor $`𝒮_i`$ not on a subcategory, but on a quotient category (in the sense of Serre) – namely, on the quotient of $`(Q)`$ over the embedded category of representations of the quiver with one vertex $`i`$ and no edges.
Authors decided not to go into the treatment of Hall mappings for quotient categories in order to save space. Interested reader can easily invent all the necessary notions.
#### 3.2.8.
Though $`𝒮_i`$ is not strictly speaking a Hall functor, one can still define the corresponding Hall map $`𝔰_i`$ if one restricts its domain.
We denote by $`{}_{}{}^{i}𝐄_{\alpha }^{nil}`$ the subvariety of $`𝐄_\alpha ^{nil}`$ given by $`{}_{}{}^{i}𝐄_{\alpha }^{nil}=\{x𝐄_\alpha ^{nil}|(^\alpha ,x){}_{}{}^{i}(Q)\}`$.
Let $`{}_{}{}^{i}_{\alpha }^{}(Q)`$ be a subspace of $`_\alpha (Q)`$ consisting of all functions with support in $`{}_{}{}^{i}𝐄_{\alpha }^{nil}`$, and $`{}_{}{}^{i}(Q)=_{\alpha _+[I]}{}_{}{}^{i}_{\alpha }^{}(Q)`$. It follows from 3.2.5.1 that $`{}_{}{}^{i}(Q)`$ is a subalgebra of $`(Q)`$ with respect to $``$-product, and a Lie subalgebra with respect to the Lie bracket (2.2.5.1).
Now one can repeat the construction of the Hall map in 3.1.3 for the functor $`𝒮_i`$ using Proposition Proposition, and get an isomorphism of algebras
$$𝔰_i:{}_{}{}^{i}(Q){}_{}{}^{i}(\sigma _iQ),$$
such that
$$(𝔰_i(g))([𝒮_i(𝐌)])=g([𝐌])$$
for $`g{}_{}{}^{i}(Q)`$.
## 4. Quiver of finite type
#### 4.0.0.
Throughout this chapter $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$ denotes a quiver of finite type with underlying Cartan datum $`(I,<,>)`$.
### 4.1. Root system for a Cartan datum of finite type
#### 4.1.1.
An element $`\alpha `$ of $`[I]`$ is called a *root* if $`<\alpha ,\alpha >=2`$. We denote the set of all roots by $`R`$. Let $`R_+`$ denote the set of *positive roots* : $`R_+=R_+[I]`$.
For example elements of $`I[I]`$ are (positive) roots. We call $`iIR_+`$ a *simple root*.
Since the symmetric form $`<,>`$ is positive definite, the set of roots is finite. Therefore the Weyl group $`W`$ (which preserves the root system $`R`$) is a finite group.
### 4.2. Lie algebra based on the Euler cocycle
#### 4.2.1.
We denote by $`𝔫^ϵ(Q)`$ a $``$-linear space with basis $`\{\stackrel{~}{e}_\alpha \}_{\alpha R_+}`$ equipped with the bilinear bracket
(4.2.1.1)
$$[\stackrel{~}{e}_\alpha ,\stackrel{~}{e}_\beta ]=\{\begin{array}{cc}ϵ(\alpha ,\beta )\stackrel{~}{e}_{\alpha +\beta }\hfill & \text{ if }\alpha +\beta R_+,\hfill \\ 0\hfill & \text{ if }\alpha +\beta R_+,\hfill \end{array}$$
where the Euler cocycle $`ϵ`$ is defined in 1.3.4.
###### Theorem (Frenkel, Kac \[FK80\], Segal \[Seg81\]).
The space $`𝔫^ϵ(Q)`$ equipped with the bracket (4.2.1.1) is a Lie algebra over $``$.
###### Proof.
One has to check that the bracket is skew-symmetric and satisfies the Jacobi identity. Both statements follow from the definition of $`ϵ`$. ∎
###### 4.2.2 Proposition.
The Lie algebra $`𝔫^ϵ(Q)`$ is generated by $`\{\stackrel{~}{e}_i\}_{iI}`$, the following relations hold in $`𝔫^ϵ(Q)`$ :
$`[\stackrel{~}{e}_i,\stackrel{~}{e}_j]=0`$ $`\text{if }<i,j>=0\text{ },`$
$`[\stackrel{~}{e}_i,[\stackrel{~}{e}_i,\stackrel{~}{e}_j]]=0`$ $`\text{if }<i,j>=1\text{ }.`$
###### Proof.
Induction using the fact that for any $`\alpha R_+`$ there exists $`iI`$ such that $`(\alpha i)(R_+0)`$ \[Bou68\]. ∎
#### 4.2.3.
Because of Proposition 4.2.2 we have a surjective homomorphism of Lie algebras
$$\mathrm{\Xi }^ϵ:𝔫(Q)𝔫^ϵ(Q)$$
induced by
$$\mathrm{\Xi }^ϵ(e_i)=\stackrel{~}{e}_i,$$
where $`𝔫(Q)`$ is introduced in 1.3.2.
It is known that $`𝔫(Q)`$ is isomorphic to the nilpotent radical of a Borel subalgebra of the semisimple Lie algebra associated to the Dynkin graph underlying the quiver $`Q`$. In particular $`dim_{}𝔫(Q)=|R_+|=dim_{}𝔫^ϵ(Q)`$.
Comparison of dimensions implies that the surjective homomorphism $`\mathrm{\Xi }^ϵ`$ is actually an isomorphism and $`\mathrm{\Xi }_Q=\mathrm{\Xi }^{}(\mathrm{\Xi }^ϵ)^1`$ is a well-defined, surjective, $`[I]`$-graded homomorphism from $`𝔫^ϵ(Q)`$ to $`𝔫^{}(Q)`$.
The following commutative diagram shows all the Lie algebras involved and the corresponding $`[I]`$-graded homomorphisms
Our goal is to study the map $`\mathrm{\Xi }_Q`$.
### 4.3. Lie algebra $`𝔫^{}`$ for a quiver of finite type
#### 4.3.1.
The following theorem is a combination of results of Gabriel \[Gab72\], and Bernstein, Gelfand, and Ponomarev \[BGP73\].
###### Theorem.
Let $`Q`$ be a quiver of finite type.
1. If $`\alpha R_+`$ then there are no indecomposable objects in $`(Q)`$ with dimension $`\alpha `$ (in other words, $`𝒯_\alpha `$ is empty).
2. There is a unique (up to an isomorphism) indecomposable object $`𝐏_\alpha `$ in $`(Q)`$ with dimension $`\alpha `$ for any $`\alpha R_+`$ In other words, $`𝒯_\alpha `$ consists of a single point. We denote the characteristic function of $`[𝐏_\alpha ]`$ by $`E_\alpha _\alpha (Q)`$.
3. For any given $`\alpha R_+`$ there exists an admissible sequence $`\{i_t\}_{t=0}^N`$ of vertices of $`Q`$ (depending on $`\alpha `$), such that
* $`\alpha =\sigma _{i_N}\sigma _{i_{N1}}\mathrm{}\sigma _{i_1}i_0`$,
* $`𝐏_\alpha =𝒮_{i_N}𝒮_{i_{N1}}\mathrm{}𝒮_{i_1}𝐏_{i_0}`$,
* $`E_\alpha =𝔰_{i_N}𝔰_{i_{N1}}\mathrm{}𝔰_{i_1}E_{i_0}`$,
* for any $`t\{1,\mathrm{},N\}`$ the element $`\sigma _{i_t}\sigma _{i_{t1}}\mathrm{}\sigma _{i_1}i_0`$ of $`R`$ belongs to $`R_+`$.
Note that the notation $`𝐏_\alpha `$ and $`E_\alpha `$ is consistent with the notation $`𝐏_i`$ and $`E_i`$ which we use above.
#### 4.3.2. Remark
It follows from Theorem 4.3.1 that in the case of a quiver of finite type there are finitely many non-isomorphic objects of $`(Q)`$ in each dimension. In other words, there are finitely many $`𝐆_\alpha `$-orbits in $`𝐄_\alpha `$ for each $`\alpha `$. Therefore $`_\alpha (Q)`$ is finite-dimensional with a basis consisting of characteristic functions of the orbits. Thus $`(Q)`$ can be thought of as a $``$-linear space with a basis given by the set $`𝒳(Q)`$ of isomorphism classes of objects of $`(Q)`$.
The $``$-product can be rewritten as
$$[𝐌_1][𝐌_2]=\underset{[𝐌_3]𝒳(Q)}{}n_{𝐌_1,𝐌_2;𝐌_3}^Q[𝐌_3],$$
where $`𝐌_1`$, $`𝐌_2`$, $`𝐌_3`$ are representatives of isomorphism classes and $`n_{𝐌_1,𝐌_2;𝐌_3}^Q`$ is defined in 2.2.7.
#### 4.3.3.
Our main result for quivers of finite type is the following
###### Theorem.
Let $`Q`$ be a quiver of finite type. Then $`\mathrm{\Xi }_Q(\stackrel{~}{e}_\alpha )=E_\alpha `$ for any $`\alpha R_+`$.
#### 4.3.4.
Before proving Theorem Theorem we give the following immediate Corollary:
###### Corollary (Ringel \[Rin90c\]).
1. The $``$-linear space $`𝔫^{}(Q)`$ coincides with the space of all $``$-valued functions on the set $`𝒯`$ of isomorphism classes of indecomposable representations.
2. The Lie bracket in $`𝔫^{}(Q)`$ is given by the Euler cocycle $`ϵ`$:
$$[E_\alpha ,E_\beta ]=\{\begin{array}{cc}ϵ(\alpha ,\beta )E_{\alpha +\beta },\hfill & \text{ if }\alpha +\beta R_+,\hfill \\ 0\hfill & \text{ if }\alpha +\beta R_+.\hfill \end{array}$$
#### 4.3.5. Remark
Statement Corollary was first established by C. M. Ringel \[Rin90c\], who used slightly different definition of the Lie algebra $`𝔫^{}(Q)`$. He found all the varieties $`N_{𝐌_1,𝐌_2;𝐌_3}^Q`$ used in the definition of the structure constants $`n_{𝐌_1,𝐌_2;𝐌_3}^Q`$ (see 4.3.2).
Our proof does not rely on the explicit form of the varieties $`N_{𝐌_1,𝐌_2;𝐌_3}^Q`$.
### 4.4. Reflections revisited or a proof of the Ringel theorem
#### 4.4.1.
Let $`i`$ be an admissible vertex. We denote by $`{}_{}{}^{i}𝔫_{}^{}(Q)`$ the intersection of $`𝔫^{}(Q)`$ with $`{}_{}{}^{i}(Q)`$ (see 3.2.8). Note that $`E_i{}_{}{}^{i}𝔫_{}^{}(Q)`$.
Let $`j`$ be a vertex connected with $`i`$ by an edge (we recall that $`Q`$ being of finite type is simply laced). Let $`E_{i+j}{}_{}{}^{i}(Q)`$ be a function on $`𝐄_{i+j}`$ equal to the characteristic function of the $`𝐆_{i+j}`$-orbit, corresponding to the unique, up to an isomorphism, indecomposable representation $`𝐏_{i+j}`$ of dimension $`i+j[I]`$.
It follows from the definition of the $``$-product that $`E_{i+j}=[E_j,E_i]`$ if $`i`$ is a source, and $`E_{i+j}=[E_i,E_j]`$ if $`i`$ is a sink. In particular, $`E_{i+j}{}_{}{}^{i}𝔫_{}^{}(Q)`$.
###### Proposition.
1. Let $`\alpha R_+`$, $`\alpha i`$. We denote by $`E_\alpha ^{}{}_{}{}^{i}(Q)`$ and $`E_\alpha ^{\prime \prime }{}_{}{}^{i}(\sigma _iQ)`$ the characteristic functions defined in 4.3.1.2 for quivers $`Q`$ and $`\sigma _iQ`$ respectively. Then
$$𝔰_i(E_\alpha ^{})=E_{\sigma _i\alpha }^{\prime \prime }.$$
2. The subspace $`{}_{}{}^{i}𝔫_{}^{}(Q)`$ is a $`_+[I]`$-graded Lie subalgebra of $`𝔫^{}(Q)`$ generated by elements $`E_k`$ for $`ki`$ and $`E_{i+j}`$ for all $`j`$ connected to $`i`$ by an edge.
3. The restriction of $`𝔰_i`$ to $`{}_{}{}^{i}𝔫_{}^{}(Q)`$ yields an isomorphism of Lie algebras $`{}_{}{}^{i}𝔫_{}^{}(Q)`$ and $`{}_{}{}^{i}𝔫_{}^{}(\sigma _iQ)`$.
###### Proof.
Note that $`\sigma _i(j)=i+j`$ because $`Q`$ is simply laced. Now the proposition follows from 3.2.8, 3.2.6.3, and 3.2.6.4. The action of $`𝔰_i`$ on the generators of $`{}_{}{}^{i}𝔫_{}^{}(Q)`$ is given by
$$𝔰_i(E_j^{})=E_j^{\prime \prime }\text{ if }i\text{ is not connected with }j\text{ by an edge },$$
$$𝔰_i(E_{i+j}^{})=E_j^{\prime \prime },𝔰_i(E_j^{})=E_{i+j}^{\prime \prime }\text{ if }i\text{ is connected with }j\text{ by an edge },$$
where the notation is as in 4.4.1.1. ∎
#### 4.4.2.
For an admissible vertex $`i`$ we denote by $`{}_{}{}^{i}𝔫_{}^{ϵ}`$ the Lie subalgebra of $`𝔫^ϵ`$, generated by elements $`\{\stackrel{~}{e}_k\}_{\begin{array}{c}kI\\ ki\end{array}}`$ and $`\{\stackrel{~}{e}_{i+j}\}_{\begin{array}{c}jI\\ i+jR_+\end{array}}`$ (cf. 4.4.1). Note that $`i+jR_+`$ if and only if $`j`$ is connected to $`i`$ by an edge. One has $`𝔫^ϵ={}_{}{}^{i}𝔫_{}^{ϵ}\stackrel{~}{e}_i`$ as a $``$-linear space, and $`\stackrel{~}{e}_{i+j}=[\stackrel{~}{e}_j,\stackrel{~}{e}_i]`$ (resp. $`\stackrel{~}{e}_{i+j}=[\stackrel{~}{e}_i,\stackrel{~}{e}_j]`$) if $`i`$ is a source (resp. a sink).
#### 4.4.3.
We denote by $`𝔰_i^ϵ`$ a $``$-linear map $`{}_{}{}^{i}𝔫_{}^{ϵ}(Q){}_{}{}^{i}𝔫_{}^{ϵ}(\sigma _iQ)`$ given by $`𝔰_i^ϵ(\stackrel{~}{e}_\alpha ^{})=\stackrel{~}{e}_{\sigma _i(\alpha )}^{\prime \prime }`$, where $`\stackrel{~}{e}_\alpha ^{}`$, $`\stackrel{~}{e}_\alpha ^{\prime \prime }`$ are basic vectors of $`{}_{}{}^{i}𝔫_{}^{ϵ}(Q)`$ and $`{}_{}{}^{i}𝔫_{}^{ϵ}(\sigma _iQ)`$ respectively. In particular, for the generators
$$𝔰_i^ϵ(\stackrel{~}{e}_j^{})=\stackrel{~}{e}_j^{\prime \prime }\text{ if }i\text{ is not connected with }j\text{ by an edge },$$
$$𝔰_i^ϵ(\stackrel{~}{e}_{i+j}^{})=\stackrel{~}{e}_j^{\prime \prime },𝔰_i^ϵ(\stackrel{~}{e}_j^{})=\stackrel{~}{e}_{i+j}^{\prime \prime }\text{ if }i\text{ is connected with }j\text{ by an edge }.$$
###### Proposition.
The map $`𝔰_i^ϵ`$ is well-defined and is an isomorphism of Lie algebras.
###### Proof.
Follows from the fact that $`\sigma _i(R_+\backslash \{i\})=R_+\backslash \{i\}`$ and from 3.2.4. ∎
#### 4.4.4.
It follows from 4.4.1, 4.4.2 that $`\mathrm{\Xi }_Q(^i𝔫^ϵ(Q)){}_{}{}^{i}𝔫_{}^{}(Q)`$, and that $`\mathrm{\Xi }_Q(\stackrel{~}{e}_{i+j})=E_{i+j}`$.
Using Propositions 4.4.1.3 and Proposition we get the following diagram of Lie algebra homomorphisms:
(4.4.4.1)
where the horizontal maps are isomorphisms. The diagram (4.4.4.1) is commutative on generators $`\stackrel{~}{e}_k`$, $`\stackrel{~}{e}_{i+j}`$ of $`{}_{}{}^{i}𝔫_{}^{ϵ}(Q)`$ and, therefore, is commutative.
#### 4.4.5. Proof of Theorem Theorem
Let $`G(Q)R_+`$ be the set of all $`\alpha R_+`$ such that $`\mathrm{\Xi }_Q(\stackrel{~}{e}_\alpha )=E_\alpha `$. Then $`IG(Q)`$ for any $`Q`$, and it follows from 4.4.1.1, 4.4.3, and the commutative diagram (4.4.4.1) that $`\sigma _i(G(Q)\backslash \{i\})=G(\sigma _iQ)\backslash \{i\}`$ for any admissible vertex $`i`$. Using the sequence of reflections (4.3.1.3) one obtains that $`\alpha G(Q)`$ for any $`\alpha R_+`$. Therefore $`\mathrm{\Xi }_Q(\stackrel{~}{e}_\alpha )=E_\alpha `$ for any $`\alpha R_+`$ , which proves the theorem. ∎
## 5. Quiver of affine type
#### 5.0.0.
In this chapter we consider the Lie algebra $`𝔫^{}(Q)`$ in the case of $`Q`$ being of affine type.
Throughout the chapter $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$ denotes a quiver of affine type with underlying Cartan datum $`(I,<,>)`$. In particular the bilinear form $`<,>`$ is positive semi-definite.
### 5.1. Root system for a Cartan datum of affine type
#### 5.1.1.
An element $`\alpha [I]`$ is called a *root* if $`<\alpha ,\alpha >2`$ and $`\alpha 0`$. We denote the set of all roots by $`R`$. Let $`R_+`$ denote the set of *positive roots* : $`R_+=R_+[I]`$.
For example elements of $`I[I]`$ are (positive) roots. We call $`iIR_+`$ a *simple* root.
It is easy to see that if $`\alpha R`$ then $`<\alpha ,\alpha >\{0,2\}`$. A root $`\alpha `$ such that $`<\alpha ,\alpha >=2`$ (resp. $`<\alpha ,\alpha >=0`$ ) is called *real* (resp. *imaginary*) root. We denote the set of all positive real (resp. positive imaginary) roots by $`R_+^{\mathrm{re}}`$ (resp. $`R_+^{\mathrm{im}}`$). So $`R_+=R_+^{\mathrm{re}}R_+^{\mathrm{im}}`$.
It is known that $`R_+^{\mathrm{im}}=(_+\backslash \{0\})\delta `$, where $`\delta R_+`$. We call $`\delta `$ the *first imaginary root*.
The components of $`\delta `$ are given in the table of affine Dynkin graphs (Figure 2).
#### 5.1.2.
A vertex $`pI`$ is called an *extending* vertex if $`\delta _p=1`$, where $`\delta `$ is the first imaginary root. Such a vertex always exists, but it is not unique in general. If one removes from $`Q`$ an extending vertex together with adjacent edges one gets an irreducible quiver of finite type.
### 5.2. Lie algebra based on the Euler cocycle
#### 5.2.1.
We consider the following set of $``$-linear spaces:
* for each positive real root $`\alpha `$ a one-dimensional $``$-linear space $`𝔫_\alpha ^ϵ(Q)=\stackrel{~}{e}_\alpha `$ with generator $`\stackrel{~}{e}_\alpha `$,
* for each positive imaginary root $`n\delta `$ a $``$-linear space $`𝔫_{n\delta }^ϵ(Q)=[I]/\delta `$, where we consider $`\delta `$ as an element of $`R_+[I][I]`$. For $`h[I]`$ we denote by $`h(n)`$ the image of $`h`$ under the natural projection map $`[I]𝔫_{n\delta }^ϵ(Q)`$.
We denote by $`𝔫^ϵ(Q)`$ the following $`R_+`$-graded $``$-linear space:
$$𝔫^ϵ(Q)=\underset{\alpha R_+}{}𝔫_\alpha ^ϵ(Q),$$
equipped with the bilinear bracket
(5.2.1.1)
$$\begin{array}{cc}\hfill [\stackrel{~}{e}_\alpha ,\stackrel{~}{e}_\beta ]& =\{\begin{array}{cc}ϵ(\alpha ,\beta )\stackrel{~}{e}_{\alpha +\beta }\hfill & \text{ if }\alpha +\beta R_+^{\mathrm{re}},\hfill \\ ϵ(\alpha ,\beta )\alpha (k)\hfill & \text{ if }\alpha +\beta =k\delta ,\hfill \\ 0\hfill & \text{ if }\alpha +\beta R_+,\hfill \end{array}\hfill \\ \hfill [h(n),\stackrel{~}{e}_\alpha ]=[\stackrel{~}{e}_\alpha ,h(n)]& =ϵ(n\delta ,\alpha )<h,\alpha >\stackrel{~}{e}_{\alpha +n\delta },\hfill \\ \hfill [h(n),h(m)]& =0,\hfill \end{array}$$
where $`h[I]`$, $`ϵ`$ is the Euler cocycle (see 1.3.4). Note that $`<h^{},\alpha >=<h^{\prime \prime },\alpha >`$ if $`h^{}(n)=h^{\prime \prime }(n)`$ because $`\delta \mathrm{Ker}<,>`$. Therefore the bracket (5.2.1.1) is well-defined.
###### Theorem.
The $``$-linear space $`𝔫^ϵ(Q)`$ equipped with the bracket (5.2.1.1) is a Lie algebra.
###### Proof.
One has to check that the bracket is skew-symmetric and satisfies the Jacobi identity. Both statements follow (after lengthy, but straightforward calculations) from the definition of $`ϵ`$. A more conceptual proof is to consider $`𝔫^ϵ`$ as embedded into an algebra of vertex operators, associated with the quiver $`Q`$ (see \[Fre85, MRY90\]). ∎
###### 5.2.2 Proposition.
The Lie algebra $`𝔫^ϵ(Q)`$ is generated by $`\{\stackrel{~}{e}_i\}_{iI}`$, the following relations hold in $`𝔫^ϵ(Q)`$ :
$$(\mathrm{ad}(\stackrel{~}{e}_i))^{1a_{ij}}\stackrel{~}{e}_j=0\text{ if }ij\text{ },$$
where $`\mathrm{ad}(x)y=[x,y]`$, and $`a_{ij}=<i,j>`$.
###### Proof.
Induction on $`n(\alpha )`$, where $`n(\alpha )=_{iI}\alpha _i`$ for a root $`\alpha `$. ∎
#### 5.2.3.
Because of Proposition 5.2.2 we have a surjective $`R_+`$-graded homomorphism of Lie algebras:
$$\mathrm{\Xi }^ϵ:𝔫(Q)𝔫^ϵ(Q)$$
induced by
$$\mathrm{\Xi }^ϵ(e_i)=\stackrel{~}{e}_i,$$
where $`𝔫(Q)`$ is the Lie algebra introduced in 1.3.2.
It is known that
$$dim_{}𝔫_\alpha (Q)=\{\begin{array}{cc}1\hfill & \text{ if }\alpha R^{\mathrm{re}}\text{ }\hfill \\ |I|1\hfill & \text{ if }\alpha R^{\mathrm{im}}\text{ }\hfill \end{array}\}=dim_{}𝔫^ϵ_\alpha (Q).$$
Comparison of dimensions implies that the surjective homomorphism $`\mathrm{\Xi }^ϵ`$ is actually an isomorphism and $`\mathrm{\Xi }_Q=\mathrm{\Xi }^{}(\mathrm{\Xi }^ϵ)^1`$ is a well-defined, surjective, $`[I]`$-graded homomorphism from $`𝔫^ϵ(Q)`$ to $`𝔫^{}(Q)`$.
As in the case of a quiver of finite type we want to study the map $`\mathrm{\Xi }_Q`$ and, in particular, its image. There are two (related) differences with the finite case. First, the description of indecomposable objects in the category $`(Q)`$ is more complicated. In particular, there are infinite families of indecomposable objects in some graded dimensions. Second, some of the graded components of $`𝔫^ϵ(Q)`$ are not one-dimensional.
Because of these complications we consider special cases first, and then build the general affine case on them.
### 5.3. Jordan quiver
#### 5.3.1.
In this subsection we consider quiver $`C_1`$ that consists of one vertex connected by an edge with itself: $`C_1=(\{\delta \},\{e\},\mathrm{In},\mathrm{Out})`$, where $`\mathrm{In}(e)=\mathrm{Out}(e)=\delta `$. One can draw $`C_1`$ as follows:
$$\text{}.$$
Though $`C_1`$ is not a Dynkin quiver, the definition of the category of nilpotent representations of $`C_1`$ makes perfect sense. Objects of $`(C_1)`$ are pairs $`(V,x)`$ consisting of a $``$-linear space $`V`$ and a nilpotent $`x\mathrm{Hom}_{}(V,V)`$.
#### 5.3.2.
An indecomposable object of $`(C_1)`$ with dimension $`n`$ is isomorphic to the pair $`(^n,J_n)=𝐉_n`$, where $`J_n`$ is the Jordan block
(5.3.2.1)
$$J_n=\left(\begin{array}{cccccc}0& 1& 0& \mathrm{}& 0& 0\\ 0& 0& 1& \mathrm{}& 0& 0\\ 6\\ 0& 0& 0& \mathrm{}& 1& 0\\ 0& 0& 0& \mathrm{}& 0& 1\\ 0& 0& 0& \mathrm{}& 0& 0\end{array}\right).$$
###### 5.3.3 Proposition.
$`dim_{}\mathrm{Hom}_{(C_1)}(𝐉_n,𝐉_n)=n`$.
###### Proof.
An element $`f`$ of $`\mathrm{Hom}_{(C_1)}(𝐉_n,𝐉_n)`$ is, by definition, an $`n\times n`$ matrix such that $`fJ_n=J_nf`$. It follows that
(5.3.3.1)
$$f=\left(\begin{array}{cccccc}a_1& a_2& a_3& \mathrm{}& a_{n1}& a_n\\ 0& a_1& a_2& \mathrm{}& a_{n2}& a_{n1}\\ 6\\ 0& 0& 0& \mathrm{}& a_2& a_3\\ 0& 0& 0& \mathrm{}& a_1& a_2\\ 0& 0& 0& \mathrm{}& 0& a_1\end{array}\right)$$
for some set $`\{a_i\}_{i=1}^n`$. ∎
#### 5.3.4.
We call the vertex $`\delta `$ of $`C_1`$ the first imaginary root when $`\delta `$ is considered as an element of $`[\delta ]`$.
One has $`e_{C_1}(\alpha ,\beta )=0`$ for any $`\alpha ,\beta [\delta ]`$. We put $`<\alpha ,\beta >_{C_1}=e_{C_1}(\alpha ,\beta )+e_{C_1}(\beta ,\alpha )=0`$ and $`ϵ_{C_1}(\alpha ,\beta )=(1)^{e_{C_1}(\alpha ,\beta )}=1`$.
We use the notation $`r`$ (rotation transformation, cf. 5.5.2) for the identity transformation of $`[\delta ]`$.
### 5.4. Kronecker quiver
#### 5.4.1.
In this section we consider a quiver $`K`$ that consists of two vertices connected by two edges, which are directed in the same way. So $`K=(\{0,1\},\{a,b\},\mathrm{In},\mathrm{Out})`$, where $`\mathrm{In}(a)=\mathrm{In}(b)=1`$, $`\mathrm{Out}(a)=\mathrm{Out}(b)=0`$:
The set of indecomposable representations of $`K`$ was described by Kronecker \[Kro90\].
#### 5.4.2.
It is easy to see that $`\alpha `$ is a positive real root (resp. positive imaginary root) if $`\alpha =(n+1,n)`$ or $`\alpha =(n,n+1)`$ (resp. $`\alpha =(n+1,n+1)`$) for some $`n_+`$. In particular $`\alpha _0=(1,0)`$ and $`\alpha _1=(0,1)`$ are the simple roots, and $`\delta =(1,1)`$ is the first imaginary root.
The Lie bracket (5.2.1.1) in $`𝔫^ϵ(K)`$ can be written as follows:
(5.4.2.1)
$$\begin{array}{cc}\hfill [\stackrel{~}{e}_{(n,n+1)},\stackrel{~}{e}_{(m,m+1)}]& =0,\hfill \\ \hfill [\stackrel{~}{e}_{(n+1,n)},\stackrel{~}{e}_{(m+1,m)}]& =0,\hfill \\ \hfill [\alpha _1(n),\alpha _1(m)]& =0,\hfill \\ \hfill [\alpha _1(n),\stackrel{~}{e}_{(m,m+1)}]& =2(1)^n\stackrel{~}{e}_{(m+n,m+n+1)},\hfill \\ \hfill [\alpha _1(n),\stackrel{~}{e}_{(m+1,m)}]& =2(1)^{n+1}\stackrel{~}{e}_{(m+n+1,m+n)},\hfill \\ \hfill [\stackrel{~}{e}_{(n,n+1)},\stackrel{~}{e}_{(m+1,m)}]& =(1)^{n+m}\alpha _1(m+n+1).\hfill \end{array}$$
Note that $`\alpha _0(n)=\alpha _1(n)`$.
#### 5.4.3.
According to Section 1.4 a representation $`𝐌`$ of the quiver $`K`$ is a quadruple $`((V_0,V_1),(x_a,x_b))`$ consisting of two $``$-linear spaces $`V_0`$ and $`V_1`$ and two linear maps $`x_a,x_b\mathrm{Hom}_{}(V_0,V_1)`$. The dimension of the representation $`𝐌`$ is given by $`dim_{(K)}𝐌=(dim_{}V_0,dim_{}V_1)`$.
#### 5.4.4.
Let us describe the structure of the category $`(K)`$. There are two simple objects $`𝐔_0^0=((,0),(0,0))`$ and $`𝐔_0^1=((0,),(0,0))`$. If one applies repeatedly the reflection functors $`𝒮_0`$ and $`𝒮_1`$ to the simple objects one obtains two families of indecomposable objects $`𝐔_n^0`$ and $`𝐔_n^1`$, given by
$$𝐔_n^0=((^{n+1},^n),(A,B)),dim_{(K)}𝐔_n^0=(n+1,n),$$
$$𝐔_n^1=((^n,^{n+1}),(A^t,B^t)),dim_{(K)}𝐔_n^1=(n,n+1),$$
where
(5.4.4.1)
$$A=\left(\begin{array}{ccccc}1& 0& \mathrm{}& 0& 0\\ 0& 1& \mathrm{}& 0& 0\\ 5\\ 0& 0& \mathrm{}& 0& 0\\ 0& 0& \mathrm{}& 1& 0\end{array}\right),B=\left(\begin{array}{ccccc}0& 1& \mathrm{}& 0& 0\\ 0& 0& \mathrm{}& 0& 0\\ 5\\ 0& 0& \mathrm{}& 1& 0\\ 0& 0& \mathrm{}& 0& 1\end{array}\right).$$
Note that the set of graded dimensions of the indecomposables obtained in this way coincides with the set of positive real roots.
#### 5.4.5.
In this section we use the argument and notation borrowed (with slight adjustments) from \[Lus92, Section 2\].
The category $`(K)`$ is isomorphic to the following category $`𝒞`$: an object of $`𝒞`$ is a triple $`((V_0,V_1),\mathrm{\Delta })`$ consisting of a pair $`(V_0,V_1)`$ of $``$-linear spaces, and a $``$-linear map $`\mathrm{\Delta }:V_0_{}^2V_1`$. A morphism from $`((V_0,V_1),\mathrm{\Delta })`$ to $`((V_0^{},V_1^{}),\mathrm{\Delta }^{})`$ is a pair $`(\varphi _0,\varphi _1)`$ of $``$-linear maps $`\varphi _i:V_iV_i^{}`$ such that $`\varphi _1\mathrm{\Delta }=\mathrm{\Delta }^{}(\varphi _0\mathrm{Id}_^2)`$. Let us introduce the following notation: $`\mathrm{\Delta }_e(v)=\mathrm{\Delta }(ve)`$, where $`e^2`$, $`vV_0`$. The isomorphism of categories $`𝒞`$ and $`(K)`$ is given by a functor $`𝒥`$ that takes $`((V_0,V_1),\mathrm{\Delta })`$ to $`((V_0,V_1),(\mathrm{\Delta }_{\left(\begin{array}{c}1\\ 0\end{array}\right)},\mathrm{\Delta }_{\left(\begin{array}{c}0\\ 1\end{array}\right)}))`$, and acts naturally on morphisms.
Let $`𝐌=((V_0,V_1),\mathrm{\Delta })\mathrm{Ob}(𝒞)`$. We denote by $`\mathrm{Spec}𝐌`$ (*spectrum* of $`𝐌`$) the set of all lines $`z^1`$ such that for $`ez`$, $`e0`$ the map $`\mathrm{\Delta }_e:V_0V_1`$ is not an isomorphism.
Let $`{}_{}{}^{0}𝒞`$ be the following full subcategory of $`𝒞`$. An object $`𝐌`$ of $`𝒞`$ belongs to $`\mathrm{Ob}({}_{}{}^{0}𝒞)`$ if $`\mathrm{Spec}𝐌`$ is a finite set. In particular, if $`𝐌=((V_0,V_1),\mathrm{\Delta })\mathrm{Ob}({}_{}{}^{0}𝒞)`$ then $`dim_{}V_0=dim_{}V_1`$, that is $`dim_{(K)}𝒥(𝐌)R_+^{\mathrm{im}}`$.
One can prove (cf. \[Lus92\]) that $`{}_{}{}^{0}𝒞=_{z^1}{}_{}{}^{0}𝒞_{z}^{}`$ (direct coproduct of abelian categories), where $`{}_{}{}^{0}𝒞_{z}^{}`$ is a full subcategory of $`{}_{}{}^{0}𝒞`$ whose objects are $`𝐌\mathrm{Ob}({}_{}{}^{0}𝒞)`$ such that $`\mathrm{Spec}𝐌=z`$.
The category $`{}_{}{}^{0}𝒞_{z}^{}`$ for any $`z^1`$ is equivalent to the category $`(C_1)`$ of nilpotent representations of the Jordan quiver $`C_1`$. Let us describe the equivalence functor $`_z:(C_1){}_{}{}^{0}𝒞_{z}^{}`$. We choose two elements $`e`$, $`e^{}^2`$ such that $`e0`$, $`e^{}0`$, $`ez`$, $`e^{}z`$. Then the functor $`_z`$ is given by the following action on objects
$$_z((V,x))=((V,V),\mathrm{\Delta }),$$
$$\text{ where }\mathrm{\Delta }(v(ce+c^{}e^{}))=cxv+c^{}v,$$
and the natural action on morphisms. The equivalence class of $`_z`$ does not depend on the choice of $`e`$, $`e^{}`$.
We denote by $`_z`$ the composition functor $`_z=𝒥_z:(C_1)(K)`$. For example, if $`z=(1:0)`$ the functor $`_{(1:0)}`$ is given by the following action on objects
$$_{(1:0)}((V,x))=((V,V),(x,\mathrm{Id}_V))$$
and the natural action on morphisms (here we have chosen $`e=(1,0)`$, $`e^{}=(0,1)`$).
The functor $`_z`$ is a Hall functor for any $`z^1`$.
It follows that $`_z(𝐉_n)`$ is indecomposable in $`(K)`$ of dimension $`n\delta `$, and all indecomposables in $`\mathrm{Ob}(𝒥({}_{}{}^{0}𝒞))`$ are of this form. One can prove (cf. \[Lus92\]) that if $`𝐌\mathrm{Ob}((K))`$ is indecomposable then either $`𝐌`$ is isomorphic to an object, obtained from a simple one by repeated applications of reflection functors, or $`𝐌`$ belongs to $`\mathrm{Ob}(𝒥({}_{}{}^{0}𝒞))`$. Therefore any indecomposable object of $`(K)`$ of dimension $`n\delta `$ belongs to $`\mathrm{Ob}(𝒥({}_{}{}^{0}𝒞))`$, and we can use the map
$$\mathrm{Spec}:𝒯_{n\delta }^1,$$
$$\mathrm{Spec}([𝐌])=\mathrm{Spec}𝐌$$
to identify $`𝒯_{n\delta }`$ and $`^1`$. Abusing notation we just write $`𝒯_{n\delta }=^1`$. In this way $`𝒯_{n\delta }`$ is endowed with a structure of an algebraic variety.
#### 5.4.6.
We summarize the above discussion in the following proposition which describes the set $`𝒯`$ of isomorphism classes of indecomposable objects of $`(K)`$.
###### Proposition.
One has:
1. if $`\alpha R_+`$ then $`𝒯_\alpha =\mathrm{}`$,
2. if $`\alpha R_+^{\mathrm{re}}`$ then $`𝒯_\alpha =pt`$, the only (up to an isomorphism) indecomposable representation with dimension $`\alpha `$ being $`𝐔_n^0`$ (resp. $`𝐔_n^1`$) for $`\alpha =(n+1,n)`$ (resp. $`\alpha =(n,n+1)`$),
3. if $`\alpha R_+^{\mathrm{im}}`$ then $`𝒯_\alpha =^1`$; the indecomposable representation of dimension $`\alpha =(n,n)`$ corresponding to $`z^1`$ is $`_z(𝐉_n)`$.
#### 5.4.7.
Let $`𝐄_\alpha ^{ind}=\{(x_a,x_b)𝐄_\alpha |((^{\alpha _0},^{\alpha _1}),(x_a,x_b))\text{ is indecomposable }\}`$. Abusing notation we denote by $`\mathrm{Spec}:𝐄_{n\delta }^{ind}^1`$ a regular map obtained by assigning to element $`(x_a,x_b)𝐄_{n\delta }^{ind}`$ the point $`z^1`$ such that $`((^n,^n),(x_a,x_b))`$ is isomorphic to $`_z(𝐉_n)`$. The map $`\mathrm{Spec}:𝐄_{n\delta }^{ind}^1`$ is a fibration with a constant fiber.
#### 5.4.8.
Let us introduce the following functions on the set of isomorphism classes of indecomposable objects of $`(K)`$ (which by our usual abuse of notation are also $`𝐆_\alpha `$-equivariant functions on $`𝐄_\alpha `$):
$$E_\alpha =\theta (𝒯_\alpha )=\theta (𝐄_\alpha ^{ind})_\alpha (K),$$
where $`\theta `$ denotes the characteristic function of a set. Note that if $`\alpha R_+^{\mathrm{im}}`$ then $`E_\alpha `$ is the constant function equal to $`1`$ on $`𝒯_\alpha =^1`$.
#### 5.4.9.
The following proposition gives some $``$-products, which we use later. Due to Proposition 2.2.8 we are only interested in the restrictions of the $``$-products to the set of indecomposable representations.
###### Proposition.
(5.4.9.1) $`E_{(0,1)}E_{(n,n1)}|_{𝐄_{(n,n)}^{ind}}`$ $`=E_{(n,n)},`$
(5.4.9.2) $`E_{(n,n1)}E_{(0,1)}|_{𝐄_{(n,n)}^{ind}}`$ $`=0,`$
(5.4.9.3) $`E_{(1,0)}E_{(n1,n)}|_{𝐄_{(n,n)}^{ind}}`$ $`=0,`$
(5.4.9.4) $`E_{(n1,n)}E_{(1,0)}|_{𝐄_{(n,n)}^{ind}}`$ $`=E_{(n,n)},`$
(5.4.9.5) $`E_{(0,1)}E_{(n,n)}|_{𝐄_{(n,n+1)}^{ind}}`$ $`=2E_{(n,n+1)},`$
(5.4.9.6) $`E_{(n,n)}E_{(0,1)}|_{𝐄_{(n,n+1)}^{ind}}`$ $`=0,`$
(5.4.9.7) $`E_{(1,0)}E_{(n,n)}|_{𝐄_{(n+1,n)}^{ind}}`$ $`=0,`$
(5.4.9.8) $`E_{(n,n)}E_{(1,0)}|_{𝐄_{(n+1,n)}^{ind}}`$ $`=2E_{(n+1,n)}.`$
###### Proof.
5.4.9.2 and 5.4.9.6 follow from the fact that $`𝐔_0^1`$ is projective.
5.4.9.3 and 5.4.9.7 follow from the fact that $`𝐔_0^0`$ is injective.
Let us prove 5.4.9.1. We need to find the value of the function $`E_{(0,1)}E_{(n,n1)}`$ on the isomorphism class of the representation $`_z(𝐉_n)`$. We give below the calculation for $`z=(1:0)`$. The case of arbitrary $`z^1`$ is completely analogous, and, moreover, the result does not depend on $`z`$. The value of $`E_{(0,1)}E_{(n,n1)}`$ on the isomorphism class of $`_{(1:0)}(𝐉_n)`$ is equal, by the definition of the $``$-product, to the Euler characteristic of the variety $`N_{𝐔_0^1,𝐔_{n1}^0,_{(1:0)}(𝐉_n)}`$ of all subrepresentations $`𝐕`$ of $`_{(1:0)}(𝐉_n)`$ such that $`𝐕`$ is isomorphic to $`𝐔_0^1`$ and $`_{(1:0)}(𝐉_n)/𝐕`$ is isomorphic to $`𝐔_{n1}^0`$. Since $`_{(1:0)}(𝐉_n)=((^n,^n),(J_n,\mathrm{Id}_^n))`$ it follows that $`N_{𝐔_0^1,𝐔_{n1}^0;_{(1:0)}(𝐉_n)}`$ coincides with the set of all lines $`l^n`$ such that $`((^n,^n/l),(pJ_n,p))`$ is indecomposable, where $`p:^n^n/l`$ is the projection. The object $`((^n,^n/l),(pJ_n,p))`$ is indecomposable if and only if $`l\mathrm{im}J_n`$. Therefore $`N_{𝐔_0^1,𝐔_{n1}^0;_{(1:0)}(𝐉_n)}=^{n1}`$, and we get 5.4.9.1.
Statement 5.4.9.4 is dual to 5.4.9.1.
Let us prove 5.4.9.5. We need to find the value of the function $`E_{(0,1)}E_{(n,n)}`$ on the isomorphism class of the representation $`𝐔_n^1`$. This value is equal, by the definition of the $``$-product, to the Euler characteristic of a variety $`X`$ of all subrepresentations $`𝐕`$ of $`𝐔_n^1`$ such that $`𝐕`$ is isomorphic to $`𝐔_0^1`$ and $`𝐔_n^1/𝐕`$ is isomorphic to $`_z(𝐉_n)`$ for some $`z^1`$. Similar to 5.4.7 one can use the $`\mathrm{Spec}`$ map to prove that the variety $`X`$ is a fibration over $`^1`$ with the fiber over $`z`$ equal to the variety $`N_{𝐔_0^1,_z(𝐉_n);𝐔_n^1}`$ of all subrepresentations $`𝐕`$ of $`𝐔_n^1`$ such that $`𝐕`$ is isomorphic to $`𝐔_0^1`$ and $`𝐔_n^1/𝐕`$ is isomorphic to $`_z(𝐉_n)`$. Let us consider the case $`z=(1:0)`$. Since $`𝐔_n^1=((^n,^{n+1}),(A^t,B^t))`$, where $`A`$ and $`B`$ are as in (5.4.4.1), and $`_{(1:0)}(𝐉_n)=((^n,^n),(J_n,\mathrm{Id}_^n))`$, it follows that $`𝐕=\{(((0,\mathrm{},0),(a,\mathrm{},0)),(0,0))\}_a𝐔_n^1`$. Therefore $`N_{𝐔_0^1,_{(1:0)}(𝐉_n);𝐔_n^1}`$ is a point. The calculation for arbitrary $`z^1`$ is completely analogous, and, moreover, the variety $`N_{𝐔_0^1,_z(𝐉_n);𝐔_n^1}`$ does not depend on $`z`$. Thus $`\chi (X)=\chi (^1)=2`$, which is the statement 5.4.9.5.
Statement 5.4.9.8 is dual to 5.4.9.5. ∎
#### 5.4.10.
Let us recall that $`𝔫^{}(K)`$ denotes the Lie algebra generated by $`E_{(0,1)}`$ and $`E_{(1,0)}`$ with respect to the bracket (2.2.5.1). In 5.2.3 we introduced a homomorphism $`\mathrm{\Xi }_K:𝔫^ϵ(K)𝔫^{}(K)`$. The following proposition is an analog of Ringel Theorem (Theorem Theorem) for the quiver $`K`$.
###### Proposition.
The homomorphism $`\mathrm{\Xi }_K`$ is given by the formulas
$$\begin{array}{cc}\hfill \mathrm{\Xi }_K(\stackrel{~}{e}_\alpha )& =E_\alpha \text{ for }\alpha R_+^{\mathrm{re}}\text{ },\hfill \\ \hfill \mathrm{\Xi }_K(\alpha _1(n))& =(1)^{n+1}E_{(n,n)}.\hfill \end{array}$$
###### Proof.
We prove the proposition by induction. It follows from the definitions that
$$\mathrm{\Xi }_K(\stackrel{~}{e}_{(0,1)})=E_{(0,1)},$$
$$\mathrm{\Xi }_K(\stackrel{~}{e}_{(1,0)})=E_{(1,0)}.$$
Suppose that
$$\mathrm{\Xi }_K(\stackrel{~}{e}_{(k,k+1)})=E_{(k,k+1)},$$
$$\mathrm{\Xi }_K(\stackrel{~}{e}_{(k+1,k)})=E_{(k+1,k)}.$$
Then using (5.4.2.1) and Proposition Proposition we have
$$\mathrm{\Xi }_K(\alpha _1(k+1))=\mathrm{\Xi }_K((1)^k[\stackrel{~}{e}_{(0,1)},\stackrel{~}{e}_{(k+1,k)}])=$$
$$=(1)^k[\mathrm{\Xi }_K(\stackrel{~}{e}_{(0,1)}),\mathrm{\Xi }_K(\stackrel{~}{e}_{(k+1,k)})]=(1)^k[E_{(0,1)},E_{(k+1,k)}]=$$
$$=(1)^{(k+1)+1}E_{(k+1,k+1)},$$
which gives $`\mathrm{\Xi }_K(\alpha _1(k+1))`$, and
$$\mathrm{\Xi }_K(\stackrel{~}{e}_{(k+1,k+2)})=\mathrm{\Xi }_K(\frac{(1)^{k+1}}{2}[\alpha _1(k+1),\stackrel{~}{e}_{(0,1)}])=$$
$$=\frac{(1)^{k+1}}{2}[\mathrm{\Xi }_K(\alpha _1(k+1)),\mathrm{\Xi }_K(\stackrel{~}{e}_{(0,1)})]=\frac{1}{2}[E_{(k+1,k+1)},E_{(0,1)}]=$$
$$=E_{(k+1,k+2)},$$
$$\mathrm{\Xi }_K(\stackrel{~}{e}_{(k+2,k+1)})=\mathrm{\Xi }_K(\frac{(1)^k}{2}[\alpha _1(k+1),\stackrel{~}{e}_{(1,0)}])=$$
$$=\frac{(1)^k}{2}[\mathrm{\Xi }_K(\alpha _1(k+1)),\mathrm{\Xi }_K(\stackrel{~}{e}_{(1,0)})]=\frac{1}{2}[E_{(k+1,k+1)},E_{(1,0)}]=$$
$$=E_{(k+2,k+1)},$$
which provides the induction step.
One can avoid some of these calculations using an argument similar to the proof of Theorem Theorem to get $`\mathrm{\Xi }_K(\stackrel{~}{e}_\alpha )=E_\alpha `$ for $`\alpha R_+^{\mathrm{re}}`$ (note that the indecomposable object with dimension $`\alpha R_+^{\mathrm{re}}`$ can be obtained by repeated applications of reflection functors to a simple object). However, one would still need explicit calculations with the $``$-product to get the statement of the theorem in the case of an imaginary root. ∎
#### 5.4.11.
Let us introduce the following function $`\xi _K:R_+\{\pm 1\}`$.
$$\xi _K(\alpha )=(1)^{(1+dim_{}\mathrm{Hom}_{(K)}(𝐏,𝐏))},$$
where $`𝐏\mathrm{Ob}((K))`$ is indecomposable and $`dim_{(K)}𝐏=\alpha `$.
Since $`𝐔_n^0`$ and $`𝐔_n^1`$ can be obtained by sequences of reflection functors from simple objects we have $`\xi _K(\alpha )=1`$ if $`\alpha R_+^{\mathrm{re}}`$. As any indecomposable object with dimension $`n\delta `$ is isomorphic to $`_z(𝐉_n)`$ for some $`z^1`$, it follows from Proposition 5.3.3 that $`\xi _K(n\delta )=(1)^{n+1}`$. In particular, $`\xi _K`$ is well-defined.
We put $`\stackrel{~}{E}_\alpha =\xi _K(\alpha )E_\alpha `$.
#### 5.4.12.
Here is our final theorem for the Kronecker quiver.
###### Theorem.
The homomorphism $`\mathrm{\Xi }_K:𝔫^ϵ(K)𝔫^{}(K)`$ is given by the formulas
$$\begin{array}{cc}\hfill \mathrm{\Xi }_K(\stackrel{~}{e}_\alpha )& =\stackrel{~}{E}_\alpha \text{ for }\alpha R_+^{\mathrm{re}}\text{ },\hfill \\ \hfill \mathrm{\Xi }_K(\alpha _1(n))& =\stackrel{~}{E}_{(n,n)}.\hfill \end{array}$$
###### Proof.
Follows from Proposition Proposition. ∎
#### 5.4.13.
Given an imaginary root $`n\delta `$, there are infinitely many indecomposable representations with dimension $`n\delta `$ and only one basic element of the Lie algebra with degree $`n\delta `$. In a sense the Lie algebra of functions $`𝔫^{}(K)`$ does not distinguish among these representations (they are ”similar”, though not isomorphic). Note also that the Euler characteristic of the set of indecomposable representations with dimension $`n\delta `$ is equal to $`2`$ (for $`𝒯_{n\delta }=^1`$). It is this $`2`$ that appears in the structure constants 5.4.2.1 !
### 5.5. Cyclic quiver
#### 5.5.1.
Let us fix $`N`$, $`N2`$. In this section we study the cyclic quiver $`C_N=(/N,/N,\mathrm{In},\mathrm{Out})`$, where $`\mathrm{In}(k)k+1(\text{mod}N)`$, $`\mathrm{Out}(k)k(\text{mod}N)`$. It is a quiver of affine type, with the underlying Dynkin graph of type $`A_{N1}^{(1)}`$.
One can draw $`C_7`$ as follows
and $`C_2`$ as follows
#### 5.5.2.
A cyclic quiver has no admissible vertices. In particular, there is no Coxeter element. However there exists the following ”rotation” transformation $`r`$ of the lattice $`[/N]`$.
$$r:[/N][/N],$$
$$r(i)i+1(\text{mod}N).$$
#### 5.5.3.
It is easy to see that if $`(a_0,a_1,\mathrm{},a_{N1})_+[/N]`$ is a positive real root then there exist unique $`i/N`$ and $`l_+`$, $`l0(\text{mod}N)`$, such that $`a_k=\mathrm{\#}\{m\{i,\mathrm{},i+l1\}|mk(\text{mod}N)\}`$. We denote the corresponding root by $`\alpha _{i,l}`$. In particular the set of simple roots is $`\{\alpha _{i,1}\}_{i/N}`$. Imaginary roots are integer multiples of $`\delta =(1,1,\mathrm{},1)[/N]`$.
#### 5.5.4.
Next we want to write down explicitly the Lie bracket (5.2.1.1) in the case of a cyclic quiver. It appears that expressions become much simpler if one changes the basis in $`𝔫^ϵ(C_N)`$ a little. Namely, let
$$\stackrel{~}{f}_{\alpha _{i,l}}=(1)^{[l/N]}\stackrel{~}{e}_{\alpha _{i,l}},$$
where $`[l/N]`$ denotes the integer part of $`l/N`$, and
$$\stackrel{~}{h}(n)=(1)^nh(n)$$
for $`h(n)𝔫_{n\delta }^ϵ`$. Then
(5.5.4.1)
$$[\stackrel{~}{f}_{\alpha _{i,l}},\stackrel{~}{f}_{\alpha _{j,k}}]=\{\begin{array}{cc}\stackrel{~}{f}_{\alpha _{i,k+l}}\hfill & \text{ if }i+lj(\text{mod}N)\text{ }\hfill \\ & \text{ and }k+l0(\text{mod}N)\text{ },\hfill \\ \stackrel{~}{f}_{\alpha _{j,k+l}}\hfill & \text{ if }j+ki(\text{mod}N)\text{ }\hfill \\ & \text{ and }k+l0(\text{mod}N)\text{ },\hfill \\ (\stackrel{~}{\alpha }_{i,1}+\mathrm{}+\stackrel{~}{\alpha }_{j1,1})(n)\hfill & \text{ if }i+lj(\text{mod}N)\text{ }\hfill \\ & \text{ and }k+l=nN\text{ },\hfill \end{array}$$
$$[\stackrel{~}{h}(n),\stackrel{~}{f}_{\alpha _{i,l}}]=<h,\alpha _{i,l}>\stackrel{~}{f}_{\alpha _{i,l+Nn}}.$$
#### 5.5.5.
Let us consider the case $`N=2`$. The Lie bracket 5.5.4.1 becomes
$$[\stackrel{~}{f}_{\alpha _{i,l}},\stackrel{~}{f}_{\alpha _{j,k}}]=\{\begin{array}{cc}\stackrel{~}{\alpha }_{i,1}(n)\hfill & \text{ if }ij(\text{mod}\mathrm{\hspace{0.25em}2})\text{ and }k+l=2n\text{ },\hfill \\ 0\hfill & \text{ otherwise },\hfill \end{array}$$
$$[\stackrel{~}{h}(n),\stackrel{~}{f}_{\alpha _{i,l}}]=<h,\alpha _{i,l}>\stackrel{~}{f}_{\alpha _{i,l+2n}}.$$
The scalar product is given by $`<\alpha _{0,1},\alpha _{0,1}>`$$`=`$ $`<\alpha _{1,1},\alpha _{1,1}>`$$`=`$ $`<\alpha _{0,1},\alpha _{1,1}>`$$`=2`$.
Note that the quivers $`C_2`$ and $`K`$ have the same underlying Dynkin graph $`A_1^{(1)}`$. Therefore the algebras $`𝔫^ϵ(C_2)`$ and $`𝔫^ϵ(K)`$ are (non-canonically) isomorphic. An isomorphism can be given by
$$\eta :𝔫^ϵ(C_2)𝔫^ϵ(K)$$
$$\begin{array}{cc}\hfill \eta (\stackrel{~}{e}_{\alpha _{0,2n+1}})& =\eta ((1)^n\stackrel{~}{f}_{\alpha _{0,2n+1}})=(1)^n\stackrel{~}{e}_{n+1,n},\hfill \\ \hfill \eta (\stackrel{~}{e}_{\alpha _{1,2n+1}})& =\eta ((1)^n\stackrel{~}{f}_{\alpha _{1,2n+1}})=(1)^{n+1}\stackrel{~}{e}_{n,n+1},\hfill \\ \hfill \eta (\alpha _{1,1}(n))& =\eta ((1)^n\stackrel{~}{\alpha }_{1,1}(n))=\alpha _1(n).\hfill \end{array}$$
#### 5.5.6.
We now turn to the Lie algebra $`𝔫^{}(C_N)`$. The category $`(C_N)`$ is isomorphic to the following category $`𝒩_N`$. Objects of $`𝒩_N`$ are pairs $`(V,x)`$, where $`V`$ is a $`/N`$-graded $``$-linear space and $`x`$ is $``$-linear nilpotent endomorphism of $`V`$, such that $`\mathrm{deg}(xv)=\mathrm{deg}(v)+1`$ for any $`vV`$, the set of morphisms from $`(V,x)`$ to $`(W,y)`$ is the set of all $`/N`$-graded $``$-linear maps $`f`$ from $`V`$ to $`W`$, such that $`yf=fx`$.
#### 5.5.7.
A simple object $`𝐏_{i,1}`$ ($`i/N`$) of the category $`𝒩_N`$ is a pair $`𝐏_{i,1}=(V,x)\mathrm{Ob}(𝒩_N)`$, such that
$$V_i=,$$
$$V_j=0\text{ for }ji\text{ },$$
$$x=0.$$
The set $`\{[𝐏_{i,1}]\}_{i/N}`$ is the complete set of isomorphism classes of simple objects of $`𝒩_N`$.
#### 5.5.8.
Given $`i/N`$ and $`l_+`$, $`l0`$, we denote by $`𝐏_{i,l}`$ the following object $`(V,x)`$ of $`𝒩_N`$:
$$V=^l,$$
$$\mathrm{deg}(ϵ_k)lk+i(\text{mod}N),$$
$$x=J_l,$$
where $`\{ϵ_1,\mathrm{},ϵ_l\}`$ is the standard basis in $`^l`$ and $`J_l`$ is the Jordan block (5.3.2.1).
Note that the notation $`𝐏_{i,l}`$ is consistent with the notation $`𝐏_{i,1}`$ used for simple objects.
The dimension of $`𝐏_{i,l}`$ is given by
$$dim_{𝒩_N}𝐏_{i,l}=\{\begin{array}{cc}\alpha _{i,l}\hfill & \text{ if }l0(\text{mod}N)\text{ },\hfill \\ n\delta \hfill & \text{ if }l=nN\text{ }.\hfill \end{array}$$
###### Proposition.
The set $`\{[𝐏_{i,l}]\}_{\begin{array}{c}i/N\\ l_+\backslash \{0\}\end{array}}`$ is the complete set of isomorphism classes of indecomposable objects of $`𝒩_N`$. In other words, $`𝒯_\alpha `$ is a one element set if $`\alpha `$ is a positive real root and an $`N`$ element set if $`\alpha `$ is a positive imaginary root.
###### Proof.
An exercise in (graded) linear algebra (see, for example, \[Lus92, 2.19\]). ∎
#### 5.5.9.
We denote by $`E_{i,l}(C_N)`$ the characteristic function of $`[𝐏_{i,l}]`$.
The following proposition describes the restriction of the $``$-product to the set of indecomposable representations.
###### Proposition.
$$(E_{i,l}E_{j,k})|_{𝐄_{\alpha _{i,l}+\alpha _{j,k}}^{ind}}=\{\begin{array}{cc}E_{j,l+k}\hfill & \text{ if }j+ki(\text{mod}N)\text{ },\hfill \\ 0\hfill & \text{ if }j+ki(\text{mod}N)\text{ }.\hfill \end{array}$$
###### Proof.
Follows from the fact that if $`𝐗\mathrm{Ob}(𝒩_N)`$ is a subobject of $`𝐏_{m,n}`$ then $`𝐗`$ is isomorphic to $`𝐏_{o,p}`$ for some $`pn`$ and $`o=m+np(\text{mod}N)`$, and $`𝐏_{m,n}/𝐗`$ is isomorphic to $`𝐏_{m,np}`$. ∎
#### 5.5.10.
Now we are ready to prove the analog of the Ringel theorem for the quiver $`C_N`$ (cf. \[Rin93\]).
###### Proposition.
The following formulas describe the map $`\mathrm{\Xi }_{C_N}:𝔫^ϵ(C_N)𝔫^{}(C_N)`$:
$$\begin{array}{cc}\hfill \mathrm{\Xi }_{C_N}(\stackrel{~}{f}_{\alpha _{i,l}})& =E_{i,l},\hfill \\ \hfill \mathrm{\Xi }_{C_N}(\stackrel{~}{\alpha }_{i,1}(n))& =E_{i+1,nN}E_{i,nN}.\hfill \end{array}$$
###### Proof.
Let us temporarily denote the map given by the formulas in the formulation of the proposition by $`\mathrm{\Xi }_{C_N}^{}`$. It follows from (5.5.4.1) and Proposition Proposition that $`\mathrm{\Xi }_{C_N}^{}:𝔫^ϵ(C_N)𝔫^{}(C_N)`$ is a homomorphism of Lie algebras, whose value on generators $`\stackrel{~}{f}_{\alpha _{i,1}}`$ of $`𝔫^ϵ(C_N)`$ coincides with the value of $`\mathrm{\Xi }_{C_N}`$. Thus $`\mathrm{\Xi }_{C_N}^{}=\mathrm{\Xi }_{C_N}`$, which proves the proposition.
Another possible proof is an induction on length $`l`$ in $`\alpha _{i,l}`$ similar to the one used in the proof of Proposition Proposition. ∎
#### 5.5.11.
Let us introduce the following function $`\xi _{C_N}:R_+\{\pm 1\}`$ (cf. 5.4.11).
$$\xi _{C_N}(\alpha )=(1)^{(1+dim_{}\mathrm{Hom}_{𝒩_N}(𝐏,𝐏))},$$
where $`𝐏\mathrm{Ob}(𝒩_N)`$ is indecomposable and $`dim_{𝒩_N}𝐏=\alpha `$.
It follows from a graded analogue of Proposition 5.3.3 that
$$\begin{array}{cc}\hfill \xi _{C_N}(\alpha _{i,l})& =(1)^{[l/N]},\hfill \\ \hfill \xi _{C_N}(n\delta )& =(1)^{n+1}.\hfill \end{array}$$
We put $`\stackrel{~}{E}_{i,l}=\xi _{C_N}(\alpha _{i,l})E_{i,l}`$ for $`l0(\text{mod}N)`$ and $`\stackrel{~}{E}_{i,nN}=\xi _{C_N}(n\delta )E_{i,nN}`$.
#### 5.5.12.
Here is our final theorem for a cyclic quiver (we return to the original generators $`\stackrel{~}{e}_\alpha `$ and $`h(n)`$ of $`𝔫^ϵ(C_N)`$).
###### Theorem.
The following formulas describe the map $`\mathrm{\Xi }_{C_N}:𝔫^ϵ(C_N)𝔫^{}(C_N)`$:
$$\begin{array}{cc}\hfill \mathrm{\Xi }_{C_N}(\stackrel{~}{e}_{\alpha _{i,l}})& =\stackrel{~}{E}_{i,l},\hfill \\ \hfill \mathrm{\Xi }_{C_N}(\alpha _{i,1}(n))& =\stackrel{~}{E}_{i,nN}\stackrel{~}{E}_{i+1,nN}.\hfill \end{array}$$
###### Proof.
Follows from Proposition Proposition. ∎
#### 5.5.13. Remark
For any real root $`\alpha `$ there is only one indecomposable representation of $`C_N`$ with dimension equal to $`\alpha `$. However, unlike the case of a finite type quiver the image under the map $`\mathrm{\Xi }_{C_N}`$ of a basic root vector $`\stackrel{~}{e}_\alpha `$ is not always the characteristic function $`E_\alpha `$, but $`\pm E_\alpha `$ (compare with imaginary roots of the Kronecker quiver).
#### 5.5.14.
We recall (see Proposition Proposition) that the set of isomorphism classes of indecomposable representations with dimension $`n\delta `$ is the $`N`$-element set $`𝒯_{n\delta }=\{[𝐏_{i,nN}]\}_{i/N}`$. Therefore the set of functions on $`𝒯_{n\delta }`$ is $`N`$-dimensional. However $`dim_{}𝔫_{n\delta }^ϵ(C_N)=N1`$. The following corollary of Theorem Theorem describes the image of the restriction to $`𝔫_{n\delta }^ϵ(C_N)`$ of the map $`\mathrm{\Xi }_{C_N}`$.
###### Corollary.
$`\mathrm{\Xi }_{C_N}(𝔫_{n\delta }^ϵ(C_N))=\{f:𝒯_{n\delta }|_{i/N}f([𝐏_{i,nN}])=0\}`$
This corollary concludes the discussion of the cyclic quiver $`C_N`$.
### 5.6. General affine quiver
#### 5.6.1.
Finally let $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$ be a quiver of affine type, $`QC_N`$, $`K`$. In particular, $`Q`$ can be a quiver with cyclic (type $`A_k^{(1)}`$) underlying Dynkin graph, but with non-cyclic orientation.
In what follows we use classification of indecomposable objects of $`(Q)`$ as given by V. Dlab and C.M. Ringel \[DR76\]. Let us note that originally the classification was obtained by L.A. Nazarova \[Naz73\], and P. Donovan and M.R. Freislich \[DF73\]. However the description in \[DR76\] is more conceptual and suitable for our purposes. The original proof of Dlab and Ringel relies on case-by-case consideration. Since then there appeared proofs based on McKay correspondence \[Lus92\], and on pure homological algebra (see, for example, \[CB92\]).
Our strategy in studying the map $`\mathrm{\Xi }_Q`$ is to consider different kinds of roots separately.
#### 5.6.2.
Given a root $`\alpha `$ we denote by $`\mathrm{Supp}(\alpha )`$ (support of $`\alpha `$) the set of all vertices $`iI`$ such that $`\alpha _i0`$.
A root $`\alpha `$ is said to be of *finite type* if there exists a proper full subquiver $`Q^{}=(I^{},\mathrm{\Omega }^{},\mathrm{In}^{}\mathrm{Out}^{})`$ of $`Q`$ such that $`\mathrm{Supp}(\alpha )I^{}`$. It follows that $`Q^{}`$ is of finite type and that $`\alpha `$ is a real root (see 1.1.6).
Note that the set $`R_+^{}`$ of all positive roots of $`Q`$ with support in $`Q^{}`$ coincides with the set of positive roots of $`Q^{}`$, and that the restriction of the Euler cocycle $`ϵ_Q`$ to $`[I^{}]\times [I^{}]`$ coincides with $`ϵ_Q^{}`$. Therefore there exists a natural embedding $`𝔦_{Q^{}Q}^ϵ:𝔫^ϵ(Q^{})𝔫^ϵ(Q)`$.
We recall that $`𝔦_{Q^{}Q}:𝔫^{}(Q^{})𝔫^{}(Q)`$ denotes the Hall map associated to the embedding functor $`_{Q^{}Q}:(Q^{})(Q)`$ (see 3.1.6).
The following diagram of Lie algebra homomorphisms is commutative on generators $`\{\stackrel{~}{e}_k\}_{kI^{}}`$ of $`𝔫^ϵ(Q^{})`$, and, therefore, is commutative.
(5.6.2.1)
###### 5.6.3 Proposition.
Let $`\alpha `$ be a positive root of finite type. Then:
1. There exists unique (up to an isomorphism) indecomposable object $`𝐏_\alpha \mathrm{Ob}((Q))`$ with $`dim_{(Q)}𝐏_\alpha =\alpha `$. In other words $`𝒯_\alpha =pt`$. We denote by $`E_\alpha (Q)`$ the characteristic function of $`𝒯_\alpha `$.
2. $`\mathrm{\Xi }_Q(\stackrel{~}{e}_\alpha )=E_\alpha `$.
###### Proof.
The proposition follows from 3.1.4, diagram (5.6.2.1), and Theorem Theorem. ∎
#### 5.6.4.
Let $`:[I]`$ be an additive function given by
$$(\alpha )=e(\delta ,\alpha ),$$
where $`\delta `$ is the first imaginary root and $`e`$ is the Euler form 1.3.4.
The integer $`(\alpha )`$ is called *defect* of $`\alpha `$.
An element $`\alpha [I]`$ is called *regular* (resp. *irregular*) if $`(\alpha )=0`$ (resp. $`(\alpha )0`$).
A root $`\alpha R`$ is called *regular root* (resp. *irregular root*) if it is regular (resp. irregular) as an element of $`[I]`$.
We denote the set of regular (resp. irregular) positive roots by $`R_+^{\mathrm{reg}}`$ (resp. $`R_+^{\mathrm{irr}}`$).
Let us note that each imaginary root is regular, because $`e(\delta ,\delta )=\frac{<\delta ,\delta >}{2}=0`$. However there are regular real roots.
#### 5.6.5.
Irregular roots behave similarly to roots of a quiver of finite type. The following proposition is proven in \[DR76, Sections 1, 2\].
###### Proposition.
Let $`Q`$ be a quiver of affine type, $`QC_N,K`$. Then:
1. A positive root $`\alpha `$ is irregular if and only if there exists an admissible sequence $`\{i_t\}_{t=0}^k`$ of vertices of $`Q`$ such that $`\alpha =\sigma _{i_k}\sigma _{i_{k1}}\mathrm{}\sigma _{i_1}i_0`$, and $`\sigma _{i_l}\sigma _{i_{l1}}\mathrm{}\sigma _{i_1}i_0R_+`$ for any $`l\{1,\mathrm{},k\}`$.
2. If $`𝐌\mathrm{Ob}((Q))`$ is indecomposable and $`dim_{(Q)}𝐌`$ is irregular then $`dim_{(Q)}𝐌`$ is an (irregular) root.
#### 5.6.6.
Having the admissible sequence of vertices 5.6.5.1 one can repeat word-by-word the proof of Theorem Theorem to get the following proposition.
###### Proposition.
Let $`\alpha R_+^{\mathrm{irr}}`$. Then:
1. There exists unique (up to an isomorphism) indecomposable object $`𝐏_\alpha \mathrm{Ob}((Q))`$ with $`dim_{(Q)}𝐏_\alpha =\alpha `$. In other words $`𝒯_\alpha =pt`$. We denote the characteristic function of $`𝒯_\alpha `$ by $`E_\alpha (Q)`$.
2. $`\mathrm{\Xi }_Q(\stackrel{~}{e}_\alpha )=E_\alpha `$.
#### 5.6.7.
Next we turn to regular roots.
Let us consider a full subcategory $`^{\mathrm{reg}}(Q)`$ of $`(Q)`$ with objects being such $`𝐌\mathrm{Ob}((Q))`$ that
* $`(dim_{(Q)}𝐌)=0`$ and
* $`(dim_{(Q)}𝐌^{})0`$ for any subobject $`𝐌^{}`$ of $`𝐌`$ in $`(Q)`$.
Objects of $`^{\mathrm{reg}}(Q)`$ are called *regular* objects (of $`(Q)`$).
The following is proved in \[DR76\].
###### Proposition.
1. The subcategory $`^{\mathrm{reg}}(Q)`$ is abelian and *épaisse* in $`(Q)`$.
2. Let $`𝐌\mathrm{Ob}((Q))`$ be indecomposable and $`(dim_{(Q)}𝐌)=0`$. Then $`𝐌\mathrm{Ob}(^{\mathrm{reg}}(Q))`$.
#### 5.6.8.
The following proposition gives a complete description of the category $`^{\mathrm{reg}}`$ due to V. Dlab and C.M. Ringel \[DR76\] (see also \[Lus92\], \[CB92\]).
###### Proposition.
There exists a set of functors $`\{𝒞_z\}_{z^1}`$, parameterized by points of $`^1`$, such that:
1. $`𝒞_z:(C_{N_z})(Q)`$ is a Hall functor from the category of nilpotent representations of a cyclic or the Jordan quiver to the category of nilpotent representations of $`Q`$.
2. $`^{\mathrm{reg}}(Q)=_{z^1}\mathrm{im}𝒞_z`$ (coproduct of abelian categories).
3. Let $`\alpha R_+^{\mathrm{reg}}`$. Then there exists $`z^1`$ such that $`\alpha =dim𝒞_z(\alpha ^{})`$, where $`\alpha ^{}[/N_z]`$.
4. $`dim𝒞_z(\delta _z)=\delta `$ for any $`z^1`$, where $`\delta _z`$ (resp. $`\delta `$) is the first imaginary root of $`C_{N_z}`$ (resp. $`Q`$).
5. $`dim𝒞_zr_z=cdim𝒞_z`$ for any $`z^1`$, where $`r_z`$ is the rotation transformation of $`[/N_z]`$ and $`c`$ is the Coxeter element of $`Q`$.
6. $`N_z=1`$ for all $`z^1`$ except for a finite number of points (actually no more then $`3`$). We denote the exceptional points by $`z_1,\mathrm{},z_L`$.
7. $`_{z^1}(N_z1)=_{i=1}^L(N_{z_i}1)=|I|2=dim_{}𝔫_{n\delta }^ϵ(Q)1`$.
We refer the reader to Corollary Corollary for a description of the set of integers $`L`$, $`\{N_{z_i}\}_{i=1}^L`$. We do not use the exact values of $`L`$ or $`N_{z_i}`$ in this chapter.
An example of explicit construction of the functors $`𝒞_z`$ in $`E_6^{(1)}`$ case can be found in \[GR92, Chapter 11\].
#### 5.6.9.
Let $`\alpha _{z,i,l}=dim𝒞_z(\alpha _{i,l})`$ and $`𝐏_{z,i,l}=𝒞_z(𝐏_{i,l})`$, where $`z^1`$, $`i/N_z`$, and $`l_+\backslash \{0\}`$ (if $`N_z=1`$ we put $`𝐏_{0,l}=𝐉_l`$; if $`l=nN_z`$ we put $`\alpha _{z,i,l}=n\delta `$). Then $`dim_{(Q)}𝐏_{z,i,l}=\alpha _{z,i,l}`$ and $`𝐏_{z,i,l}`$ is indecomposable because $`𝒞_z`$ is a Hall functor. Moreover it follows from 5.6.5.2, 5.6.6.1, 5.6.7.2, and 5.6.8.2 that $`\{[𝐏_\alpha ]\}_{\alpha R_+^{\mathrm{irr}}}\{[𝐏_{z,i,l}]\}_{\begin{array}{c}z^1\\ i/N_z\\ l_+\backslash \{0\}\end{array}}`$ is the complete set of isomorphism classes of indecomposable representations.
#### 5.6.10.
The set of functors $`𝒞_z`$ satisfying conditions of Proposition Proposition is defined uniquely up to a permutation of the set of points of $`^1`$ and up to automorphisms of categories $`(C_{N_z})`$. More explicitly, one can first identify the set of isomorphism classes of simple objects of $`^{\mathrm{reg}}(Q)`$. This set splits into orbits of the Coxeter functor (product of reflection functors, corresponding to the reflections in the Coxeter element). Each orbit is the set $`\{[𝐏_{z,i,1}]\}_{i/N_z}`$ for some $`z`$. Then $`\mathrm{Ob}(\mathrm{im}𝒞_z)`$ is the set of such $`𝐌\mathrm{Ob}(^{\mathrm{reg}}(Q))`$ that Jordan-Hölder series (in $`^{\mathrm{reg}}(Q)`$) of $`𝐌`$ contains only objects $`𝐏_{z,i,1}`$.
#### 5.6.11.
The following proposition lists some properties of the dimension maps $`dim𝒞_z`$.
###### Proposition.
One has
1. $`ϵ_Q(dim𝒞_z(\alpha ),dim𝒞_w(\beta ))=1`$ if $`zw`$.
2. $`ϵ_Q(dim𝒞_z(\alpha ),dim𝒞_z(\beta ))=ϵ_{C_{N_z}}(\alpha ,\beta )`$.
3. $`<dim𝒞_z(\alpha ),dim𝒞_w(\beta )>_Q=0`$ if $`zw`$.
4. $`<dim𝒞_z(\alpha ),dim𝒞_z(\beta )>_Q=<\alpha ,\beta >_{C_{N_z}}`$.
5. $`\{\alpha _{z,i,l}\}_{\begin{array}{c}z^1\\ i/N_z\\ l_+\backslash \{0\}\end{array}}`$ is the complete list of regular roots, the only equalities being $`\alpha _{z,i,nN_z}=\alpha _{w,j,nN_w}(=n\delta )`$.
6. Let
$$V=[\{\alpha _{z,i,1}mod\delta \}_{\begin{array}{c}z^1\\ i/N_z\end{array}}]=$$
$$=[\{\alpha _{z_k,i,1}mod\delta \}_{\begin{array}{c}k=1,\mathrm{},L\\ i/N_{z_k}\end{array}}][I]/\delta ,$$
where $`[S]`$ denotes the $``$-linear span of a set $`S`$. Then
$$dim_{}V=|I|2=dim_{}𝔫_{n\delta }^ϵ(Q)1.$$
The following is a generating set of linear relations among elements of the set $`\{\alpha _{z_k,i,1}mod\delta \}_{\begin{array}{c}k=1,\mathrm{},L\\ i/N_{z_k}\end{array}}`$.
$$\{\underset{i/N_{z_k}}{}\alpha _{z_k,i,1}0(mod\delta )\}_{k\{1,\mathrm{},L\}}$$
7. Let $`z=z_k`$ (i.e. $`N_z1`$). Then there exists $`i_0/N_z`$ such that $`\alpha _{z,i_0,1}`$ is of finite type.
###### Proof.
5.6.11.1 and 5.6.11.3 follow from 1.4.6 and 5.6.8.2.
5.6.11.2 and 5.6.11.4 follow from 5.6.8.1 and 3.1.5.
5.6.11.5 follows from the definition of root, 5.6.8.2, 5.6.8.3, 5.6.11.3, and 5.6.11.4.
5.6.11.6 follows from 5.6.8.7, 5.6.11.3, and 5.6.11.4.
Let us prove 5.6.11.7. Let $`pI`$ be an extending vertex (i.e. $`\delta _p=1`$). It follows from the following equality
$$\underset{i/N_z}{}\alpha _{z,i,1}=dim𝒞_z(\underset{i/N_z}{}\alpha _{i,1})=dim𝒞_z(\delta _z)=\delta $$
that either $`N_z=1`$ or there exists $`i_0/N_z`$ such that $`(\alpha _{z,i_0,1})_p=0`$, which proves 5.6.11.7. ∎
#### 5.6.12.
Now we are ready to give the classification of indecomposable objects of $`(Q)`$ (cf. \[DR76\]).
###### Theorem.
Let $`Q`$ be a quiver of affine type, $`QC_N,K`$. Let $`𝒯`$ be the set of isomorphism classes of indecomposable objects of $`(Q)`$. Then
1. $`𝒯=\{[𝐏_\alpha ]\}_{\alpha R_+^{\mathrm{irr}}}\{[𝐏_{z,i,l}]\}_{\begin{array}{c}z^1\\ i/N_z\\ l_+\backslash \{0\}\end{array}}`$ ,
2. $`𝒯_\alpha =\mathrm{}`$ if $`\alpha R_+`$ ,
3. $`𝒯_\alpha =pt`$ if $`\alpha R_+^{\mathrm{re}}`$ ,
4. $`𝒯_{k\delta }=\{[𝐏_{z,i,kN_z}]\}_{\begin{array}{c}z^1\\ i/N_z\end{array}}`$ .
###### Proof.
The theorem follows from 5.6.9 and 5.6.11.5. ∎
#### 5.6.13.
Theorem Theorem allows us to introduce the following function $`\xi _Q:R_+\{\pm 1\}`$ (compare with 5.4.11, 5.5.11).
$$\xi _Q(\alpha )=(1)^{(1+dim_{}\mathrm{Hom}_{(Q)}(𝐏,𝐏))},$$
where $`𝐏\mathrm{Ob}((Q))`$ is indecomposable and $`dim_{(Q)}𝐏=\alpha `$.
The function $`\xi _Q(\alpha )`$ is well-defined for $`\alpha R_+^{\mathrm{re}}`$ because there is unique (up to an isomorphism) indecomposable object with dimension $`\alpha `$. In the case of an imaginary root one has
(5.6.13.1)
$$\begin{array}{c}\xi _Q(n\delta )=(1)^{(1+dim_{}\mathrm{Hom}_{(Q)}(𝐏_{z,i,nN_z},𝐏_{z,i,nN_z}))}=\hfill \\ \hfill =(1)^{(1+dim_{}\mathrm{Hom}_{(C_{N_z})}(𝐏_{i,nN_z},𝐏_{i,nN_z}))}=(1)^{n+1}.\end{array}$$
This expression does not depend on $`z`$ or $`i`$.
If $`\alpha `$ is an irregular root then $`\xi _Q(\alpha )=1`$, because $`𝐏_\alpha `$ can be obtained from a simple object by a sequence of reflection functors. We put $`\stackrel{~}{E}_\alpha =\xi _Q(\alpha )E_\alpha =E_\alpha `$ for $`\alpha R_+^{\mathrm{irr}}`$.
#### 5.6.14.
We return to the map $`\mathrm{\Xi }_Q`$.
Let $`E_{z,j,l}(Q)`$ be the characteristic function of $`[𝐏_{z,j,l}]𝒯_{\alpha _{z,j,l}}`$ and $`\stackrel{~}{E}_{z,j,l}=\xi _Q(\alpha _{z,j,l})E_{z,j,l}`$.
###### Proposition.
Let $`1iL`$ and $`j/N_{z_i}`$. Then
$$\mathrm{\Xi }_Q(\stackrel{~}{e}_{\alpha _{z_i,j,1}})=E_{z_i,j,1}=\stackrel{~}{E}_{z_i,j,1}.$$
###### Proof.
Given $`i`$ the statement is true for one $`j`$, say $`j_0`$, because of 5.6.11.7 and Proposition 5.6.3. Any other $`\alpha _{z_i,j,1}`$ can be obtained from $`\alpha _{z_i,j_0,1}`$ by repeated applications of the Coxeter element (see 5.6.8.5). Then one can use the sequence of reflection functors corresponding to the sequence of reflections in the Coxeter element, and reason similarly to the proof of Theorem Theorem to get the first equality in the statement of the proposition for any $`j`$.
The second equality in the statement of the proposition follows from the fact that $`\xi _Q(\alpha _{z_i,j,1})=\xi _{C_{N_{z_i}}}(\alpha _{j,1})=1`$. ∎
#### 5.6.15.
Next we consider arbitrary regular positive real roots.
###### Proposition.
The map $`𝔠_i^ϵ:𝔫^ϵ(C_{N_{z_i}})𝔫^ϵ(Q)`$ given by
$$\begin{array}{cc}\hfill 𝔠_i^ϵ(\stackrel{~}{e}_{\alpha _{j,l}})& =\stackrel{~}{e}_{\alpha _{z_i,j,l}},\hfill \\ \hfill 𝔠_i^ϵ(\alpha _{j,1}(n))& =\alpha _{z_i,j,1}(n)\hfill \end{array}$$
is an injective homomorphism of Lie algebras for any $`i\{1,\mathrm{},L\}`$.
###### Proof.
The proposition follows from 5.6.11.2, 5.6.11.4, and 5.6.8.4.
#### 5.6.16.
We denote by $`𝔠_i:𝔫^{}(C_{N_{z_i}})𝔫^{}(Q)`$ the Hall map associated to the functor $`𝒞_{z_i}:(C_{N_{z_i}})(Q)`$. In particular, $`𝔠_i(\stackrel{~}{E}_{j,l})=\stackrel{~}{E}_{z_i,j,l}`$.
The following diagram of Lie algebra homomorphisms
is commutative on generators $`\stackrel{~}{e}_{\alpha _{j,1}}`$ of $`𝔫^ϵ(C_{N_{z_i}})`$ (because of Proposition Proposition), and, therefore, is commutative, and using Theorem Theorem we get the following:
###### Proposition.
Let $`1iL`$. Then
$$\begin{array}{c}\hfill \mathrm{\Xi }_Q(\stackrel{~}{e}_{\alpha _{z_i,j,l}})=\stackrel{~}{E}_{z_i,j,l}\text{ for }l0(\text{mod}N_{z_i})\text{ },\\ \hfill \mathrm{\Xi }_Q(\alpha _{z_i,j,1}(n))=\stackrel{~}{E}_{z_i,j,nN_{z_i}}\stackrel{~}{E}_{z_i,j+1,nN_{z_i}}.\end{array}$$
#### 5.6.17.
At this point we know values of $`\mathrm{\Xi }_Q`$ on $`𝔫_\alpha ^ϵ(Q)`$ for any $`\alpha R_+^{\mathrm{re}}`$. As for an imaginary root space $`𝔫_{n\delta }^ϵ(Q)`$ we only know values of $`\mathrm{\Xi }_Q`$ on the $``$-linear span of $`\{\alpha _{z_k,i,1}(n)\}_{\begin{array}{c}k=1,\mathrm{},L\\ i/N_{z_k}\end{array}}`$, which has codimension $`1`$ in $`𝔫_{n\delta }^ϵ(Q)`$ due to Proposition 5.6.11.6.
Note also that all functions in the part of the image of $`\mathrm{\Xi }_Q`$ that we know by now vanish on $`[𝐏_{z,0,n}]`$ for any $`zz_1,\mathrm{},z_L`$ (i.e. when $`N_z=1`$).
#### 5.6.18.
The following proposition is due to V. Dlab and C.M. Ringel \[DR76\].
###### Proposition.
There exist
* a functor $`𝒦:(K)(Q)`$ from the category of representations of the Kronecker quiver to the category of representations of $`Q`$,
* a set of elements $`k_i/N_{z_i}`$, one for each $`i\{1,\mathrm{},L\}`$
such that
1. $`𝒦`$ is a Hall functor.
2. $`𝒞_z=𝒦_z`$ for any $`zz_1,\mathrm{},z_L`$.
3. $`𝒦_{z_i}(𝐉_n)=𝐏_{z_i,k_i,nN_{z_i}}`$ for any $`i\{1,\mathrm{},L\}`$,
4. Let $`\alpha _0=dim_{(Q)}𝒦(𝐔_0^0)`$. Then $`\alpha _0R_+^{\mathrm{irr}}`$ and $`dim_{(Q)}𝒦(𝐔_0^1)=(\delta \alpha _0)R_+^{\mathrm{irr}}`$.
Let us remark that by redefining functors $`𝒞_{z_i}`$ (”rotating cyclic quivers”) one can make $`k_i=0`$ for all $`i`$.
#### 5.6.19.
As the authors could not find a reference for Proposition Proposition except for \[DR76\] which relies on case-by-case considerations, let us give a sketch of a construction of the functor $`𝒦`$.
We choose an extending vertex $`pI`$, which is admissible. If the underlying Dynkin graph of $`Q`$ is not $`A_n^{(1)}`$ then any extending vertex has only one adjacent edge, and, therefore, is admissible. If the underlying Dynkin graph of $`Q`$ is $`A_n^{(1)}`$ then one can still choose an admissible extending vertex because $`Q`$ has non-cyclic orientation. We assume that $`p`$ is a sink. If $`p`$ is a source the construction is analogues or one can compose the functor $`𝒦`$ described below with the reflection functor $`𝒮_p`$.
We fix an indecomposable object $`(^{\delta p},y)\mathrm{Ob}((Q))`$ (note that $`\delta p`$ is a positive root of finite type). Let $`q`$ ($`q_1`$, $`q_2`$ in the $`A_n^{(1)}`$ case) be the vertex (vertices) connected with $`p`$ by an edge (edges). We denote by $`h_0`$ (resp. $`h_1`$, $`h_2`$) the edge (resp. edges) such that $`\mathrm{In}(h_0)=p`$ (resp. $`\mathrm{In}(h_1)=\mathrm{In}(h_2)=p`$) and $`\mathrm{Out}(h_0)=q`$ (resp. $`\mathrm{Out}(h_1)=q_1`$, $`\mathrm{Out}(h_2)=q_2`$).
Instead of $`𝒦`$ we construct a functor $`𝒦^{}=𝒦𝒥:𝒞(Q)`$ (see 5.4.5 for the definitions of $`𝒞`$ and $`𝒥`$). Then $`𝒦=𝒦^{}𝒥^1`$ (being an isomorphism of categories $`𝒥`$ has the inverse).
Let $`Q`$ have the underlying graph not of $`A_n^{(1)}`$ type. Then it is known that $`(\delta p)_q=\delta _q=2`$. The functor $`𝒦^{}:𝒞(Q)`$ is given by the following action on objects:
$$𝒦^{}((V_0,V_1),\mathrm{\Delta })=(W,z),$$
where
$$W_i=V_0_{}^{(\delta p)_i}\text{ if }ip\text{ },$$
$$W_p=V_1,$$
$$z_h=\mathrm{Id}_{V_0}y_h:V_0_{}^{(\delta p)_{\mathrm{Out}(h)}}V_0_{}^{(\delta p)_{\mathrm{In}(h)}}\text{ if }hh_0\text{ },$$
$$z_{h_0}=\mathrm{\Delta }:V_0_{}^2,V_1,$$
and the natural action on morphisms.
Let $`Q`$ have the underlying graph of $`A_n^{(1)}`$ type. Then $`(\delta p)_{q_1}=(\delta p)_{q_2}=1`$. The functor $`𝒦^{}:𝒞(Q)`$ is given by the following action on objects:
$$𝒦^{}((V_0,V_1),\mathrm{\Delta })=(W,z),$$
where
$$W_i=V_0_{}^{(\delta p)_i}\text{ if }ip\text{ },$$
$$W_p=V_1,$$
$$z_h=\mathrm{Id}_{V_0}y_h:V_0_{}^{(\delta p)_{\mathrm{Out}(h)}}V_0_{}^{(\delta p)_{\mathrm{In}(h)}}\text{ if }hh_1,h_2\text{ },$$
$$z_{h_1}z_{h_2}=\mathrm{\Delta }:V_0_{}()V_1,$$
and the natural action on morphisms.
One can check that the functor $`𝒦`$ defined above satisfies all the conditions in Proposition Proposition (see \[CB92, §9\] for a similar argument).
#### 5.6.20.
Functor $`𝒦:(K)(Q)`$ satisfying conditions of Proposition Proposition is not unique (cf. 5.6.10). For example, in construction 5.6.19 one needs to choose an extending vertex, and, moreover, one can compose the functor $`𝒦`$ described in 5.6.19 with a sequence of reflection functors.
On the other hand a choice of functor $`𝒦`$ fixes the parameterization of the functors $`𝒞_z`$ by points in $`^1`$. More precisely, given a functor $`𝒦:(K)(Q)`$, the set of functors $`\{𝒞_z:(C_{N_z})(Q)\}_{z^1}`$, such that Propositions Proposition and Proposition hold true is uniquely defined up to automorphisms of categories $`(C_{N_z})`$.
From now on we fix some particular choice of the functors $`𝒦`$ and $`𝒞_z`$.
#### 5.6.21.
Since $`𝒦`$ is a Hall functor, the following proposition follows from 3.1.5.
###### Proposition.
A map $`𝔨^ϵ:𝔫^ϵ(K)𝔫^ϵ(Q)`$ given by
$$\begin{array}{cc}\hfill 𝔨^ϵ(\stackrel{~}{e}_{(n+1,n)})& =\stackrel{~}{e}_{\alpha _0+n\delta },\hfill \\ \hfill 𝔨^ϵ(\stackrel{~}{e}_{(n,n+1)})& =\stackrel{~}{e}_{\alpha _0+(n+1)\delta },\hfill \\ \hfill 𝔨^ϵ(\alpha _0^{}(n))& =\alpha _0(n)\hfill \end{array}$$
is a Lie algebra homomorphism. Here $`\alpha _0^{}=dim_{(K)}𝐔_0^0=(1,0)`$ and $`\alpha _0=dim_{(Q)}𝒦(𝐔_0^0)`$.
#### 5.6.22.
Let $`E_0(n)`$ be the characteristic function of the set $`\{[𝒦_z(𝐉_n)]\}_{z^1}𝒯_{n\delta }`$. In particular, according to 5.6.18.3,
$$E_0(n)([𝐏_{z_i,j,nN_{z_i}}])=\{\begin{array}{cc}1\hfill & \text{ if }j=k_i\text{ },\hfill \\ 0\hfill & \text{ if }jk_i\text{ }.\hfill \end{array}$$
We put $`\stackrel{~}{E}_0(n)=\xi _Q(n\delta )E_0(n)=(1)^{n+1}E_0(n)`$.
#### 5.6.23.
Let $`𝔨:𝔫^{}(K)𝔫^{}(Q)`$ be the Hall map associated with the functor $`𝒦`$.
Using 5.6.18.4 and Proposition we get the following equalities
$$\begin{array}{cc}\hfill \mathrm{\Xi }_Q(\stackrel{~}{e}_{\alpha _0})& =E_{\alpha _0}=\stackrel{~}{E}_{\alpha _0},\hfill \\ \hfill \mathrm{\Xi }_Q(\stackrel{~}{e}_{\delta \alpha _0})& =E_{\delta \alpha _0}=\stackrel{~}{E}_{\delta \alpha _0},\hfill \end{array}$$
which guarantee that the following diagram of Lie algebra homomorphisms is commutative.
Using this diagram and Theorem Theorem we get the following proposition.
###### Proposition.
$`\mathrm{\Xi }_Q(\alpha _0(n))=\stackrel{~}{E}_0(n)`$.
#### 5.6.24.
The last proposition concludes description of the map $`\mathrm{\Xi }_Q`$ due to the following
###### Proposition.
The set
$$\{\alpha _{z_i,j,1}(n)\}_{\begin{array}{c}1iL\\ j/N_{z_i}\end{array}}\{\alpha _0(n)\}𝔫_{n\delta }^ϵ(Q)=[I]/\delta $$
spans the whole space $`𝔫_{n\delta }^ϵ(Q)=[I]/\delta `$ as a $``$-linear space.
###### Proof.
Note that $`\alpha _0`$ is linearly independent with $`\{\alpha _{z_i,j,1}\}_{\begin{array}{c}1iL\\ j/N_{z_i}\end{array}}\{\delta \}`$ because $`(\alpha _0)0`$ whereas $`(\alpha _{z_i,j,1})=\delta =0`$. Now the proposition follows from 5.6.11.6. ∎
#### 5.6.25.
Let us simplify our notation a little. We put $`N_i=N_{z_i}`$, $`\stackrel{~}{E}_{i,j}(n)=\stackrel{~}{E}_{z_i,j,nN_i}_{n\delta }(Q)`$, and $`\alpha _{i,j}=\alpha _{z_i,j,1}[I][I]`$.
#### 5.6.26.
Here is our main theorem for a general affine quiver, which follows from Propositions Proposition, Proposition, 5.6.11.5, Proposition, Proposition.
###### Theorem.
Let $`Q`$ be a quiver of affine type, $`QC_N,K`$. Then the following equalities completely describe the map $`\mathrm{\Xi }_Q`$:
$$\begin{array}{cc}\hfill \mathrm{\Xi }_Q(\stackrel{~}{e}_\alpha )& =\stackrel{~}{E}_\alpha \text{ for any }\alpha R_+^{\mathrm{re}}\text{ },\hfill \\ \hfill \mathrm{\Xi }_Q(\alpha _{i,j}(n))& =\stackrel{~}{E}_{i,j}(n)\stackrel{~}{E}_{i,j+1}(n),\hfill \\ \hfill \mathrm{\Xi }_Q(\alpha _0(n))& =\stackrel{~}{E}_0(n).\hfill \end{array}$$
#### 5.6.27.
Theorem Theorem is an affine analog of the Ringel theorem. Let us state a few corollaries.
###### Corollary.
The map $`\mathrm{\Xi }_Q:𝔫^ϵ(Q)𝔫^{}(Q)`$ is an isomorphism of Lie algebras.
###### Proof.
The map $`\mathrm{\Xi }_Q`$ is surjective by construction. Theorem Theorem implies that $`\mathrm{Ker}\mathrm{\Xi }_Q=0`$. ∎
Since $`\mathrm{\Xi }_Q`$ is an isomorphism given by explicit formulas (see Theorem Theorem) one can use $`\mathrm{\Xi }_Q`$ together with (5.2.1.1) to obtain explicit expression for a Lie bracket of any two elements of $`𝔫^{}(Q)`$ (cf. Corollary Corollary).
#### 5.6.28.
The next corollary of Theorem Theorem is a kind of inverse to Proposition 2.2.8.4.
###### Corollary.
The support of the space of functions $`𝔫^{}(Q)`$ is equal to the set of all isomorphism classes of indecomposable objects of $`(Q)`$. In other words, for any indecomposable object $`𝐌`$ there exists an element $`f`$ of the Lie algebra $`𝔫^{}(Q)`$ such that $`f([𝐌])0`$.
#### 5.6.29.
The most interesting feature of the category $`(Q)`$ for a general affine quiver $`Q`$ is the structure of the set $`𝒯_{n\delta }`$ of indecomposable objects with dimension equal to an imaginary root (it does not depend on the root in question) and the restriction of the image of the map $`\mathrm{\Xi }_Q`$ to this set. We refer the reader to Theorem and Theorem for the corresponding results.
Let us consider a map
$$\begin{array}{c}\hfill \mu :𝒯_{n\delta }^1,\\ \hfill \mu ([𝐏_{z,i,nN_z}])=z.\end{array}$$
Note that $`\mu ^1(z)`$ consists of $`N_z`$ points. In particular, the map $`\mu `$ is injective on $`^1\backslash \{z_1,\mathrm{},z_L\}`$.
We denote by $`\mu _{}`$ the following map from the space of $``$-valued functions on $`𝒯_{n\delta }`$ to the space of $``$-valued functions on $`^1`$:
$$\mu _{}(\varphi )(z)=\underset{[𝐏]\mu ^1(z)}{}\varphi ([𝐏]).$$
#### 5.6.30.
The following is a corollary of Theorem Theorem.
###### Corollary.
$$𝔫_{n\delta }^{}(Q)=\mathrm{im}(\mathrm{\Xi }_Q|_{𝔫_{n\delta }^ϵ(Q)})=\{\varphi :𝒯_{n\delta }|\mu _{}(\varphi )\text{ is a constant function}\}.$$
#### 5.6.31.
Instead of introducing the functions $`\stackrel{~}{E}_\alpha `$ we could change the definition of the Lie algebra $`𝔫^ϵ(Q)`$. Namely let $`𝔫^ϵ^{}(Q)`$ be the Lie algebra defined in the same way as $`𝔫^ϵ(Q)`$ (see 5.2.1) but using the following cocycle
(5.6.31.1)
$$ϵ_Q^{}(\alpha ,\beta )=ϵ_Q(\alpha ,\beta )\xi _Q(\alpha +\beta )(\xi _Q(\alpha ))^1(\xi _Q(\beta ))^1$$
instead of the Euler cocycle $`ϵ_Q`$. Here we use arbitrary extension of $`\xi _Q`$ from $`R_+`$ to $`[I]`$. To avoid confusion we denote the generators of $`𝔫^ϵ^{}(Q)`$ by $`\stackrel{~}{f}_\alpha `$ and $`\stackrel{~}{h}(n)`$ instead of $`\stackrel{~}{e}_\alpha `$ and $`h(n)`$. Then it follows from Theorem Theorem that the map $`\mathrm{\Xi }_Q^{}:𝔫^ϵ^{}(Q)𝔫^{}(Q)`$ given by
$$\begin{array}{cc}\hfill \mathrm{\Xi }_Q^{}(\stackrel{~}{f}_\alpha )& =E_\alpha \text{ for any }\alpha R_+^{\mathrm{re}}\text{ },\hfill \\ \hfill \mathrm{\Xi }_Q^{}(\stackrel{~}{\alpha }_{i,j}(n))& =E_{i,j}(n)E_{i,j+1}(n),\hfill \\ \hfill \mathrm{\Xi }_Q^{}(\stackrel{~}{\alpha }_0(n))& =E_0(n)\hfill \end{array}$$
is an isomorphism of Lie algebras.
Let us remark that equation (5.6.31.1) means that the cocycles $`ϵ_Q`$ and $`ϵ_Q^{}`$ differ by a coboundary, that is belong to the same cohomology class in $`H^2([I],/2)`$.
### 5.7. Lie algebras over $``$
#### 5.7.1.
Throughout this section $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$ is a quiver of affine type, $`QC_N,K`$.
#### 5.7.2.
Let $`\alpha R_+^{\mathrm{re}}`$. We put $`\mathrm{}_\alpha ^ϵ(Q)=\stackrel{~}{e}_\alpha 𝔫_\alpha ^ϵ(Q)`$. We denote by $`\mathrm{}_{n\delta }^ϵ(Q)𝔫_{n\delta }^ϵ(Q)`$ the lattice additively generated by $`\{imod\delta \}_{iI}`$. The lattice $`\mathrm{}_{n\delta }^ϵ(Q)`$ is isomorphic to the root lattice $`[I^{}]`$ of the quiver $`Q^{}=(I^{},\mathrm{\Omega }^{},\mathrm{In}^{}\mathrm{Out}^{})`$ obtained by removing from $`Q`$ an extending vertex and the adjacent edges.
We put $`\mathrm{}^ϵ(Q)=_{\alpha R_+}\mathrm{}_\alpha ^ϵ(Q)`$. The lattice $`\mathrm{}^ϵ(Q)`$ is a Lie algebra over $``$ with respect to the Lie bracket (5.2.1.1).
###### Proposition.
The Lie algebra $`\mathrm{}^ϵ(Q)`$ is generated by $`\{\stackrel{~}{e}_i\}_{iI}`$.
###### Proof.
Induction on $`n(\alpha )`$, where $`n(\alpha )=_{iI}\alpha _i`$ for a root $`\alpha `$. It is crucial that $`Q`$ is simply laced. ∎
#### 5.7.3.
Let $`\alpha R_+^{\mathrm{re}}`$. We denote by $`\mathrm{}_\alpha ^{}(Q)`$ the set of all $``$-valued functions on $`𝒯_\alpha =pt`$. We put (cf. Corollary 5.6.30)
$$\mathrm{}_{n\delta }^{}(Q)=\{\varphi :𝒯_{n\delta }|\mu _{}(\varphi )\text{ is a constant function}\}.$$
Let $`\mathrm{}^{}(Q)=_{\alpha R_+}\mathrm{}_\alpha ^{}(Q)`$. Our goal is to prove that $`\mathrm{\Xi }_Q(\mathrm{}^ϵ(Q))=\mathrm{}^{}(Q)`$.
#### 5.7.4.
Let $`\widehat{\mathrm{}}^{}(Q)=_{\alpha R_+}\widehat{\mathrm{}}_\alpha ^{}(Q)`$, where $`\widehat{\mathrm{}}_\alpha ^{}(Q)`$ is the set of all constructible $`𝐆_\alpha `$-invariant functions $`f:𝐄_\alpha ^{nil}`$ such that $`f(x)=0`$ if $`(^\alpha ,x)`$ is decomposable.
It follows from the definition of the $``$-product and Proposition 2.2.8 that $`\widehat{\mathrm{}}^{}(Q)`$ is a Lie algebra over $``$ with respect to the bracket (2.2.5.1).
Using definition of the map $`\mathrm{\Xi }_Q`$ and Proposition Proposition one gets $`\mathrm{\Xi }_Q(\mathrm{}^ϵ(Q))\widehat{\mathrm{}}^{}(Q)`$. On the other hand it follows from definitions that $`\mathrm{\Xi }_Q(\mathrm{}^ϵ(Q))𝔫^{}(Q)`$ and from Corollary 5.6.30 that $`\mathrm{}^{}(Q)=𝔫^{}(Q)\widehat{\mathrm{}}^{}(Q)`$. Thus we get an inclusion
(5.7.4.1)
$$\mathrm{\Xi }_Q(\mathrm{}^ϵ(Q))\mathrm{}^{}(Q).$$
It follows from Theorem Theorem that
$$\stackrel{~}{E}_{i,j}(n)\stackrel{~}{E}_{i,j+1}(n)\mathrm{\Xi }_Q(\mathrm{}^ϵ(Q)),$$
$$\stackrel{~}{E}_0(n)\mathrm{\Xi }_Q(\mathrm{}^ϵ(Q)),$$
$$\stackrel{~}{E}_\alpha \mathrm{\Xi }_Q(\mathrm{}^ϵ(Q))\text{ for }\alpha R_+^{\mathrm{re}}\text{ }.$$
Since $`\{\stackrel{~}{E}_{i,j}(n)\stackrel{~}{E}_{i,j+1}(n)\}_{\begin{array}{c}1iL\\ j/N_i\end{array}}\{E_0(n)\}`$ is a set of additive generators of $`\mathrm{}_{n\delta }^{}(Q)`$ and $`\stackrel{~}{E}_\alpha `$ is an additive generator of $`\mathrm{}_\alpha ^{}(Q)`$ for $`\alpha R_+^{\mathrm{re}}`$, we conclude that
(5.7.4.2)
$$\mathrm{}^{}(Q)\mathrm{\Xi }_Q(\mathrm{}^ϵ(Q)).$$
Combining (5.7.4.1) and (5.7.4.2) we get the following
###### Proposition.
$`\mathrm{\Xi }_Q(\mathrm{}^ϵ(Q))=\mathrm{}^{}(Q)`$.
It follows, in particular, that $`\mathrm{}^{}(Q)`$ is closed with respect to the Lie bracket (2.2.5.1) and is generated as a Lie algebra over $``$ by $`\{E_i\}_{iI}`$.
## 6. Fine structure of affine root systems
In this section we deliberate on relations of our results to the structure theory of affine Lie algebras.
#### 6.1.1.
Let $`Q=(I,\mathrm{\Omega },\mathrm{In},\mathrm{Out})`$ be a quiver of affine type, $`QC_N,K`$.
The set $`\{\alpha _{i,j}[I]\}_{\begin{array}{c}1iL\\ j/N_i\end{array}}`$ (notation as in 5.6.25) plays a prominent role in our consideration of an imaginary root space $`𝔫_{n\delta }^{}(Q)`$. This set of elements of the root lattice $`[I]`$ of $`Q`$ is canonically associated to the quiver. We call $`\alpha _{i,j}`$ a *cyclic root*.
Tables of cyclic roots for particular quivers of affine type can be found in \[DR76, Section 6\].
Let us give another description of the set of cyclic roots. It follows from 5.6.8.5 and 5.6.11.5 that regular roots have finite orbits under the action of the Coxeter element $`c`$. The converse is also true, that is any root having finite $`c`$-orbit is regular (see \[DR76, Section 1\]).
We call a finite $`c`$-orbit lowest if an element of this orbit cannot be represented as a sum of two regular positive roots (this condition does not depend on the choice of the element of the orbit). Let $`𝒪_i=\{\alpha _{i,j}\}_{j/N_i}`$. It follows from 5.6.8.5, 5.6.11.5 and 5.6.11.6 that $`\{𝒪_i\}_{1iL}`$ is the complete set of the lowest $`c`$-orbits.
Let $`R_i[I]`$ be the additive span of elements of an orbit $`𝒪_i`$. It follows from Proposition Proposition that $`<R_i,R_j>_Q=0`$ if $`ij`$, $`_{i=1}^LR_i_{}`$ has codimension $`1`$ in $`[I]`$, and each $`R_i`$ is isomorphic to the root lattice of type $`A_{N_i1}^{(1)}`$. In other words there are $`L`$ root lattices of type $`A_n^{(1)}`$ ”hidden” inside the root lattice of $`Q`$, and to uncover them one can use finite orbits of the Coxeter element.
#### 6.1.2.
It follows from Proposition Proposition that the set $`\{\alpha _{i,j}mod\delta =\mathrm{\Xi }_Q^1(\stackrel{~}{E}_{i,j}(n)\stackrel{~}{E}_{i,j+1}(n))\}_{\begin{array}{c}1iL\\ j/N_i\end{array}}\{\alpha _0mod\delta =\mathrm{\Xi }_Q^1(\stackrel{~}{E}_0(n))\}`$ is a set of additive generators of $`\mathrm{}_{n\delta }^ϵ(Q)`$. Since $`\delta _p=1`$ for an extending vertex $`pI`$ it follows that if $`\alpha `$, $`\beta [I][I]`$ and $`\alpha \beta (\text{mod}\delta )`$ then $`\alpha \beta (\text{mod}\delta )`$. Combining with 5.6.11.6, 5.6.8.4 and Proposition Proposition we get the following
###### Proposition.
The abelian group $`[I]`$ has the following presentation:
$$<\{\alpha _{i,j}\}_{\begin{array}{c}1iL\\ j/N_i\end{array}}\{\alpha _0\}\{\delta \}|\{\underset{j/N_i}{}\alpha _{i,j}=\delta \}_{1iL}>,$$
that is $`\{\alpha _{i,j}\}_{\begin{array}{c}1iL\\ j/N_i\end{array}}\{\alpha _0\}\{\delta \}`$ is a set of additive generators of $`[I]`$ and the set $`\{_{j/N_i}\alpha _{i,j}=\delta \}_{1iL}`$ is a generating set of relations among the generators.
#### 6.1.3.
From now till the end of the chapter we assume that $`\delta \alpha _0`$ is a simple root, $`\delta \alpha _0=pI`$. We also assume that $`p`$ is an extending vertex. Then it follows from the definition of $`\alpha _0`$ that $`p`$ is a sink.
Let $`Q^{}=(I^{},\mathrm{\Omega }^{},\mathrm{In}^{},\mathrm{Out}^{})`$ be the quiver obtained from $`Q`$ by removing the vertex $`p`$ and adjacent edges. It is a quiver of finite type. The lattice $`\mathrm{}_{n\delta }^ϵ(Q)`$ is isomorphic to $`[I^{}]`$.
Assume $`Q`$ is not of $`A_n^{(1)}`$ type. It follows from 5.6.8.4 that $`(_{j/N_i}\alpha _{i,j})_p=\delta _p=1`$. Hence for each $`i\{1,\mathrm{},L\}`$ there exists $`n_i/N_i`$ such that
$$(\alpha _{i,n_i})_p=1,$$
$$(\alpha _{i,l})_p=0\text{ if }ln_i\text{ }.$$
Note that $`p`$ being a sink we may assume that the reflection $`\sigma _p`$ is the first reflection in the Coxeter element $`c`$. Therefore
$$0=(\alpha _{i,n_i+1})_p=(c\alpha _{i,n_i})_p=(\alpha _{i,n_i})_q(\alpha _{i,n_i})_p=(\alpha _{i,n_i})_q1,$$
where $`q`$ is the vertex connected with $`p`$ by an edge. We conclude that $`(\alpha _{i,n_i})_q=1`$. Note that $`(_{j/N_i}\alpha _{i,j})_q=\delta _q=2`$. Hence there exists $`m_i/N_i`$, $`m_in_i`$ such that
$$(\alpha _{i,m_i})_q=(\alpha _{i,n_i})_q=1$$
$$(\alpha _{i,l})_q=0\text{ if }ln_i\text{}m_i\text{ }.$$
Now we again use the fact that $`\sigma _p`$ is the first reflection in the Coxeter element $`c`$:
$$(\alpha _{i,m_i+1})_p=(c\alpha _{i,m_i})_p=(\alpha _{i,m_i})_q(\alpha _{i,m_i})_p=1.$$
Therefore $`n_i=m_i+1`$.
It follows from the above considerations and from Proposition Proposition that the abelian group $`[I^{}]`$ is freely generated by the set $`\{\alpha _{i,j}\}_{\begin{array}{c}1iL\\ j/N_i,jn_i\end{array}}\{\alpha _0\}`$ and that
(6.1.3.1)
$$\begin{array}{cc}\hfill e_Q(\alpha _0,\alpha _{i,m_i})=e_Q(p,\alpha _{i,m_i})& =0,\hfill \\ \hfill e_Q(\alpha _{i,m_i},\alpha _0)=e_Q(\alpha _{i,m_i},p)& =1,\hfill \\ \hfill e_Q(\alpha _0,\alpha _{i,j})=e_Q(p,\alpha _{i,j})& =0\text{ if }jn_i\text{}m_i\text{ },\hfill \\ \hfill e_Q(\alpha _{i,j},\alpha _0)=e_Q(\alpha _{i,j},p)& =0\text{ if }jn_i\text{}m_i\text{ }.\hfill \end{array}$$
In the $`A_n^{(1)}`$ case one also has a set of pairs $`\{(n_i,m_i)\}_{i=1}^L`$ such that the above statement is true. The proof is similar, except that one should consider two vertices $`q_1`$ and $`q_2`$ connected to $`p`$ instead of one vertex $`q`$. From now on we do not distinguish between $`A_n^{(1)}`$ and non-$`A_n^{(1)}`$ cases.
Let $`\widehat{Q}=(\widehat{I},\widehat{\mathrm{\Omega }},\widehat{\mathrm{In}},\widehat{\mathrm{Out}})`$ be a quiver given as follows. The set of vertices $`\widehat{I}`$ is equal to $`\{\widehat{\alpha }_{i,j}\}_{\begin{array}{c}1iL\\ j/N_i,jn_i\end{array}}\{\mathrm{}\}`$, the set of edges is equal to the set of all ordered pairs $`(h,t)`$, where either $`h=\widehat{\alpha }_{i,j}`$ and $`t=\widehat{\alpha }_{i,j1}`$, or $`h=\mathrm{}`$ and $`t=\widehat{\alpha }_{i,m_i}`$, and the maps $`\widehat{\mathrm{In}}`$ and $`\widehat{\mathrm{Out}}`$ are given by $`\widehat{\mathrm{In}}((h,t))=h`$, $`\widehat{\mathrm{Out}}((h,t))=t`$. One can draw the quiver $`\widehat{Q}`$ as follows.
We summarize the results of this subsection in the following proposition.
###### Proposition.
The map $`\nu :[\widehat{I}][I^{}]`$ given by $`\nu (\widehat{\alpha }_{i,j})=\alpha _{i,j}`$, $`\nu (\mathrm{})=\alpha _0`$ is an isomorphism of lattices and $`e_Q^{}(\nu (\alpha ),\nu (\beta ))=e_{\widehat{Q}}(\alpha ,\beta )`$ for any $`\alpha `$, $`\beta [\widehat{I}]`$.
#### 6.1.4.
It follows from Proposition Proposition that $`<\nu (\alpha ),\nu (\beta )>_Q^{}=<\alpha ,\beta >_{\widehat{Q}}`$ for any $`\alpha `$, $`\beta [\widehat{I}]`$. In other words, $`\nu `$ is an isometry of lattices.
We recall (see \[Bou68\]) that the set of simple roots in a root system of finite type is unique up to an action of the semi-direct product of the automorphism group of the Dynkin graph and the Weyl group. In our situation we state it as follows (abusing notation we denote by the same letter $`\pi `$ a morphism of Dynkin graphs $`\pi :(I^{},\mathrm{\Omega }^{},\{\mathrm{In}^{},\mathrm{Out}^{}\})(\widehat{I},\widehat{\mathrm{\Omega }},\{\widehat{\mathrm{In}},\widehat{\mathrm{Out}}\})`$ and the induced morphism of the root lattices $`\pi :[I^{}][\widehat{I}]`$).
###### Proposition.
There exist a morphism of Dynkin graphs $`\pi :(I^{},\mathrm{\Omega }^{},\{\mathrm{In}^{},\mathrm{Out}^{}\})(\widehat{I},\widehat{\mathrm{\Omega }},\{\widehat{\mathrm{In}},\widehat{\mathrm{Out}}\})`$ and an element $`w`$ of the Weyl group of $`(\widehat{I},\widehat{\mathrm{\Omega }},\{\widehat{\mathrm{In}},\widehat{\mathrm{Out}}\})`$ such that $`\pi \nu =w`$, or, more explicitly,
$$\pi (\alpha _{i,j})=w(\widehat{\alpha }_{i,j}),$$
for any $`1iL`$, $`j/N_i,jn_i`$ and
$$\pi (\alpha _0)=w(\mathrm{}).$$
A similar result for such an orientation of $`Q`$ that each vertex is admissible has been proven by R. Steinberg \[Ste85\].
In particular, the underlying Dynkin graphs of $`\widehat{Q}`$ and $`Q^{}`$ coincide, and we get the following corollary which describes the set of numbers $`L`$, $`\{N_i\}_{i=1}^L`$ for a quiver with a given underlying Dynkin graph (cf. \[DR76\]).
###### Corollary.
One has:
1. If $`Q`$ is of $`A_n^{(1)}`$ type then either $`L=1`$ and $`N_1=n`$, or $`L=2`$ and $`N_1+N_2=n+1`$.
2. If $`Q`$ is of $`D_n^{(1)}`$ type then $`L=3`$, $`N_1=N_2=2`$, and $`N_3=n2`$ (up to a permutation of $`\{N_i\}`$).
3. If $`Q`$ is of $`E_n^{(1)}`$ type $`(6n8)`$ then $`L=3`$, $`N_1=2`$, $`N_2=3`$, and $`N_3=n3`$ (up to a permutation of $`\{N_i\}`$).
Corollary Corollary does not give the numbers $`L`$, $`\{N_i\}_{i=1}^L`$ in the $`A_n^{(1)}`$ case. Let us describe these numbers for the sake of completeness (see \[DR76\] for the proof). We denote by $`J_1`$ (resp $`J_2`$) the number of clockwise (resp. counterclockwise) oriented arrows in $`Q`$. Note that $`J_1,J_20`$ since $`Q`$ is not cyclic. If $`J_1=1`$ (resp. $`J_2=1`$), then $`L=1`$ and $`N_1=J_2`$ (resp. $`N_1=J_1`$). If $`J_11`$ and $`J_21`$, then $`L=2`$ and $`N_i=J_i`$ up to the transposition of $`N_1`$ and $`N_2`$.
#### 6.1.5.
Finally we would like to discuss relation of the quiver construction of the Lie algebra $`𝔫`$ and the representation theory of finite groups. In this section we assume that $`Q`$ is a quiver of $`A_{2n+1}^{(1)}`$, $`D_n^{(1)}`$, or $`E_n^{(1)}`$ type with such an orientation that each vertex is admissible. Lusztig \[Lus92\] has reinterpreted representations of such quivers using McKay correspondence \[McK80\].
Let us recall that McKay correspondence establishes an isomorphism of lattices $`\eta :K\mathrm{\Gamma }[I]`$, where $`K\mathrm{\Gamma }`$ is the Grothendieck group of the category of $``$-linear finite dimensional representations of a finite subgroup $`\mathrm{\Gamma }`$ of $`SL(2,)`$ and $`[I]`$ is the root lattice of an affine Dynkin graph $`(I,E,\mathrm{Ends})`$. In particular, $`\eta ^1(i)`$ is the class of a simple $`\mathrm{\Gamma }`$-module, which we denote by $`\rho _i`$. The following formula describes the pull-back of the bilinear form $`<,>`$ under $`\eta ^{}`$.
$$<\eta (X),\eta (Y)>=dim_{}\mathrm{Hom}_\mathrm{\Gamma }(X^2,Y)dim_{}\mathrm{Hom}_\mathrm{\Gamma }(X\rho ,Y),$$
where $`X,Y`$ are $`\mathrm{\Gamma }`$-modules, $`^2`$ is equipped with the trivial $`\mathrm{\Gamma }`$-action, and $`\rho `$ is the natural $`2`$-dimensional representation of $`\mathrm{\Gamma }SL(2,)`$.
The isomorphism of lattices yields the following bijection between the set of affine Dynkin graphs and the set of finite subgroups of $`SL(2,)`$.
$`A_n^{(1)}/(n+1)`$ \- cyclic group of order $`n+1`$,
$`D_n^{(1)}𝔻_{n2}`$ \- binary dihedral group,
$`E_6^{(1)}𝕋`$ \- binary tetrahedral group,
$`E_7^{(1)}𝕆`$ \- binary octahedral group,
$`E_8^{(1)}𝕀`$ \- binary icosahedral group.
Note that affine Dynkin graphs of type $`A_{2n}^{(1)}`$ do not have an orientation such that each vertex is admissible and are excluded from our considerations. The rest of the graphs admit precisely two such orientations.
For an affine quiver $`Q`$ with the special orientation Lusztig was able to re-obtain the classification of indecomposable representations of $`Q`$ entirely in terms of the representation theory of $`\mathrm{\Gamma }`$. Below we discuss relation of maximal cyclic subgroups of $`\mathrm{\Gamma }`$ and cyclic roots of $`Q`$ following \[Lus92\].
Let $`F`$ be the set of lines $`l`$ in $`P(\rho )`$ whose isotropy group $`\mathrm{\Gamma }_l`$ (which is a cyclic group) has order greater than $`2`$. The set $`F`$ is finite. Let us choose representatives $`\{l_i\}_{i=1}^L^{}`$ for $`\mathrm{\Gamma }`$-orbits in $`F`$. We assume that if $`\mathrm{\Gamma }_l^{}=\mathrm{\Gamma }_{l_i}`$ for some $`i`$, but $`l^{}`$ and $`l_i`$ are not in the same $`\mathrm{\Gamma }`$-orbit then $`l^{}=l_j`$ for some $`j`$. Let $`\mathrm{\Gamma }_i=\mathrm{\Gamma }_{l_i}`$. It follows from our assumptions on the quiver $`Q`$ that the order of $`\mathrm{\Gamma }_i`$ is even. We put $`N_i^{}=\frac{|\mathrm{\Gamma }_i|}{2}`$. Note that $`l_i`$ is a one-dimensional $`\mathrm{\Gamma }_i`$-module, and that $`l_i^{(2N_i^{})}`$ is a trivial $`\mathrm{\Gamma }_i`$-module. Let $`\varkappa _{i,j}=l_i^{(2j)}l_i^{(2j+1)}`$ for $`j/N_i^{}`$.
One can identify the number of exceptional points of $`^1`$ and their multiplicities appearing in the classification of indecomposable representations of $`Q`$ in terms of the group $`\mathrm{\Gamma }`$ and its maximal cyclic subgroups (see \[Lus92, Section 1\]). We recall that $`\{\alpha _{i,j}[I]\}_{\begin{array}{c}1iL\\ j/N_i\end{array}}`$ is the set of cyclic roots.
###### Proposition.
One has
1. $`L=L^{}`$.
2. $`N_i=N_i^{}`$ up to a permutation of indexes $`i\{1,\mathrm{},L\}`$.
3. $`\alpha _{i,j}=\eta (\mathrm{Ind}_{\mathrm{\Gamma }_i}^\mathrm{\Gamma }\varkappa _{i,j})`$ up to a permutation of indexes $`i\{1,\mathrm{},L\}`$ and a cyclic permutation of indexes $`j/N_i`$.
In particular, combining Propositions Corollary and Proposition one can identify the lengths of the branches of the Dynkin graph with the orders of the maximal cyclic subgroups of $`\mathrm{\Gamma }`$.
Using Frobenius Reciprocity we get the following formula for the cyclic roots in terms of the representation theory of $`\mathrm{\Gamma }`$ and its cyclic subgroups.
$$(\alpha _{i,j})_k=dim_{}\mathrm{Hom}_\mathrm{\Gamma }(\mathrm{Ind}_{\mathrm{\Gamma }_i}^\mathrm{\Gamma }\varkappa _{i,j},\rho _k)=dim_{}\mathrm{Hom}_{\mathrm{\Gamma }_i}(\varkappa _{i,j},\rho _k).$$
We would like to remark that both our construction of the Lie algebra $`𝔫`$ and the proof of Theorem Theorem can be reformulated in the language of the representation theory of $`\mathrm{\Gamma }`$. In particular, instead of the reflection functors which involve all possible quiver orientations one might use “the square root” of the Coxeter functor and its inverse employed by Lusztig \[Lus92\].
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# Gravity-driven Dense Granular Flows
\[
## Abstract
We report and analyze the results of numerical studies of dense granular flows in two and three dimensions, using both linear damped springs and Hertzian force laws between particles. Chute flow generically produces a constant density profile that satisfies scaling relations suggestive of a Bagnold grain inertia regime. The type of force law has little impact on the behavior of the system. Bulk and surface flows differ in their failure criteria and flow rheology, as evidenced by the change in principal stress directions near the surface. Surface-only flows are not observed in this geometry.
\]
Understanding the behavior of granular materials has been a great challenge to scientists and engineers. One major hurdle has been the lack of a formal connection between the complicated but relatively well-understood world of contact mechanics, which describes the nature and dynamics of intergranular interactions, and empirical continuum models that describe the macroscopic behavior of the system.
We aim to isolate the essential features of granular flow, unencumbered by complicated boundary effects. To this end, we perform simulations of granular dynamics in a simple geometry: gravity driven dense flow down an inclined plane, denoted henceforth as “chute flow”. Several remarkable features emerge:
(1) In steady-state, the packing fraction $`\varphi `$ remains constant as a function of depth, beyond a dilatant surface region a few layers thick (Fig. 1.) The compacting influence of increasing stress due to the weight of grains overhead is balanced by increasing velocity fluctuations towards the bottom of the assembly.
(2) Unlike Couette flows, the entire assembly is in motion and surface-only flows are not observed.
(3) Components of the stress tensor and the square of the strain rate grow linearly with depth, indicative of Bagnold grain-inertia behavior.
(4) Normal stresses differ from each other in a systematic way (Fig. 2), which we do not fully understand.
We report results of large scale molecular dynamics simulations of chute flow in two and three dimensions (2D and 3D), with interparticle interactions betweeen the (monodisperse) spheres modeled using both damped linear springs and Mindlin-Hertz contact forces, with static friction. Detailed results of the simulations will be presented elsewhere. The main obstacle for experiments and simulations so far has been the difficulty of reaching and maintaining steady state. Previous simulations employed very few particles or did not reach steady state. All of these simulations were in 2D with the exception of Walton. Experiments on chute flow did not involve deep assemblies. Different effects of flow were also studied in simulation, such as size segregation.
The 3D simulation cells contain spheres of diameter $`d`$ and mass $`m`$, supported by a fixed bottom on the $`xy`$ plane. The bottom wall is constructed from a cross-section of a random close packing of identical spheres, providing a rough surface. Periodic boundary conditions are imposed along $`x`$ and $`y`$ directions. 2D simulations follow the same procedure, except that particles are restricted to the $`xz`$ plane and the bottom consists of a regular array of particles of diameter $`2d`$. In both cases, there is no slip observed at the bottom. (There is some slip in 2D with a regular array of particles of diameter $`d`$ at the bottom.) The gravity vector g is rotated by the tilt angle $`\theta `$ away from the $`(\widehat{z})`$ direction in the $`xz`$ plane, so that the free surface is normal to the $`\widehat{z}`$ axis. In 3D, most of our results are for $`\mathrm{8\hspace{0.17em}000}`$ particle systems with a simulation cell of size $`L_x=18.6d`$ and $`L_y=9.3d`$, resulting in an assembly roughly 40 particles deep at rest. In 2D, $`L_x=100d`$ and the number of particles varied from a few hundred to $`\mathrm{20\hspace{0.17em}000}`$. In this Letter, we present results for $`N=\mathrm{10\hspace{0.17em}000}`$, with a depth of about 100 particles.
We use the contact force model of Cundall and Strack. Static friction is implemented by keeping track of the elastic shear displacement throughout the lifetime of a contact. For two contacting particles at positions r<sub>1</sub> and r<sub>2</sub>, with velocities v<sub>1,2</sub> and angular velocities $`𝝎_{1,2}`$, the force on particle 1 is computed as follows: The normal compression $`\delta `$, normal velocity $`𝐯_n`$, relative surface velocity $`𝐯_t`$, and the rate of change of the elastic tangential displacement $`𝐮_t`$, set to 0 at the initiation of a contact, are given by
$`\delta `$ $`=`$ $`d𝐫_{12},`$ (1)
$`𝐯_n`$ $`=`$ $`(𝐯_{12}\widehat{r}_{12})\widehat{r}_{12},`$ (2)
$`𝐯_t`$ $`=`$ $`𝐯_{12}𝐯_n(𝝎_1+𝝎_2)\times 𝐫_{12}/2,`$ (3)
$`{\displaystyle \frac{d𝐮_t}{dt}}`$ $`=`$ $`𝐯_t{\displaystyle \frac{(𝐮_t𝐯_{12})𝐫_{12}}{𝐫_{12}^2}},`$ (4)
where $`𝐫_{12}=𝐫_1𝐫_2`$, $`\widehat{r}_{12}=𝐫_{12}/𝐫_{12}`$, and $`𝐯_{12}=𝐯_1𝐯_2`$. The second term in Eq.(4) arises from the rigid body rotation around the contact point and assures that $`𝐮_t`$ always remains in the local tangent plane of contact. Normal and tangential forces acting on particle 1 are given by
$`𝐅_n`$ $`=`$ $`k_nf(\delta /d)\left(\delta \widehat{r}_{12}\tau _n𝐯_n\right),`$ (5)
$`𝐅_t`$ $`=`$ $`k_sf(\delta /d)\left(𝐮_t\tau _s𝐯_t\right),`$ (6)
where $`k_{n,s}`$ and $`\tau _{n,s}`$ are elastic constants and viscoelastic relaxation times respectively; $`f(x)=1`$ for damped linear springs or $`f(x)=\sqrt{x}`$ for Hertzian contacts between spheres. A local Coulomb yield criterion, $`F_t<\mu F_n`$, is satisfied by truncating the magnitude of $`𝐮_t`$ as necessary. Thus, the contact surfaces are treated as “stuck” while $`F_t<\mu F_n`$, and slipping while the yield criterion is satisfied. This “proportional loading” approximation is a simplification of the much more complicated and hysteretic behavior of real contacts. The force on particle 2 is determined from Newton’s third law. Each particle is also subject to a body force
$$𝐅_{\mathrm{body}}=mg(\widehat{z}\mathrm{cos}\theta +\widehat{x}\mathrm{sin}\theta ).$$
(7)
All results are given in terms of non-dimensionalized quantities: Distances, times, velocities, forces, elastic constants and stresses are reported in units of $`d`$, $`t_0\sqrt{d/g}`$, $`v_0\sqrt{gd}`$, $`F_0mg`$, $`k_0mg/d`$ and $`\sigma _0mg/d^2`$, respectively. The summary of parameters used in the simulations are shown in Table I. In these units, the correct elastic constant for glass spheres with $`d=100\mu m`$ would have been $`k_n^{\mathrm{glass}}/k_03\times 10^{10}`$, which would have prohibited a large-scale simulation. We use $`k_n/k_0=2\times 10^5`$, while controlling the coefficient of restitution for binary collisions through $`\tau _n`$, assuming that we remain sufficiently close to the $`k_n\mathrm{}`$ limit of small deformations. Simulations for $`k_n/k_0=2\times 10^4`$ and $`2\times 10^6`$ in 2D gave essentially the same results, supporting this assumption. For Hertzian contacts, the ratio $`k_s/k_n`$ depends on the Poisson ratio of the material, and is about $`2/3`$ for most materials. We use the value $`k_s/k_n=2/7`$, since this makes the period of normal and shear contact oscillations equal to each other in the damped linear springs case. The collisional dynamics are not very sensitive to the precise value of this ratio.
The equations of motion for the translational and rotational degrees of freedom were integrated with either a third-order Gear predictor-corrector or velocity-Verlet scheme with a time step $`\delta t=1\times 10^4`$. Typically it was necessary to run between 5 and $`20\times 10^6\delta t`$ to reach steady state, particularly when starting from a non-flowing state. All coarse grained quantities have been averaged both temporally (typically 2 to $`8\times 10^6\delta t)`$ and spatially over slices of constant $`z`$.
The main characteristics of all the flows are: (i) The existence of an “angle of repose” $`\theta _r`$, such that granular flows can not be sustained for $`\theta <\theta _r`$, (ii) a steady-state flow with a packing fraction independent of depth for $`\theta _r<\theta <\theta _{\mathrm{max}}`$, and (iii) for $`\theta >\theta _{\mathrm{max}}`$, development of a shear thinning layer at the bottom of the assembly that results in lift-off and unstable acceleration of the entire assembly. For very thin assemblies, less than about $`20`$ layers, the value of $`\theta _r`$ depends on their depth, in agreement with experiment. Here we consider only deep assemblies where the value of $`\theta _r`$ is independent of depth, and focus our attention on region (ii). As seen in Fig. 1, the packing fraction $`\varphi `$ remains constant as a function of depth, away from the free surface and the bottom wall. Its value is shown as a function of $`\theta `$ in Fig. 2. Results in 2D for systems of size $`\mathrm{5\hspace{0.17em}000}`$ and $`\mathrm{20\hspace{0.17em}000}`$ demonstrate that the thicknesses of the boundary layers at the bottom and top are independent of the height of the assembly. The data suggest an approximate tilt dependence for $`\varphi `$ of the form
$`\varphi _{2D}(\theta )`$ $``$ $`\varphi _{2D}^{\mathrm{max}}c_{2D}(\theta \theta _{r,2D})^2,`$ (8)
$`\varphi _{3D}(\theta )`$ $``$ $`\varphi _{3D}^{\mathrm{max}}c_{3D}(\theta \theta _{r,3D}),`$ (9)
where $`\varphi _{2D}^{\mathrm{max}}=0.810(5)`$ and $`\varphi _{3D}^{\mathrm{max}}=0.595(5)`$.
In 2D, upon lowering the tilt angle below $`\theta _r`$, we observe a compaction to a polycrystalline triangular lattice with $`\varphi _{2D}0.9`$. This causes considerable hysteresis in 2D simulations as $`\theta `$ is subsequently increased beyond $`\theta _r`$: There is no flow until a maximum angle of stability is exceeded. Initial failure always occurs at the bottom, followed by movement of a dilation front towards the top. Once the system reaches steady state, $`\theta `$ can be reduced and the system continues to flow while $`\theta >\theta _r`$. On the other hand, in 3D there is no jump in $`\varphi `$ as the system comes to a stop, and no detectable hysteresis as the system is stopped and restarted by varying $`\theta `$. Thus, Mohr-Coulomb analysis of the stress tensor can be used to relate the flow rheology near $`\theta _r`$ to the Coulomb failure criterion as follows: The Mohr circle is the set of normal and shear stresses ($`\sigma `$ and $`\tau `$, see inset in Fig. 3) associated with all possible shear planes. The points $`A`$ and $`B`$ at coordinates $`(\sigma _{zz},\sigma _{xz})`$ and $`(\sigma _{xx},\sigma _{xz})`$ in the $`(\sigma ,\tau )`$ plane form a diameter of the circle. At a given tilt angle, $`\sigma _{zz}`$ and $`\sigma _{xz}`$ are determined by force balance, which pins down the location of point $`A`$ ($`\sigma _{xz}/\sigma _{zz}=\mathrm{tan}\theta `$; see Fig. 3, inset.)
However, $`\sigma _{xx}`$, thus the location of point $`B`$, remains indeterminate from these considerations. The “gap angle” $`\delta \theta \theta _{MC}\theta `$ is a measure of the difference between $`\tau /\sigma `$ along the slip plane parallel to the surface ($`=\mathrm{tan}\theta `$) and the largest value of $`\tau /\sigma `$ experienced by the system (equal to the slope $`\mathrm{tan}\theta _{MC}`$ of the line that passes through the origin and that is tangent to the Mohr circle.) For an ideal Coulomb material with a uniform yield criterion, $`\delta \theta =0`$ at $`\theta =\theta _r`$, when the static system is at incipient yield. The behavior of $`\delta \theta `$ as a function of depth at different tilt angles, as shown in Figure 3, reveal that: (i) $`\delta \theta `$ becomes independent of depth below a transitional surface layer about $`5d`$ to $`8d`$ in thickness, (ii) the gap angle remains finite at the bottom and in the bulk but vanishes at the top surface when $`\theta =\theta _r`$. It thus seems that although the bulk of the system has the intrinsic capability to withstand slightly larger tilt angles, the destabilization of the surface at $`\theta =\theta _r`$ is ultimately responsible for the failure and initiation of flow in the entire system. Note that the transitional surface layer is not directly related to the dilatant layer; it is much thicker near $`\theta =\theta _r`$ and penetrates well into the region of constant density. Interestingly, in 2D the gap angle at $`\theta _r`$ is actually larger at the surface compared to the bulk, and both gap angles remain finite. However, the presence of hysteresis precludes us from applying the Mohr-Coulomb analysis in 2D.
A feature that distinguishes granular flows from Newtonian fluid flows is that normal stresses, i.e., diagonal terms in the stress tensor, are in general not equal to each other. Although $`\sigma _{xx}\sigma _{zz}`$ in our simulations, we observe small but systematic deviations, which are depicted in Fig. 2 by plotting the normal stress anomaly in the bulk, $`(\sigma _{xx}\sigma _{zz})/\sigma _{zz}`$, as a function of tilt angle. These deviations are likely to be due to a constitutive equation of the flow rheology which we have not yet been able to determine. In addition, in all 3D runs, $`\sigma _{yy}`$ is smaller than the other normal stresses by $`1015\%`$, suggesting that consolidation and compaction normal to the shear plane is poorer.
Another question of particular interest is the relationship between the stresses and strain rate. Shear stress $`\sigma _{xz}`$ is a linear function of strain rate $`\dot{\gamma }v_x/z`$ for viscous fluids, and quadratic in $`\dot{\gamma }`$ for granular systems in the Bagnold grain-inertia regime. The latter result is rather general: When typical stresses on grains are large compared to the weight of individual particles but not large enough to significantly distort the spheres $`(1\sigma /\sigma _0k_n/k_0)`$, the only relevant time scale is $`\dot{\gamma }^1`$, which forces $`\sigma _{xz}\dot{\gamma }^2`$ simply by dimensional analysis. As an example, Fig. 4 shows the relationship between shear strain rate $`\dot{\gamma }`$ and shear stress $`\sigma _{xz}`$ for 2D and 3D cases. Below the transitional surface layer, and away from the bottom wall, both systems exhibit Bagnold scaling, indicated by the dashed lines in Fig. 4. Data for 2D suggest an offset of about $`1.5\sigma _0`$ in the stress, possibly due to corrections from the body force on individual spheres. Such an offset is not needed for an acceptable fit to the 3D data.
In order to probe the sensitivity of the results to the exact form of the stiff elastic response, we have also performed runs in the 3D system with linear damped springs, keeping all other parameters fixed. Remarkably, the packing fraction profiles and the normal stress anomaly remained virtually the same, and strain rate profiles were changed only by a global factor of about 1.35, suggesting that results are not too sensitive to the particular force scheme selected.
The lack of a regime with only surface flow is in contrast with experimental observations of avalanche flows in rotating drums. Since periodic boundary conditions are used, our simulations correspond to an infinitely large system with finite depth and fixed surface tilt, and therefore constitutes a different system. Experiments on chute flows also lack a regime of surface only flow; although this might be attributed to side-wall friction that necessitates higher tilt angles to initiate flow. Although there is no fundamental reason we can find that prohibits surface-only flows in this geometry, it appears that they are unlikely to occur in an assembly of monodisperse spheres.
DL is supported by the Israel Science Foundation under grant 211/97. Sandia is a multiprogram laboratory operated by Sandia Corporation, a Lockheed Martin Company, for the United States Department of Energy under Contract DE-AC04-94AL85000.
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# Polarized Sub-Millimetre Emission from Filamentary Molecular Clouds
## 1 Introduction
Molecular clouds are highly filamentary structures that are supported against rapid, global collapse by a combination of ordered magnetic fields as well as by non-thermal, hydromagnetic turbulence of some kind. However, determining the structure of the field has proven difficult. Direct Zeeman measurements have been limited to a few points per cloud rather than full-scale maps, and rarely trace the highest density gas (Crutcher et al. 1993, 1996, 1999). Polarization maps in optical and near-infrared extinction fail to reveal the structure of the field in the dense gas, perhaps as a result of the poor polarizing power of grains at these wavelengths in regions of high optical depth (Goodman et al. 1995).
Observations at far-infrared and sub-millimetre wavelengths have demonstrated that the thermal emission from dust grains is often partially polarized (eg. Schleuning 1998, Matthews & Wilson 2000, Hildebrand et al. 1999; see also review by Weintraub, Goodman, & Akeson 2000), which is usually attributed to the partial alignment of the emitting grains by the magnetic field. The mechanisms that are responsible for grain alignment are not entirely understood, but it is most often argued that dissipative processes gradually cause rapidly spinning grains to relax to a state in which grains are, on average, aligned perpendicular to the magnetic field (cf. Lazarian, Goodman, & Myers 1997). The possibility now exists, therefore, of using completely sampled JCMT polarimetry maps to place strong constraints on the structure of the magnetic fields in molecular clouds.
Current observations reveal the presence of well-ordered magnetic fields within molecular clouds and their clumps and cores. However, the results are quite puzzling in that they often show that the percent polarization $`P`$ decreases toward the highest density and most heavily obscured regions (eg. Schleuning 1998, Matthews & Wilson 2000; see also Weintraub et al. 2000 and references therein). This so-called “polarization hole” effect appears to be an almost universal feature of polarization maps at these wavelengths, although its precise origin has been a mystery. Nevertheless, two possible explanations have been proposed. Depolarization at high optical depths has often been attributed to the poor alignment of grains in dense regions (Goodman et al. 1995). The second possibility is that the effect might be due to disordered, sub-resolution scale variations in the orientation of aligned grains in dense regions, whose combined polarized emission would tend to mostly cancel (see review by Weintraub et al. 2000, for example). Recent high resolution interferometric polarization maps at millimetre wavelengths by Rao et al. (1998) have revealed strong, ordered polarization patterns within the polarization hole found by Schleuning. The position angle of the polarization vectors changes rapidly on the plane of the sky, which might be due to magnetic fields that are twisted on small scales.
The KL core is embedded in the “integral-shaped” filament of Orion A (See Johnstone & Bally 1999). There is some evidence that the polarization hole in the Orion filament is a global feature of the filament, and not localized to the KL region. Part of this filament has recently been mapped in 850 $`\mu m`$ polarization by Matthews and Wilson (2000), who find that the polarization vectors are remarkably ordered and generally aligned with the filament. They also find evidence for a decrease in the polarization percentage along the symmetry axis of the filament (See Matthews & Wilson 2000, Figure 2).
This paper demonstrates that a very simple explanation of the polarization hole effect is possible, namely that it arises as a consequence of a more general configuration of the magnetic field within filamentary clouds than has previously been contemplated by many authors. We show in this work that this effect can be rather naturally explained as arising from grain alignment by large-scale, ordered, helical magnetic fields threading filamentary molecular clouds. The Matthews and Wilson (2000) data suggest that the magnetic field threading the integral-shaped filament is probably very well-ordered. It is therefore conceivable that the polarization hole in this region might be due to a large scale, orderly twist in the magnetic field, rather than disordered magnetic fields on small scales. In addition, Vallée and Bastien (1999) compared the observed 760$`\mu m`$ polarization patterns of a number of sources with the qualitative patterns expected for several field geometries. For the subset of maps which show an ordered magnetic field, they argued that helical fields are most consistent with the data.
We present here, a very simple model of the polarized thermal emission of aligned grains to calculate the polarization patterns that would be expected from our helically magnetized filamentary cloud models (Fiege & Pudritz 2000a,b; hereafter FP1 and FP2). Polarimetric maps are difficult to interpret because very different three-dimensional field geometries can give rise to qualitatively similar polarization patterns (see the discussion in Sections 4.1 and 4 for example). This degeneracy is compounded by the high degree of uncertainty regarding the composition and distribution of grains, their polarization cross sections, and their alignment efficiencies. However, we show that under the simplest possible set of assumptions, helical field models can produce polarization patterns that are not unlike the observed patterns. Our most important result is that a large fraction of our models are depolarized along the symmetry axis of the filamentary cloud, as a result of the twisted field lines. Thus helical magnetic fields might provide a very simple, and purely geometric, explanation for the polarization hole phenomenon.
## 2 Filamentary Clouds and Helical Magnetic Fields
We first summarize the basic properties of our model of magnetized filamentary clouds (see FP1 and FP2 for details). Our models assume that filamentary clouds are self-gravitating, magnetized, truncated by an external pressure, and in a state of radial quasi-equilibrium between gravity, magnetic stresses, and internal pressure gradients. The ISM provides the external pressure for many filaments, but some are embedded in molecular clouds (ie. the integral-shaped filament of Orion A). We assume that the magnetic field is cylindrically symmetric, with both a poloidal and toroidal component, so that the field is helical in general. Note that the toroidal field component must vanish on the filament axis, where the poloidal field component is strongest. As one moves out in radius, the winding of the field on each subsequent cylinder increases as the toroidal field strength relative to the poloidal field increases. At large radii, the toroidal field decreases, but less rapidly than the poloidal field.
We constructed a three parameter family of models for helically magnetized filamentary clouds, which we constrained with an observational data set collected from the literature. We refer the reader to FP1 for the precise definition and a full discussion of the model parameters. Briefly, however, our model involves a concentration parameter $`C`$, which defines the radius of pressure truncation, and two flux to mass ratios $`\mathrm{\Gamma }_z`$ and $`\mathrm{\Gamma }_\varphi `$, for the poloidal and toroidal field components respectively, which are assumed to be constant within each filament. The models are characterized by a core radius $`r_0`$ defined as $`\sigma (4\pi G\rho _c)^{1/2}`$, where $`\sigma `$ is the one dimensional velocity dispersion of the gas (assumed to be constant), $`G`$ is the gravitational constant, and $`\rho _c`$ is the central density along the axis of the filament. The core radius characterizes the approximate radial scale at which the density structure changes from being nearly constant inside of $`r_0`$, to an approximately $`r^2`$ profile in the less dense outer envelope. The fact that we obtain an approximately $`r^2`$ profile is significant because three recent observations (Alves et al. 1998, Lada, Alves, & Lada 1999, and Johnstone & Bally 1999) all find density gradients that are consistent with an $`r^2`$ profile. Our models are therefore in good agreement with the existing data on the radial density structure of filamentary clouds.
## 3 A Simple Model for Sub-Millimetre Polarization by Aligned Dust Grains
The model that we adopt for the sub-millimetre polarization arising from our filamentary cloud models most closely resembles the analysis by Wardle and Königl (1990, hereafter WK90) for the polarized emission from the Galactic centre disk. The main difference is that our approach allows us to estimate the magnitude of the polarization percentage, whereas WK90 worked with normalized quantities. Another difference is that we include the possibility of several grain species, with the assumption that the emitting grains are a uniformly mixed population with the number density of each individual species proportional to the gas density. In addition, we require each grain species to be of uniform temperature and aligned to the same extent throughout the cloud.
Our assumptions regarding the constancy of grain properties and their alignment efficiencies throughout molecular clouds is undoubtedly a simplification of the true state of the grains in some clouds. As an example of the debate concerning grain properties, it has often been argued that the grains within dense regions of molecular clouds do not efficiently polarize starlight in extinction because they are either nearly spherical or poorly aligned at high optical depths (Goodman et al. 1995). There are several reasons to think that this might be the case. For instance, grain growth due to coagulation tends to produce slightly larger, and possibly more spherical grains in regions of higher density (Vrba et al. 1993), thus decreasing their contribution to the polarized emission. This effect might be compounded by the growth of ice mantles at high optical depths (Eiroa & Hodapp 1989, Goodman et al. 1995). The nature and efficiency of grain alignment mechanisms at high optical depths is perhaps even less certain than the grain properties themselves. The alignment might be poor because all efficient alignment mechanisms require supra-thermal rotational velocities, and neither the Purcell (1979) mechanism nor the radiative alignment mechanism (Draine & Weingartner 1996, 1997) are likely to be efficient at high optical depths. <sup>1</sup><sup>1</sup>1The radiative mechanism might be important within $`A_v2`$ of the surface of the cloud or any embedded sources (Draine & Weingartner 1996). Nevertheless, simple, quantitative estimates of the alignment efficiency and distribution of grain shapes for various physical conditions are presently lacking. It is for this reason that we mainly consider the simplest possible model of constant magnetic alignment throughout the cloud. We note however that our models can be readily extended to explicitly include radiative alignment.
### 3.1 Analysis
Grains result in polarized thermal emission if they are aspherical and partially aligned, most probably by a magnetic field. We consider the contribution to the polarized emission from several grain species $`j`$, whose absorption cross sections parallel and perpendicular to the symmetry axis are given by $`C_{,j}`$ and $`C_{,j}`$. The thermal emission from these grains is polarized if $`C_{,j}`$ and $`C_{,j}`$ are not equal. It is convenient to define the polarization cross section
$$C_{pol,j}=\{\begin{array}{cc}C_{,j}C_{,j}\hfill & \text{(oblate grains)}\hfill \\ 1/2(C_{,j}C_{,j})\hfill & \text{(prolate grains)}\hfill \end{array}$$
(1)
(Lee & Draine 1985, hereafter LD85; WK90).
To calculate the Stokes parameters $`I`$, $`Q`$, and $`U`$, which completely specify the intensity of the emission and its state of linear polarization, we sum the contributions arising from each species of dust within the cloud. We assume that species $`j`$ has temperature $`T_j`$, polarization cross section $`C_{pol,j}`$, and number density $`n_j=c_j\rho `$, where $`c_j`$ is a constant and $`\rho `$ is the total mass density of the cloud. We also assume that the mean alignment of the spin angular momentum with the magnetic field is given by
$$cos^2\gamma _j=\frac{(𝐁𝐉)^2}{B^2J^2}.$$
(2)
A full analysis of the polarized emission applicable at all wavelengths would include the effects of absorption and scattering. However, both of these effects can be neglected if we restrict our analysis to sub-millimetre wavelengths longer than about $`100\mu m`$ since the thermal emission is almost always optically thin at these long wavelengths (Hildebrand 1983) and scattering is completely insignificant (Novak et al. 1989). For example, the peak optical depth of even the dense KL core in Orion is only $`0.25`$ at $`100\mu m`$ (Novak et al. 1989, Schleuning 1998), and falls off with increasing wavelength as $`\lambda ^\beta `$, where $`\beta 1.52`$ (Hildebrand 1983, André, Ward-Thompson, & Barsony 2000). Therefore, it is appropriate to treat the sub-millimetre polarization as arising purely from emission in most circumstances.
Following DL85, we define the polarization reduction factor $`\mathrm{\Phi }_j`$ for dust grain species $`j`$ as
$$\mathrm{\Phi }_j=R_jF_j\mathrm{cos}^2\gamma ,$$
(3)
where $`R_j`$ is the Rayleigh polarization reduction factor due to imperfect grain alignment, $`F_j`$ is the polarization reduction due to the turbulent component of the magnetic field, and $`\gamma `$ is the angle between the plane of the sky and the local direction of the magnetic field. We assume that $`R_j`$ and $`F_j`$ are constants for each grain species. The combination of variables $`C_{pol,j}\mathrm{\Phi }_j`$ plays the role of the effective polarization cross-section, so that the contributions to the Stokes parameters $`Q`$ and $`U`$ from grains of species $`j`$ are given by the following:
$`Q_j`$ $`=`$ $`C_{pol,j}R_jF_jc_jB_\nu (T_j)c_jq`$ (4)
$`U_j`$ $`=`$ $`C_{pol,j}R_jF_jc_jB_\nu (T_j)c_ju,`$ (5)
where $`B_\nu `$ is the Planck function of the grain at temperature $`T_j`$, $`ds`$ is a distance element along the line of sight. The quantities $`q`$ and $`u`$ are integrals along the line of sight $`\widehat{s}`$, which depend on the density structure and the geometry of the magnetic field, but not the properties of the grains. We define the $`\widehat{x}\widehat{y}`$ plane as being parallel to the plane of the sky, with $`\widehat{x}`$ pointing West and $`\widehat{y}`$ pointing North. The angle $`\psi `$ is defined as the angle (counter-clockwise) between $`\widehat{y}`$ and the projection of the magnetic field in the plane of the sky. Then, $`q`$ and $`u`$ are respectively given by
$`q`$ $`=`$ $`{\displaystyle \rho \mathrm{cos}2\psi \mathrm{cos}^2\gamma ds}`$ (6)
$`u`$ $`=`$ $`{\displaystyle \rho \mathrm{sin}2\psi \mathrm{cos}^2\gamma ds}.`$ (7)
The angle $`\psi `$ varies as a function of $`s`$ along the line of sight for 3-dimensional field geometries where the local direction of the plane-of-sky field component varies along the line of sight. Note that $`q`$ and $`u`$ are identical in form to the “relative Stokes parameters” defined by WK90.
The polarization angle $`\chi `$ is determined by solving the equations
$`\mathrm{cos}2\chi ={\displaystyle \frac{q}{\sqrt{q^2+u^2}}},`$ (8)
$`\mathrm{sin}2\chi ={\displaystyle \frac{u}{\sqrt{q^2+u^2}}}`$ (9)
Since equations 8 and 9 are multi-valued, we restrict $`\chi `$ to be on the range $`[0,\pi )`$. The special case where $`u=0`$ and $`q`$ remains finite is very important to the analysis that follows. For this case, the polarization angle is given by $`\chi =0`$ when $`q>0`$ and $`\chi =\pi /2`$ when $`q<0`$.
We proceed to calculate the intensity of the emission, which requires the extinction cross-section from equation 3.14 of LD85:
$$C_{ext,j}=C_{+,j}\left[1\alpha _j\left(\frac{\mathrm{cos}^2\gamma }{2}\frac{1}{3}\right)\right],$$
(10)
where $`C_{+,j}`$ and $`\alpha _j`$ are defined as
$`C_{+,j}`$ $`=`$ $`{\displaystyle \frac{2C_{,j}+C_{,j}}{3}}`$
$`\alpha _j`$ $`=`$ $`{\displaystyle \frac{C_{pol,j}R_jF_j}{C_{+,j}}}.`$ (11)
Note that $`\alpha _j`$ does not vary with position, since we assume that $`R_j`$ and $`F_j`$ are constants for each grain species. The contribution of grain species $`j`$ to the intensity $`I`$ is given by
$$I_j=C_{+,j}B_\nu (T_j)c_j\left(\mathrm{\Sigma }\alpha _j\mathrm{\Sigma }_2\right),$$
(12)
where $`\mathrm{\Sigma }`$ is the surface density and $`\mathrm{\Sigma }_2`$ is a related quantity defined by
$`\mathrm{\Sigma }`$ $`=`$ $`{\displaystyle \rho 𝑑s}`$
$`\mathrm{\Sigma }_2`$ $`=`$ $`{\displaystyle \rho \left(\frac{\mathrm{cos}^2\gamma }{2}\frac{1}{3}\right)𝑑s}.`$ (13)
Note that $`\mathrm{\Sigma }`$ and $`\mathrm{\Sigma }_2`$ depend on the density and magnetic structure of the cloud, but not the grain properties.
The polarization percentage is defined by
$$p=\frac{\sqrt{Q^2+U^2}}{I},$$
(14)
where $`Q`$, $`U`$, and $`I`$ are obtained by summing equations 4, 5, and 12 over grain species $`j`$. It is easy to show that equation 14 becomes
$$p=\alpha \frac{\sqrt{q^2+u^2}}{\mathrm{\Sigma }\alpha \mathrm{\Sigma }_2},$$
(15)
where $`\alpha `$ is a weighted mean of $`\alpha _j`$ defined as follows:
$$\alpha =\frac{\mathrm{\Sigma }_j\alpha _jC_{+,j}B_\nu (T_j)c_j}{\mathrm{\Sigma }_jC_{+,j}B_\nu (T_j)c_j}$$
(16)
Equations 6, 7, and 15 show that the polarization pattern is determined by the density and magnetic structure of the gas, plus a single parameter $`\alpha `$ related to the grain cross-sections and alignment properties. We estimate $`\alpha `$ as follows. We assume that the optimal magnetic field geometry, in which $`\gamma =0`$ and $`\psi =const`$, gives rise to the maximum polarization percentage $`p_{max}`$ that is normally observed. The maximum polarization percentage is obtained from equation 15, with the help of equations 6 and 7:
$$p_{max}=\frac{\alpha }{1\alpha /6}.$$
(17)
Since the maximum polarization percentage observed in sub-millimetre polarization maps is rarely greater than $`10\%`$, equation 17 suggests that $`0.1`$ is a reasonable choice for $`\alpha `$. Using this choice, our approach predicts the magnitude of the polarization percentage obtained from our models, which is generally less than $`p_{max}`$. This is in contrast to WK90 who work with normalized quantities. We caution the reader that $`\alpha `$ could, in principle, vary from region to region, which would affect the magnitude of the polarization percentage that we predict. It would also slightly affect the shapes of the polarization profiles, since $`\alpha `$ apppears in the denominator of equation 15 and cannot be factored out.
## 4 Results
### 4.1 Transverse Fields Across Filamentary Clouds?
Polarimetry only maps the component of the magnetic field parallel to the plane of the sky. Suppose that the magnetic field threading the integral-shaped filament in Orion A contains only this plane-of-sky component. If this were true, the Matthews and Wilson (2000) map, as well as the smaller scale map by Schleuning (1998), would suggest that we are seeing either a filament or an edge-on sheet being impaled by a well-ordered field that is perpendicular to the filament axis or the midplane of the sheet. While this might explain the overall sense of the polarization vectors, it cannot easily account for the depolarization of these maps seen toward the centre of the putative sheet or filament. For such a transverse field model, it is easy to verify from equations 6, 7, and 15 that $`p`$ would be constant over the entire filament, since $`\zeta =0`$ everywhere and $`\psi =const`$. The analysis of Section 3.1 would not predict the observed depolarization along the axis of the filament or the midplane of the sheet. It could, however, be accomplished if the polarization were due to grains that are preferentially more spherical or poorly aligned in dense regions.
We show in Section 4.2 that helical fields threading a filamentary cloud can result in depolarization toward the filament axis, otherwise similar in appearance to what would be expected from the above transverse field scenario. However, the depolarization is due to the 3-dimensional structure of the field in this case and does not depend on the grain shapes or their alignment.
### 4.2 Polarization Patterns and Depolarization Effects for Helical Fields
We present polarization maps of our helically magnetized models of filamentary clouds using the method discussed in Section 3.1. We first restrict our parameter space by only showing results for filaments whose axes lie in the plane of the sky. In Section 4.4, we show maps for a filament at several inclination angles on the sky.
Figure 1 shows three representative models to illustrate the general types of behaviour that we find in our maps. We refer to these qualitative patterns as types 1 to 3 for simplicity, although we emphasize that the underlying models represent a continuum in parameter space (see FP1). In panel a) we show a model (type 1) where the polarization vectors are everywhere parallel to the filament. The most striking feature of the map is the depolarization along the axis. Qualitatively, this map shares some features with the Matthews and Wilson (2000) map of the Orion filament, specifically the overall orientation of the polarization vectors and the depolarization toward the central regions. Panel b) shows a model (type 2) whose polarization vectors are oriented opposite to the type 1 model in panel a). If such a pattern were observed, it could be misinterpreted as the result of a purely poloidal field. These models have two depolarized regions, with one on either side of the filament axis. The polarization on the filament axis is a local maximum, although the strongest polarization in the map is usually at the surface of the filament. Panel c) shows a more complicated case (type 3), where there are $`90^{}`$ flips in the orientation of the polarization vectors. The polarization percentage passes through zero at all locations where a flip in orientation occurs. We classify any patterns containing such flips in orientation as type 3.
It may seem surprising that filamentary clouds that are threaded by helical fields always have polarization vectors that are aligned parallel or perpendicular to the filament. This result is obtained regardless of whether the filament is parallel to the plane of the sky or inclined at at some angle. It is easy to understand this by the following argument (see Carlqvist & Kristen (1997) and FP1). Any line of sight passing through a filament encounters a given magnetic flux tube twice; first on the front side of the filament, and then on the back. By symmetry, the plane-of-sky components of the magnetic field vectors are mirror images of each other at these positions, so that $`\psi `$ is an odd function of the position $`s`$ along each line of sight. Therefore, equations 6 and 7 show that the contributions to $`u`$ exactly cancel, while the contributions to $`q`$ add in equal proportion.
From the definition of the polarization angle (equations 8 and 9), we observe that the polarization vectors are oriented perpendicular to the filament ($`\chi =0`$) when $`q>0`$ and parallel to the filament ($`\chi =\pi /2`$) when $`q<0`$. The dominant magnetic field component on each line of sight determines the orientation of the field vectors. This is easily seen by the following argument. Consider first a purely poloidal field threading a filament oriented in the North-South direction, parallel to the plane of the sky. The Stokes vector $`q>0`$ in this case, so that the polarization vectors are perpendicular to the filament, since $`\mathrm{cos}2\psi =1`$ everywhere. On the other hand, if the field were purely toroidal, $`\mathrm{cos}2\psi =1`$, so that $`q`$ is negative and the polarization vectors are parallel to the filament. Now consider the case of a helical field threading a filament oriented parallel to the plane of the sky. It is obvious that the poloidal and toroidal components of the magnetic field provide competing contributions to $`q`$ through $`\psi `$, and that the polarization direction at each position is determined by the field component that makes the dominant contribution to equation 6 along each line of sight.
The competing contributions to $`q`$ from the poloidal and toroidal field components may partially or completely cancel at some positions, whenever the poloidal and toroidal fields are of comparable strengths. This is the explanation for the depolarized regions seen in all three panels of Figure 1.
### 4.3 Exploration of the Parameter Space
We explore our parameter space in the same spirit as we did in FP1, FP2, and Fiege & Pudritz 2000c (Hereafter FP3). We compute polarization maps for models scattered randomly throughout our 3-dimensional parameter space ($`C`$, $`\mathrm{\Gamma }_z`$, $`\mathrm{\Gamma }_\varphi `$), subject to the observational constraints discussed in FP1. A slight complication is that we exclude some models that were found to be highly unstable to short wavelength, rapidly growing MHD instabilities in FP2. We note that a similar criterion was used in FP3 to select “parent” filaments for our prolate core models. Specifically, we exclude models for which $`\omega _{max}^2/(4\pi G\rho _c)>0.05`$, where $`\omega _{max}^2`$ is the maximum squared growth rate, and $`\rho _c`$ is the central density of the filament (See FP2, Figure 13).
In Figure 2, we show how the polarization pattern depends on $`B_{z,S}/B_{\varphi ,S}`$, where $`B_{z,S}`$ and $`B_{\varphi ,S}`$ are respectively the poloidal and toroidal field components, evaluated at the surface of the filament:
$$\frac{B_{z,S}}{B_{\varphi ,S}}=\frac{e^C\mathrm{\Gamma }_z}{\mathrm{\Gamma }_\varphi }.$$
(18)
The dots, squares, and x’s represent models of types 1, 2, and 3 respectively. Note that $`B_{z,S}/B_{\varphi ,S}`$ is small for many of our models because the toroidal magnetic field dominates in the outer envelope, near the radius of pressure truncation. The ratio of $`B_z/B_\varphi `$ is substantially higher in the interior regions of the filament.
An interesting feature of Figure 2 is that the polarization patterns are of the first type for a large portion of our models. Thus, many of our models are qualitatively similar to the Matthews and Wilson (2000) map of the Orion filament, discussed in Section 4.2. We find this type of pattern for most models with $`B_{z,S}/B_{\varphi ,S}\stackrel{<}{}0.1`$. Generally, polarization patterns of the second type occur when $`B_{z,S}/B_{\varphi ,S}\stackrel{>}{}0.33`$, and the third type occurs for intermediate values between about 0.1 and 0.33. This may be understood as follows. When $`B_{z,S}/B_{\varphi ,S}`$ is small, the models are dominated by the toroidal field component so that the poloidal field is ineffective at canceling the polarization due to the toroidal field, as discussed in Section 4.2. Thus, the polarization pattern is of type 1. However, as $`B_{z,S}`$ is increased relative to $`B_{\varphi ,S}`$, the competition becomes stronger until $`B_{z,S}`$ becomes dominant along some line of sight, which first occurs when $`B_{z,S}/B_{\varphi ,S}0.1`$. This results in a $`90^{}`$ flip in the polarization vectors at this position, which we categorize as the first Type 3 model. The exact opposite occurs when $`B_{z,S}/B_{\varphi ,S}`$ is increased past about $`0.33`$. Past this point, the toroidal field component becomes too weak compared to the poloidal field to produce a flip in the polarization vectors resulting in Type 2 patterns.
Figure 3 shows the degree to which the emission is depolarized for all three types of polarization pattern. We define the polarization hole depth as $`(P_{max}P_{min})/P_{max}`$, where $`P_{max}`$ is the maximum polarization percentage in the map and $`P_{min}`$ is the local minimum polarization percentage at the location of the polarization hole with the lowest polarization. We find that the polarization hole depth generally increases, with scatter, as a function of $`B_{z,S}/B_{\varphi ,S}`$ for Type 1 models from $`0\%`$ when $`B_{z,S}/B_{\varphi ,S}=0`$ to $`100\%`$ when $`B_{z,S}/B_{\varphi ,S}=0.1`$. This is easily understood by essentially the same argument given in the previous paragraph. There is no significant polarization hole when $`B_{z,S}/B_{\varphi ,S}`$ is small because contributions to the polarization arising from the poloidal field do not effectively cancel the larger contributions from the toroidal field along any line of sight. The depth of the polarization hole increases until $`B_{z,S}/B_{\varphi ,S}0.1`$, where the first Type 3 pattern emerges. Note that the polarization hole depth is always $`100\%`$ for Type 3 models, since $`P_{min}=0`$ at the locations where the oriention of the polarization vectors flips by $`90^{}`$. Increasing $`B_{z,S}/B_{\varphi ,S}`$ decreases the depth of the polarization hole for Type 2 models, since the contribution to the polarization from the toroidal field becomes progressively less effective at canceling the polarization due to the dominant poloidal field.
We define the width of the depolarized region as the distance between the polarization maxima on either side of the polarization hole. Panel 3b shows the ratio of this width divided by the filament diameter. Generally, this ratio increases with $`B_{z,S}/B_{\varphi ,S}`$ for Type 1 models, from nearly $`0`$ to about $`0.4`$, which occurs at the point where the models change to Type 3. This fractional width continues to increase past this point until $`B_{z,S}/B_{\varphi ,S}0.23`$, where the polarization width jumps discontinuously to the full width of the filament. This happens because of a change in the qualitative behaviour of the polarization vectors near the edge of the filament. Note that the polarization percentage shown in Figure 1 is a minimum at the outer edge of the filament for Type 1 models, but a maximum for Type 2 models. As we move through a sequence of Type 3 models, the polarization at the edge gradually increases as a function of $`B_{z,S}/B_{\varphi ,S}`$ until $`B_{z,S}/B_{\varphi ,S}0.23`$. Past this point, the maximum polarization for the filament occurs at the edge, so the width of the depolarized region becomes equal to the width of the filament.
### 4.4 The Effect of Inclination Angle
The polarization patterns due to helically magnetized filaments become more complicated for filaments that are inclined relative to the plane of the sky. As an example, in Figure 4, we show maps of one filament tilted $`10^{}`$, $`30^{}`$, and $`60^{}`$ relative to the plane of the sky. The filament model shown in Figure 4 is the same one as shown in the top two panels of Figure 1. Each polarization vector always remains parallel or perpendicular to the symmetry axis for the reasons discussed in Section 4.2 above, but the patterns become asymmetric relative to the central axis of the filament. It is particularly notable that filaments inclined relative to the line of sight usually result in asymmetric polarization patterns that include a $`90^{}`$ flip in the orientation of the polarization vectors on one or both sides of the filament.
## 5 Discussion & Summary
The main purpose of this paper is to calculate polarization maps for the models of filamentary molecular clouds that we presented in FP1 and FP2. When the filament is oriented parallel to the plane of the sky, the polarization patterns resulting from our models can be qualitatively classified into three main types of behaviour. The polarization vectors are parallel to the symmetry axis of the filament (Type 1) when the toroidal field is dominant, and perpendicular (Type 2) for models whose poloidal field dominates. The third type of pattern is more complicated, with polarization vectors that flip from being aligned parallel to the filament axis to perpendicular at some radius or radii (in projection). Generally, the first type of pattern occurs when $`B_{z,S}/B_{\varphi ,S}\stackrel{<}{}0.1`$, while the second type occurs $`B_{z,S}/B_{\varphi ,S}\stackrel{>}{}0.33`$. The third type of polarization pattern occurs at intermediate values of $`B_{z,S}/B_{\varphi ,S}`$ Our most important result is that helical fields result in “polarization holes,” in which the emission is depolarized at some positions in the interior of the filament. Our Type 1 models qualitatively agree with the polarization structure of the “integral-shaped” filament in Orion A (Matthews & Wilson 2000).
Many of our models result in polarization patterns that contain $`90^{}`$ flips in the orientation of the polarization vectors. The flips are symmetric about the symmetry axis when a filament with comparable poloidal and toroidal magnetic field strengths is oriented parallel to the plane of the sky. Filaments that are inclined relative to the line of sight generally result in asymmetric polarization flips on one or both sides of the central axis of the filament.
We have considered only models of non-fragmented, cylindrically symmetric filaments in this paper. This is obviously an idealization of real filamentary clouds, many of which have suffered gravitational fragmentation and formed strings of embedded cores (cf. Schneider & Elmegreen 1979, Dutrey et al. 1991; see also FP2 and FP3), and may also deviate from perfect cylindrical symmetry. We will address the polarization patterns of embedded cores in a future analysis. It is difficult to comment on the general effects of asymmetry. However, it is worthwhile to consider the simplest type of non-axisymmetric perturbation, that corresponding to a kink mode in which the filament is bent into a transverse sinusoidal wave. If the sinusoid lies in the plane of the sky, the front to back symmetry is preserved and the polarization vectors remain locally parallel or perpendicular to the filament. The depolarization along the central axis of the filament relies on a cancellation between contributions to the polarization from the central “backbone” of poloidal flux and the surrounding envelope dominated by the toroidal field component. Since these features would be preserved in such a perturbation, the polarization hole would also likely remain. If, on the other hand, the sinusoidal perturbation is perpendicular to the plane of the sky, then the front to back symmetry is broken and the polarization vectors would no longer be strictly parallel or perpendicular to the filament. A segment which is bent away from the observer is more compressed on the front side and more rarified on the back, so the front side would contribute more to the polarization. Assuming that the filament is oriented in the north-south direction and the field lines have right-handed (positive) helicity, the polarization vectors would shift toward the north-east/south-west plane if they were originally oriented in the north-south direction. The polarization hole would be preserved for the same reasons as in the case of a perturbation parallel to the plane of the sky. Thus, we see that non-axisymmetric distortions of a filament containing a helical field still ought to show a “polarization hole”.
Ignoring systematic variations of grain shapes and alignment efficiencies with density, we have shown that the “polarization hole” effect has a very simple and entirely geometric explanation if the field lines are twisted into helices. It might alternately be explained if the emitting dust grains are poorly aligned in the central regions of filamentary clouds, but this requires a greater understanding of dust grain physics.
More rigorous tests of our models will require a full set of polarization maps of filamentary clouds. Such a data set, when combined with well determined radial density profiles, will provide an excellent challenge to all cloud models, not only our own. Readers wishing to use our models to compare with data should contact J. Fiege to obtain fits of our model to data, or polarization maps for any inclination angle and set of input parameters ($`C`$, $`\mathrm{\Gamma }_z`$, and $`\mathrm{\Gamma }_\varphi `$). It is probably true that even with the density profiles and polarization structure of filamentary clouds determined by observations, it would still be necessary to make spatially resolved Zeeman measurements in order to conclusively prove the existence of helical fields. By making such maps, we might one day decide whether grain properties or the structure of the magnetic field is responsible for the polarization hole effect.
## 6 Acknowledgements
The authors thank B.C. Matthews and the anonymous referee for useful comments on a draft of the manuscript. J.D.F. acknowledges the financial support of McMaster University, CITA, and the Natural Sciences and Engineering Research Council of Canada (NSERC) during this research. The research of R.E.P. is supported by a grant from NSERC.
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# 1 Introduction
## 1 Introduction
One of the preoccupations of the circular accelerator designer is to estimate the influence of nonlinear forces on the single particle’s motion. These nonlinear forces manifest themselves as the systematic and random errors from optical elements, the voluntarily introduced functional ones, such as the sextupoles used for the chromaticity corrections and the octupoles used for stabilizing the particles’ collective motion, or from nonlinear beam-beam interaction forces. Even though the nonlinear forces mentioned above compared with the linear forces are usually very small, what is observed in reality, however, is that when the amplitudes of the transverse oscillation of a particle are large enough, the transverse motion might become unstable and the particle itself will finally be lost on the vacuum chamber. Apparently, the above implied maximum oscillation amplitudes, $`A_{x,y}`$, corresponding to the stable motions are functions of the specific longitudinal position, $`s`$, along the machine, and these functions $`A_{x,y}(s)`$ are the so-called dynamic apertures of the machine. A reasonably designed machine should satisfy the condition $`A_{x,y}(s)M_{x,y}(s)`$, where $`M_{x,y}(s)`$ are the mechanical cross-section dimensions of the vacuum chamber.
Needless to say the dynamic aperture problem in circular accelerators is one of the most challenging research topics for accelerator physicists, and the relevant methods adopted to treat this problem are quite different analytically and numerically. In this paper we will show how single sextupole, single octupole, and single decapole (single 2$`m`$ pole in general) in the machine
limit the dynamic aperture and what is their combined effect if there are more than one nonlinear element. From the established analytical formulae for the dynamic aperture one gets the scaling laws which relate the nonlinear perturbation strengths, beta functions, and the dynamic apertures. To test the validity of these formulae, comparisons with some numerical simulation have been made. As an interesting application the beam-beam limited dynamic apertures in a circular collider have been discussed.
## 2 Hamiltonian Formalism
The Hamiltonian formulation of dynamics is best known not only for its deep physical and philosophical inspiration to the physicists, but also for its technical convenience in solving various nonlinear dynamical problems. The general Hamiltonian for a particle of rest mass, $`m_0`$, and charge, $`e`$, in a magnetic vector potential, $`𝐀`$, and electric potential, $`\mathrm{\Phi }`$, is expressed as:
$$(q,p,t)=e\mathrm{\Phi }+c\left((𝐩e𝐀)^2+m_0c^2\right)^{1/2}$$
(1)
where $`c`$ is the velocity of light and $`𝐩`$ is the momentum with its components, $`p_i`$, conjugate to the space coordinates, $`q_i`$. The equations of motion can be readily written in terms of Hamilton’s equations:
$$\frac{dp_i}{dt}=\frac{}{q_i}$$
(2)
$$\frac{dq_i}{dt}=\frac{}{p_i}$$
(3)
For our specific dynamic problems in a circular accelerator it is convenient to chose curvilinear coordinates instead of Cartesian ones for us to describe the trajectory of a particle near an a priori known closed orbit. The Hamiltonian in the new system, $`(x,s,y)`$ (where $`x,s`$, and $`y`$ denote the coordinates in the Frenet-Serret normal, tangent, and binormal triorthogonal and right-handed coordinate system), is given by <sup>-</sup>:
$$H_t(q,p,t)=e\mathrm{\Phi }+c\left((p_xeA_x)^2+(p_yeA_y)^2+\left(\frac{p_seA_s}{1+x/\rho }\right)^2+m_0c^2\right)^{1/2}$$
(4)
where $`\rho `$ is the radius of curvature and the torsion of the closed orbit is everywhere zero. Since it is useful to use the variable, $`s`$, as the independent variable rather than time, $`t`$, one gets the new Hamiltonian by using a simple canonical transformation:
$$H_s=eA_s(1+x/\rho )\left(\frac{1}{c^2}(E^2m_0^2c^4)(p_xeA_x)^2(p_yeA_y)^2\right)^{1/2}e\mathrm{\Phi }$$
(5)
Noting that the term $`\frac{1}{c^2}(E^2m_0^2c^4)`$ in eq. 5 is equal to $`P^2`$ with $`P`$ being the total mechanical momentum of the particle, by using another trivial canonical transformation:
$$\overline{q}=q,\overline{s}=s,\overline{p}_{x,y}=\frac{p_{x,y}}{P_0},H=\frac{H_s}{P_0}$$
(6)
one gets another Hamiltonian:
$$H=\frac{eA_s}{P_0}(1+x/\rho )\left(\frac{P}{P_0}(\overline{p}_x\frac{eA_x}{P_0})^2(\overline{p}_y\frac{eA_y}{P_0})^2\right)^{1/2}\frac{e\mathrm{\Phi }}{P_0}$$
(7)
where $`P_0`$ is the mechanical momentum of the reference particle, and $`P=P_0+\mathrm{\Delta }P`$. Inserting $`eA_s/P_0`$ in eq. 7 by:
$$\frac{eA_s}{P_0}=\frac{B_yx^2}{2\rho ^2B_0}\frac{1}{B_0\rho }\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{^{n1}B_y}{x^{n1}}|_{x=0,y=0}(x+iy)^n$$
(8)
one gets finally the Hamiltonian which serves as the starting point of most of the dynamical problems concerning circular accelerators:
$$H=\frac{x^2B_y|_{x=0,y=0}}{2\rho ^2B_0}+\frac{1}{B_0\rho }\underset{n=1}{\overset{\mathrm{}}{}}\frac{1}{n!}\frac{^{n1}B_y}{x^{n1}}|_{x=0,y=0}(x+iy)^n$$
$$(1+x/\rho )\left(1+\frac{\mathrm{\Delta }P}{P_0}\left(\overline{p}_x\frac{eA_x}{P_0}\right)^2\left(\overline{p}_y\frac{eA_y}{P_0}\right)^2\right)^{1/2}\frac{e\mathrm{\Phi }}{P_0}$$
(9)
where $`B_0`$ is the bending magnetic field on the orbit of the reference particle, and $`B_y`$ in general is a complex variable.
## 3 Analytical formulae for dynamic apertures
To start with we consider the linear horizontal motion of the reference particle (no energy deviation) in the horizontal plane (y=0) assuming that the magnetic field is only transverse ($`A_x=A_y=0`$) and has no skew fields, and $`\mathrm{\Phi }`$ is a constant. The Hamiltonian can be simplified as
$$H=\frac{p^2}{2}+\frac{K(s)}{2}x^2$$
(10)
where $`x`$ denotes normal plane coordinate, $`p=dx/ds`$, and $`K(s)`$ is a periodic function satisfying the relation
$$K(s)=K(s+L)$$
(11)
where $`L`$ is the circumference of the ring. The solution of the deviation, $`x`$, is found to be
$$x=\sqrt{ϵ_x\beta _x(s)}\mathrm{cos}(\varphi (s)+\varphi _0)$$
(12)
where
$$\varphi (s)=_0^s\frac{ds}{\beta _x(s)}$$
(13)
As an essential step towards further discussion on the motions under nonlinear perturbation forces, we introduce action-angle variables and the Hamiltonian expressed in these new variables:
$$\mathrm{\Psi }=_0^s\frac{ds^{}}{\beta _x(s^{})}+\varphi _0$$
(14)
$$J=\frac{ϵ_x}{2}=\frac{1}{2\beta _x(s)}\left(x^2+\left(\beta _x(s)x^{}\frac{\beta _x^{}x}{2}\right)^2\right)$$
(15)
$$H(J,\mathrm{\Psi })=\frac{J}{\beta _x(s)}$$
(16)
Since $`H(J,\mathrm{\Psi })=J/\beta _x(s)`$ is still a function of the independent variable, $`s`$, we will make another canonical transformation to freeze the new Hamiltonian:
$$\mathrm{\Psi }_1=\mathrm{\Psi }+\frac{2\pi \nu }{L}_0^s\frac{ds^{}}{\beta _x(s^{})}$$
(17)
$$J_1=J$$
(18)
$$H_1=\frac{2\pi \nu }{L}J_1$$
(19)
Before going on further, let’s remember the relation between the last action-angle variables and the particle deviation $`x`$:
$$x=\sqrt{2J_1\beta _x(s)}\mathrm{cos}\left(\mathrm{\Psi }_1\frac{2\pi \nu }{L}s+_0^s\frac{ds^{}}{\beta _x(s^{})}\right)$$
(20)
Being well prepared, we start our journey to find out the limitations of the nonlinear forces on the stability of the particle’s motion. To facilitate the analytical treatment of this complicated problem we consider at this stage only sextupoles and octupoles (no skew terms) and assume that the contributions from the sextupoles and octupoles in a ring can be made equivalent to a point sextupole and a point octupole. The perturbed one dimensional Hamiltonian can thus be expressed:
$$H=\frac{p^2}{2}+\frac{K(s)}{2}x^2+\frac{1}{3!B\rho }\frac{^2B_z}{x^2}x^3L\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\delta (skL)+\frac{1}{4!B\rho }\frac{^3B_z}{x^3}x^4L\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\delta (skL)$$
(21)
Representing eq. 21 by action-angle variables ($`J_1`$ and $`\mathrm{\Psi }_1`$), and using
$$B_z=B_0(1+xb_1+x^2b_2+x^3b_3)$$
(22)
one has
$$H=\frac{2\pi \nu }{L}J_1+\frac{(2J_1\beta _x(s_1))^{3/2}}{3\rho }b_2L\mathrm{cos}^3\mathrm{\Psi }_1\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\delta (skL)$$
$$+\frac{(J_1\beta _x(s_2))^2}{\rho }b_3L\mathrm{cos}^4\mathrm{\Psi }_1\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\delta (skL)$$
(23)
where $`s_1`$ and $`s_2`$ are just used to differentiate the locations of the sextupole and the octupole perturbations. By virtue of Hamiltonian one gets the differential equations for $`\mathrm{\Psi }_1`$ and $`J_1`$
$$\frac{dJ_1}{ds}=\frac{H_1}{\mathrm{\Psi }_1}$$
(24)
$$\frac{d\mathrm{\Psi }_1}{ds}=\frac{H_1}{J_1}$$
(25)
$$\frac{dJ_1}{ds}=\frac{(2J_1\beta _x(s_1))^{3/2}}{3\rho }b_2L\frac{d\mathrm{cos}^3\mathrm{\Psi }_1}{d\mathrm{\Psi }_1}\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\delta (skL)$$
$$\frac{(J_1\beta _x(s_2))^2}{\rho }b_3L\frac{d\mathrm{cos}^4\mathrm{\Psi }_1}{d\mathrm{\Psi }_1}\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\delta (skL)$$
(26)
$$\frac{d\mathrm{\Psi }_1}{ds}=\frac{2\pi \nu }{L}+\frac{\sqrt{2}J_1^{1/2}\beta _x(s_1)^{3/2}}{\rho }b_2L\mathrm{cos}^3\mathrm{\Psi }_1\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\delta (skL)$$
$$+\frac{2\beta _x^2(s_2)}{\rho }J_1b_3L\mathrm{cos}^4\mathrm{\Psi }_1\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\delta (skL)$$
(27)
We now change these differential equations to the difference equations which are suitable to analyse the possibilities of the onset of stochasticity . Since the perturbations have a natural periodicity of $`L`$ we will sample the dynamic quantities at a sequence of $`s_i`$ with constant interval $`L`$ assuming that the characteristic time between two consecutive adiabatic invariance breakdown intervals is shorter than $`L/c`$. The differential equations in eqs. 26 and 27 are reduced to
$$\overline{J_1}=\overline{J_1}(\mathrm{\Psi }_1,J_1)$$
(28)
$$\overline{\mathrm{\Psi }_1}=\overline{\mathrm{\Psi }_1}(\mathrm{\Psi }_1,J_1)$$
(29)
where the bar stands for the next sampled value after the corresponding unbarred previous value.
$$\overline{J_1}=J_1\frac{(2J_1\beta _x(s_1))^{3/2}}{3\rho }b_2L\frac{d\mathrm{cos}^3\mathrm{\Psi }_1}{d\mathrm{\Psi }_1}\frac{(J_1\beta _x(s_2))^2}{\rho }b_3L\frac{d\mathrm{cos}^4\mathrm{\Psi }_1}{d\mathrm{\Psi }_1}$$
(30)
$$\overline{\mathrm{\Psi }_1}=\mathrm{\Psi }_1+2\pi \nu +\frac{\sqrt{2}\beta _x(s_1)^{3/2}\overline{J_1}^{1/2}}{\rho }b_2L\mathrm{cos}^3\mathrm{\Psi }_1+\frac{2\beta _x(s_2)^2}{\rho }\overline{J_1}b_3L\mathrm{cos}^4\mathrm{\Psi }_1$$
(31)
Eqs. 30 and 31 are the basic difference equations to study the nonlinear resonance and the onset of stochasticities considering sextupole and octupole perturbations. By using trigonometric relation
$$\mathrm{cos}^m\theta \mathrm{cos}n\theta =2^m\underset{r=0}{\overset{m}{}}\frac{m!}{(mr)!r!}\mathrm{cos}(nm+2r)\theta $$
(32)
one has
$$\mathrm{cos}^3\theta =\frac{2}{2^3}(\mathrm{cos}3\theta +3\mathrm{cos}\theta )$$
(33)
$$\mathrm{cos}^4\theta =\frac{1}{2^4}(2\mathrm{cos}4\theta +8\mathrm{cos}2\theta +\frac{4!}{((4/2)!)^2})$$
(34)
If the tune $`\nu `$ is far from the resonance lines $`\nu =m/n`$, where $`m`$ and $`n`$ are integers, the invariant tori of the unperturbed motion are preserved under the presence of the small perturbations by virtue of the Kolmogorov-Arnold-Moser (KAM) theorem. If, however, $`\nu `$ is close to the above mentioned resonance line, under some conditions the KAM invariant tori can be broken.
Consider first the case where there is only one sextupole located at $`s=s_1`$ with $`\beta _x(s_1)`$. Taking the third order resonance, $`m/3`$, for example, we keep only the sinusoidal function with phase $`3\mathrm{\Psi }_1`$ in eq. 30 and the dominant phase independent nonlinear term in eq. 31, and as the result, eqs. 30 and 31 become
$$\overline{J_1}=J_1+A\mathrm{sin}3\mathrm{\Psi }_1$$
(35)
$$\overline{\mathrm{\Psi }_1}=\mathrm{\Psi }_1+B\overline{J_1}$$
(36)
with
$$A=\frac{(2J_1\beta _x(s_1))^{3/2}}{4}\left(\frac{b_2L}{\rho }\right)$$
(37)
$$B=\sqrt{2}\beta _x(s_1)^{3/2}J_1^{1/2}\left(\frac{b_2L}{\rho }\right)$$
(38)
where we have dropped the constant phase in eq. 31 and take the maximum value of $`\mathrm{cos}^3(\mathrm{\Psi }_1)`$, 1. It is helpful to transform eqs. 37 and 38 into the form so-called standard mapping expressed as
$$\overline{I}=I+K_0\mathrm{sin}\theta $$
(39)
$$\overline{\theta }=\theta +\overline{I}$$
(40)
with $`\theta =3\mathrm{\Psi }`$, $`I=3BJ_1`$ and $`K_0=3AB`$. By virtue of the Chirikov criterion it is known that when $`|K_0|0.97164`$ resonance overlapping occurs which results in particles’ stochastic motions and diffusion processes. Therefore,
$$|K_0|1$$
(41)
can be taken as a natural criterion for the determination of the dynamic aperture of the machine. Putting eqs. 37 and 38 into eq. 41, one gets
$$|K_0|=3J_1\beta _x(s_1)^3\left(\frac{|b_2|L}{\rho }\right)^21$$
(42)
and consequently, one finds maximum $`J_1`$ corresponding to a $`m/3`$ resonance
$$J_1J_{max,sext}=\frac{1}{3\beta _x(s_1)^3}\left(\frac{\rho }{|b_2|L}\right)^2$$
(43)
The dynamic aperture of the machine is therefore
$$A_{dyna,sext}=\sqrt{2J_{max,sext}\beta _x(s)}=\frac{\sqrt{2\beta _x(s)}}{\sqrt{3}\beta _x(s_1)^{3/2}}\left(\frac{\rho }{|b_2|L}\right)$$
(44)
Eq. 44 gives the dynamic aperture of a sextuple strength determined case. The reader can confirm that if we keep $`\mathrm{sin}(\mathrm{\Psi }_1)`$ term instead of $`\mathrm{sin}(3\mathrm{\Psi }_1)`$ in eq. 35, one arrives at the same expression for $`A_{dyna,sext}`$ as expressed in eq. 44.
Secondly, we consider the case of single octupole located at $`s=s_2`$ with $`\beta _x(s_2)`$. Taking the forth order resonance, $`m/4`$, for example, we keep only the sinusoidal function with phase $`4\mathrm{\Psi }_1`$ in eq. 30 and the dominant phase independent nonlinear term in eq. 31, and as the result, we have eqs. 30 and 31 reduced to
$$\overline{J_1}=J_1+A\mathrm{sin}4\mathrm{\Psi }_1$$
(45)
$$\overline{\mathrm{\Psi }_1}=\mathrm{\Psi }_1+B\overline{J_1}$$
(46)
with
$$A=\frac{(J_1\beta _x(s_2))^2}{2}\left(\frac{b_3L}{\rho }\right)$$
(47)
$$B=2\beta _x(s_2)^2\left(\frac{b_3L}{\rho }\right)$$
(48)
where we have dropped the constant phase in eq. 31 and take the maximum value of $`\mathrm{cos}^4(\mathrm{\Psi }_1)`$, 1. By using Chirikov criterion, one gets
$$J_1J_{max,oct}=\frac{1}{2\beta _x(s_2)^2}\left(\frac{\rho }{|b_3|L}\right)$$
(49)
and the corresponding dynamic aperture:
$$A_{dyna,oct}=\sqrt{2J_{max,oct}\beta _x(s)}=\frac{\sqrt{\beta _x(s)}}{\beta _x(s_2)}\sqrt{\frac{\rho }{|b_3|L}}$$
(50)
Thirdly, without repeating, we give directly the dynamic aperture due to a decapole located at $`s=s_3`$:
$$A_{dyna,deca}=\sqrt{2\beta _x(s)}\left(\frac{1}{5\beta _x^5(s_3)}\right)^{1/6}\left(\frac{\rho }{|b_4|L}\right)^{1/3}$$
(51)
where $`b_4`$ is the coefficient of the decapole strength. Finally, we give the general expression of the dynamic aperture in the horizontal plane ($`z=0`$) of a single $`2m`$ ($`m3`$) pole component:
$$A_{dyna,2m}=\sqrt{2\beta _x(s)}\left(\frac{1}{m\beta _x^m(s(2m))}\right)^{\frac{1}{2(m2)}}\left(\frac{\rho }{|b_{m1}|L}\right)^{1/(m2)}$$
(52)
where $`s(2m)`$ is the location of this multipole. Eq. 52 gives us useful scaling laws, such as $`A_{dyna,2m}\left(\frac{\rho }{|b_{m1}|L}\right)^{1/(m2)}`$, and $`A_{dyna,2m}\left(\frac{1}{\beta _x^m(s(2m))}\right)^{\frac{1}{2(m2)}}`$.
If there is more than one nonlinear component, how can one estimate their collective effect ? Fortunately, one can distinguish two cases:
* If the components are independent, i.e. there are no special phase and amplitude relations between them, the total dynamic aperture can be calculated as:
$$A_{dyna,total}=\frac{1}{\sqrt{_i\frac{1}{A_{dyna,sext,i}^2}+_j\frac{1}{A_{dyna,oct,j}^2}+_k\frac{1}{A_{dyna,deca,k}^2}+\mathrm{}}}$$
(53)
* If the nonlinear components are dependent, i.e. there are special phase and amplitude relations between them (for example, in reality, one use some additional sextupoles to cancel the nonlinear effects of the sextupoles used to make chromaticity corrections), there is no general formula as eq. 53 to apply.
In the above discussion we have restricted us to the case where particles are moving in the horizontal plane, and the one dimensional dynamic aperture formulae expressed in eqs. 44, 50, 51 and 53, are the maximum stable horizontal excursion ranges with the vertical displacement $`y=0`$. In the following we will show briefly how to estimate the dynamic aperture in 2 dimensions when there is coupling between the horizontal and vertical planes. Now we consider the case where only one sextupole is located at $`s=s_1`$, and we have the corresponding Hamiltonian expressed as follows:
$$H=\frac{p_x^2}{2}+\frac{K_x(s)}{2}x^2+\frac{p_y^2}{2}+\frac{K_y(s)}{2}y^2+\frac{1}{3!B\rho }\frac{^2B_z}{x^2}(x^33xy^2)L\underset{k=\mathrm{}}{\overset{\mathrm{}}{}}\delta (skL)$$
(54)
Generally speaking, there exists no universal criterion to determine the start up of stochastic motions in 2D. Fortunately, in our specific case, we find out the similarity between the Hamiltonian expressed in eq. 54 and that of the Hénon and Heiles problem which has been much studied in literature . The Hénon and Heiles problem’s Hamiltonian is given by
$$H_{H\&H}=\frac{1}{2}\left(x^2+p_x^2+y^2+p_y^2+2y^2x\frac{2}{3}x^3\right)$$
(55)
when $`H_{H\&H}>1/6`$ the motion becomes unstable. The intuition we get from this conclusion is that there should exist a similar criterion for our problem, i.e. to have stable 2D motion one should have $`HH_{max}`$. Note that $`K_x(s)`$ and $`K_y(s)`$ in eq. 54 are equal to unity in the Hénon and Heiles problem’s Hamiltonian. The previous one dimensional result helps us now to find $`H_{max}`$. When $`y=0`$ one has $`H_{max}A_{dyna,sext,x}^2`$, since $`xA_{dyna,sext,x}`$. When $`y0`$, the crossing terms in eqs. 54 and 55 will play the role of exchanging energy between the two planes, and for a given set of $`x`$ and $`y`$ the total energy of the coupled system can not exceed $`H_{max}`$. If we define $`A_{dyna,sext,y}`$ is the dynamic aperture in $`y`$-plane, one has:
$$\beta _x(s_1)A_{dyna,sext,x}^2=\beta _y(s_1)A_{dyna,sext,y}^2+\beta _x(s_1)x^2$$
(56)
or:
$$A_{dyna,sext,y}=\sqrt{\frac{\beta _x(s_1)}{\beta _y(s_1)}(A_{dyna,sext,x}^2x^2)}$$
(57)
where $`\beta _y(s_1)`$ is the vertical beta function where the sextupole is located and $`A_{dyna,sext,x}`$ is given by eq. 44. Different from eq. 52, the derivation of eq. 57 is quite intuitive, hinted by the Hénon and Heiles problem which has been studied numerically instead of analytically in literature. From eq. 57 one understands that the difference between $`A_{dyna,sext,y}`$ and $`A_{dyna,sext,x}`$ comes from $`\sqrt{\beta _x(s_1)/\beta _y(s_1)}`$. If there are many sextupoles in a ring one usually has $`A_{dyna,sext,x}A_{dyna,sext,y}`$ since $`\beta _x(s_i)`$ will not be always larger or smaller than $`\beta _y(s_i)`$.
## 4 Comparison with simulation results
To verify the validity of eqs. 44, 50, 53, and 57, we compare the dynamic apertures of some special cases by using these analytical formulae with a computer code called BETA . Taking the lattice of Super-ACO as an example, we show the schematic layout of the machine in Fig. 1.
The horizontal beta function distribution and the working point in the third order tune diagram are given in Figs. 2 and 3, respectively. From Fig. 2 one finds that the horizontal beta function at the beginning and the end of the one turn mapping is 5.6 m. In the following numerical simulations the dynamic apertures correspond to $`\beta _x(s)=\beta _x(0)=5.6`$ m and it will also be used in the analytical formulae. Defining the sextupole, octupole, and decapole strengths $`S=b_2L/\rho `$ (1/m<sup>2</sup>), $`O=b_3L/\rho `$ (1/m<sup>3</sup>), $`D=b_4L/\rho `$ (1/m$`{}_{}{}^{4})`$, respectively, we make now a rather systematic comparison.
* A sextupole is located at $`s=s_1`$ with $`S(s_1)=1`$ and $`\beta _x(s_1)=13.6`$ m, and its influence on the horizontal dynamic aperture is illustrated in Figs. 4 and 5. From eq. 44 one gets analytically that $`A_{dyna,sext}=0.0385`$ m compared with the numerical value of 0.04 m.
* An octupole is located at $`s=s_1`$ with $`O(s_1)=10`$ and $`\beta _x(s_1)=13.6`$ m, and its influence on the horizontal dynamic aperture is illustrated in Figs. 6 and 7. From eq. 50 one gets analytically that $`A_{dyna,oct}=0.055`$ m compared with the numerical value of 0.054 m.
* The validity of eq. 51 has been checked also. If a decapole with strength $`D=1000`$ is located at $`s=s_1`$, the dynamic aperture in horizontal plane is shown in Figs. 8 and 9. From eq. 51 one gets $`A_{dyna,deca}=0.022`$ m compared with the numerical value of 0.024 m.
* A sextupole of $`S=2`$ and a octupole of $`O=62`$ are located at $`s=s_1`$ and $`\beta _x(s_1)=13.6`$ m, and their combined influence on the horizontal dynamic aperture is shown in Figs. 10 and 11. From eq. 53 one gets $`A_{dyna,total}=0.0145`$ m compared with the numerical value of 0.016 m.
* A sextupole of $`S=2`$ and a octupole of $`O=62`$ are located at $`s=s_1`$, $`\beta _x(s_1)=13.6`$ m, and $`s=s_2`$, $`\beta _x(s_2)=15.18`$ m, respectively, and their combined influence on the horizontal dynamic aperture is shown in Figs. 12 and 13. From eq. 53 one gets $`A_{dyna,total}=0.0138`$ m compared with the numerical value of 0.0135 m.
* Four sextupoles of $`S=2`$ are located at $`s=s_{1,2,3,4}`$ with $`\beta _x(s_1)=13.6`$ m, $`\beta _x(s_2)=15.18`$ m, $`\beta _x(s_3)=7.8`$ m, and $`\beta _x(s_4)=6.8`$ m, respectively, and their combined influence on the horizontal dynamic aperture is given in Figs. 14 and 15. From eq. 53 one obtained the analytical value $`A_{dyna,total}=0.012`$ m.
* We show how the dynamic apertures depend on the strengths of sextupole and octupole. Fig. 16 shows the comparison between the analytical (solid line) and the numerical (dotted) results for a sextupole located at $`s=s_1`$ in the machine shown in Fig. 2. Obviously, $`A_{dyna,sext}`$ scales with $`1/S`$. Fig. 17 gives the similar comparison for an octupole located at $`s=s_2`$, and confirms that $`A_{dyna,oct}`$ scales with $`1/\sqrt{S}`$.
* Now we change the tune of the machine a little bit (from $`\nu _x=1.7`$ to $`\nu _x=1.565`$), and the corresponding horizontal beta function distribution and the third order tune diagram are shown in Figs. 18 and 19. It is known that in this case $`\beta _x(0)=5.1`$ m. A sextupole of $`S=2`$ is located a $`s=s_1`$ with $`\beta _x(s_1)=12.42`$ m, and its influence on the dynamic aperture is shown in Figs. 20 and 21. From eq. 44 one finds $`A_{dyna,sext}=0.021`$ m compared with the numerical value of 0.02 m.
* Finally, a 2D dynamic aperture is calculated numerically and analytically. If a sextupole of $`S=2`$ is located at $`s_2`$ in the same lattice as in case (1) with $`\beta _x=15.18`$ m and $`\beta _y=4.26`$ m, the 2D dynamic aperture is calculated by using BETA and eq. 57 as shown in Fig. 22 and Fig. 23, respectively. The peak analytical dynamic apertures in horizontal and vertical planes are 0.0163 m and 0.031 m, compared with the numerical results of 0.017 m and 0.034 m, respectively.
To make the comparison much more clear we illustrate the machine parameters in Table 1 and the comparison results in Table 2.
The agreement between the analytical and the numerical simulation results is quite good. The fact that the analytically estimated 2D dynamic aperture agrees well with that calculated by the numerical method shows that the method we have used to treat the two coupled nonlinear oscillators is a reasonable heuristic procedure.
## 5 Beam-beam interaction limited dynamic apertures
Since we are interested in the single particle motion under beam-beam forces, the incoherent force should be taken into account :
$$F_{r,in}(r)=\pm \frac{nq^2(1+\beta ^2)}{2\pi ϵ_0r}\left(1\mathrm{exp}\left(\frac{r^2}{2\sigma ^2}\right)\right)$$
(58)
where $`n`$ is the line particle number density, $`\beta `$ is the particle’s velocity in units of the speed of light, $`q`$ is the particle’s electric charge, $`ϵ_0`$ is the permitivity in vacuum, $`r`$ is the transverse offset of a particle with respect to the center of the counter-rotating colliding bunch, $`\sigma `$ is the standard deviation of the Gaussian transverse charge density distribution, the positive and the negative signs correspond to the colliding bunches with the same or opposite charges, respectively. Now we expand $`F_{r,in}(r)`$ into Taylor series:
$$F_{r,in}(r)=\pm \frac{nq^2}{\pi ϵ_0}\left(\frac{1}{2\sigma ^2}r\frac{1}{8\sigma ^4}r^3+\frac{1}{48\sigma ^6}r^5\frac{1}{384\sigma ^8}r^7+\mathrm{}\right)$$
(59)
where we take $`\beta =1`$. To start with, we consider the particle’s motion in horizontal plane ($`y=0`$) and consider only the delta function nonlinear beam-beam forces coming from one IP. The differential equation of motion can be expressed as:
$$\frac{d^2x}{ds^2}+K_x(s)x=\pm \frac{nq^2}{\pi ϵ_0m_0c^2\gamma }(\frac{1}{2\sigma ^2}x\frac{1}{8\sigma ^4}x^3+\frac{1}{48\sigma ^6}x^5$$
$$\frac{1}{384\sigma ^8}x^7+\mathrm{})L_{k=\mathrm{}}^{\mathrm{}}\delta (skL)$$
(60)
where $`K_x(s)`$ describes the linear focusing of the lattice in the horizontal plane, $`m_0c^2`$ is the rest energy of the particle, $`\gamma `$ is the normalized particle’s energy, and $`L`$ is the circumference of the circular collider. The corresponding Hamiltonian is expressed as:
$$H=\frac{p_x^2}{2}+\frac{K_x(s)}{2}x^2\frac{nq^2}{\pi ϵ_0m_0c^2\gamma }(\frac{1}{4\sigma ^2}x^2\frac{1}{32\sigma ^4}x^4+\frac{1}{288\sigma ^6}x^6$$
$$\frac{1}{3072\sigma ^8}x^8+\mathrm{})L_{k=\mathrm{}}^{\mathrm{}}\delta (skL)$$
(61)
where $`p_x=dx/ds`$.
To make use of the general dynamic aperture formulae shown in section 3, one needs only to find the equivalence relations by comparing two Hamiltonians expressed in eqs. 21 and 61, respectively, and it is found that:
$$\frac{b_3}{\rho }L=\frac{N_eq^2}{8\pi ϵ_0m_0c^2\gamma \sigma ^4}$$
(62)
$$\frac{b_5}{\rho }L=\frac{N_eq^2}{48\pi ϵ_0m_0c^2\gamma \sigma ^6}$$
(63)
$$\frac{b_7}{\rho }L=\frac{N_eq^2}{384\pi ϵ_0m_0c^2\gamma \sigma ^8}$$
(64)
and so on, where we have replaced $`nL`$ by $`N_e`$ which is the particle population inside a bunch. Till now one can calculate all kinds of dynamic apertures due to nonlinear beam-beam forces. For example, one can get the dynamic apertures due to the beam-beam octupole nonlinear force:
$$A_{dyna,8,x}=\frac{\sqrt{\beta _x(s)}}{\beta _x(s^{})}\sqrt{\frac{\rho }{|b_3|L}}$$
$$=\frac{\sqrt{\beta _x(s)}}{\beta _x(s^{})}\left(\frac{8\pi ϵ_0m_0c^2\gamma \sigma ^4}{N_eq^2}\right)^{1/2}$$
(65)
and
$$A_{dyna,8,y}=\sqrt{\frac{\beta _x(s^{})}{\beta _y(s^{})}(A_{dyna,8,x}^2x^2)}$$
(66)
where $`s^{}`$ is the IP position. If we measure dynamic apertures by the beam sizes, one gets:
$$_{x,8}=\frac{A_{dyna,8,x}}{\sigma _x(s)}=\left(\frac{8\pi ϵ_0m_0c^2\gamma ϵ_x}{N_eq^2}\right)^{1/2}$$
(67)
where $`ϵ_x`$ is the bunch horizontal emittance. When the higher order multipoles effects ($`2m>8`$) can be neglected eqs. 65 and 66 give very good approximations to the 2D dynamic apertures limited by one beam-beam IP. If there are $`N_{IP}`$ interaction points in a ring the dynamic apertures described in eqs. 65 and 66 will be reduced by a factor of $`\sqrt{N_{IP}}`$ (if these $`N_{IP}`$ interaction points can be regarded as independent). Given the dynamic aperture of the ring without the beam-beam effect as $`A_{x,y}`$, the total dynamic aperture including the beam-beam effect can be estimated usually as:
$$A_{total,x,y}=\frac{1}{\sqrt{\frac{1}{A_{x,y}^2}+\frac{1}{A_{bb,x,y}^2}}}$$
(68)
Taking PEP-II B-Factory design parameters for example and assuming that the beams are round at IP, for the high energy ring, $`N_e=2.8\times 10^{10}`$, $`\gamma =1.76\times 10^4`$, and $`ϵ_x=49`$ nm, one gets from eq. 67, $`_{x,8}=3.2`$, and for the low energy ring, $`N_e=6\times 10^{10}`$, $`\gamma =6.07\times 10^3`$, and $`ϵ_x=49`$ nm, one finds $`_{x,8}=2.7`$.
## 6 Conclusion
We have derived the analytical formulae for the dynamic apertures in circular accelerators due to single sextupole, single octupole, single decapole (single 2$`m`$ pole in general), and the combination of many independent multipoles. The analytical results have been systematically compared with the numerical ones and the agreement is quite satisfactory. These formulae are very useful both for the physical insight and in the practical machine design and operation. One application of these formulae is to estimate analytically the beam-beam interaction determined dynamic aperture in a circular collider.
## 7 Acknowledgements
The author of this paper thanks B. Mouton for his help in using BETA program and also for his generating a flexible lattice based on the original Super-ACO one. The fruitful discussions with A. Tkachenko are very much appreciated.
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# Present limits to cosmic bubbles from the COBE-DMR three point correlation function.
## 1 Introduction
In the recent past the number of papers devoted to non-Gaussian anisotropies on the Cosmic Microwave Background (CMB) has increased dramatically. This new investigation field is, in fact, a powerful tool to distinguish between the theories of structure formation based on inflation and those based on topological defects. Quantum fluctuations produced in inflationary models are scale invariant and have a Gaussian distribution. Thus we expect that three-point correlation function of the CMB temperature vanish (Falk et al. 1993; Luo & Schramm 1993; Gangui & al. 1994).
On the contrary in models with topological defects the primordial density perturbations are scale dependent and non-Gaussian (Avelino et al. 1998): hence we expect some deviations from Gaussianity in higher order correlation functions. In this context we may also include the extended inflation model (La & Steinhardt 1989), because during the inflationary epoch we have a first order phase transition, that generates bubbles of true vacuum. These voids contribute together with ordinary quantum fluctuations to structure formation. This possibility has been investigated (Occhionero & Amendola 1994; Amendola et al. 1996): it has been shown (Occhionero et al. 1984, 1997) that primordial bubbles may be associated with the observation of large scale voids in several galaxy surveys (Kirshner et al. 1981; de Lapparent et al. 1989; da Costa et al. 1996; El Ad, Piran & da Costa 1996, 1997). Since these defects can also produce non-Gaussian anisotropies on the CMB, we may obtain some limits on the bubble parameters comparing observations with non-Gaussian predictions.
So far differents statistical tests have been applied to COBE-DMR sky maps (Kogut et al. 1996) and the results have been in agreement with the Gaussian models. Although recently two groups have detected non-Gaussian signal in COBE data (Ferreira et al. 1998; Pando et al. 1998), subsequently this has been shown to derive from a systematic effect in the data (Banday et al. 1999). On the other hand, Bromley & Tegmark (1999) tried to argue that the COBE 4 year data were in fact Gaussian.
In this paper we compare the level of non-Gaussianity produced in bubble models with the COBE data. We evaluate the three point correlation function in CDM models that contain also primordial bubbles. Comparing the numerical results with the COBE-DMR measures (Kogut et al. 1996; see also Hinshaw et al. 1994, 1995) we obtain upper limits on the parameters of the voids in agreement with galaxy surveys observations.
## 2 Method
The imprints of bubbles on the CMB has been studied in several papers (Baccigalupi, Amendola & Occhionero 1997 ; Amendola, Baccigalupi & Occhionero 1998; Baccigalupi & Perrotta 1999). The presence of the primordial voids induces a Sachs-Wolfe effect and an acoustic perturbation propagating up to the sound horizon on the photon distribution. As a consequence, the induced temperature fluctuations of individual bubbles are composed of a central spot and some concentric hotter isothermal rings; this pattern has been calculated by numerical integration of the Boltzmann equations in Amendola, Baccigalupi & Occhionero (1998). It is found that the bubble signal depends on the radius $`R`$ of the void and on the central (negative) density contrast $`\delta `$. We shall distribute $`N`$ voids on the CMB sky and use the fraction $`X`$ of the space that the voids fill today as a free parameter, where $`X=NR^3/3L_h^2\mathrm{\Delta }L_h`$ (Amendola et al. 1998) where $`L_h`$ is the horizon radius and $`\mathrm{\Delta }L_h`$ is the thickness of the last scattering surface. We consider bubbles of size $`R=30h^1`$Mpc at decoupling: due to their overcoming growth, these voids have today radii around $`2060h^1`$ Mpc, like those observed in galaxy surveys (da Costa et al. 1996; El Ad et al. 1996; 1997).
In simulated COBE maps, due to the low resolution of the satellite, the signal of the individual bubbles looks like dark spots confused amidst Gaussian anisotropies; their effect appears only when caculating the correlation functions of maps containing many bubbles. The temperature fluctuation may be decoupled in two terms:
$$\mathrm{\Delta }(\theta ,\phi )=\mathrm{\Delta }_{Gauss}(\theta ,\phi )+\mathrm{\Delta }_V(\theta ,\phi ).$$
(1)
The first term $`\mathrm{\Delta }_{Gauss}(\theta ,\phi )`$ is the Gaussian temperature fluctuation field produced by the primary anisotropies; the second term is the voids signal, that vanishes in directions where there are not bubbles. In order to compare the predictions of the model with experimental data, we calculate the collapsed three-point function
$`C_3(\alpha )`$ $`=`$ $`{\displaystyle \frac{1}{4\pi }}{\displaystyle \underset{l_1,l_2,l_3}{}}{\displaystyle \underset{m_1,m_2,m_3}{}}P_{l_1}(\mathrm{cos}\alpha )a_{l_1}^{m_1}a_{l_2}^{m_2}a_{l_3}^{m_3}`$ (2)
$`\times 𝒲_{l_1}𝒲_{l_2}𝒲_{l_3}_{l_1l_2l_3}^{m_1m_2m_3}.`$
where $`𝒲^l`$ is the window function of the experiment, $`P_l`$ are the Legendre polynomials, $`a_l^m`$ are the multipole coefficients of the spherical harmonic expansion and where
$`_{l_1l_2l_3}^{m_1m_2m_3}`$ $`=`$ $`(1)^{m_1}{\displaystyle \frac{\sqrt{(2l_1+1)(2l_2+1)(2l_3+1)}}{\sqrt{4\pi }}}\times `$ (3)
$`\times \left({\displaystyle \genfrac{}{}{0pt}{}{l_1\text{ }l_2\text{ }l_3}{0\text{ }0\text{ }0}}\right)\left({\displaystyle \genfrac{}{}{0pt}{}{l_{1\text{ }}l_2\text{ }l_3}{m_1m_2m_3}}\right).`$
with $`\left(\genfrac{}{}{0pt}{}{l_{1\text{ }}l_2\text{ }l_3}{m_1m_2m_3}\right)`$ the Wigner 3J symbol. We compare $`C_3(\alpha )`$ to the pseudo three-point collapsed function of Kogut et al. (1996): obviously the two are identical due to the absence of noise in our case. Since the Gaussian term and the signal of the bubbles are not correlated, we may write the three-point correlation function as sum of two separate contributions:
$$C_3=C_3^{Gauss}+C_3^V.$$
(4)
The contribution to $`C_3`$ from gaussian fluctuations is not zero. This contribution may arise from non-linearities in the inflationary dynamics or from non-linear growth of the perturbations. However, using the analytical expression for the $`C_3^{Gauss}(\alpha )`$ computed in Gangui et al. (1994) it can be seen that the level of non-Gaussianity produced by the non-linearities in the inflation dynamics is smaller than that arising from the non-linear growth (Mollerach et al. 1995), and that the latter is much smaller that produced by the bubbles. In fact, on the angular scales probed by COBE-DMR, Mollerach et al. (1995) found an amplitude $`C_3^{RS}(\alpha )0.1`$ $`\mu K^3`$, while we find that the contribution of the voids is larger by several orders of magnitude: $`C_3^V(\alpha )10^4`$ $`\mu K^3`$. Therefore, we neglect $`C_3^{Gauss}`$ in the following.
To calculate $`C_3^V(\alpha )`$ we use the same approch of texture-spot anisotropies (Magueijo 1995, Gangui & Mollerach 1996, 1997). The temperature fluctuations produced by a random distribution of bubbles, in the $`\widehat{\gamma }`$ direction, is simply the superposition of the signal of all the bubbles, and can be written as $`\mathrm{\Delta }_V(\widehat{\gamma })=_nb_nf_n(\widehat{\gamma }_n,\alpha )`$. Here, the signal of the n-th bubble has been decomposed as a overall amplitude $`b_n`$ (corresponding to the central temperature fluctuation) and a density profile $`f_n(\widehat{\gamma }_n,\alpha )`$ where $`\alpha `$ is the angle measured from the bubble center. The dependence on the parameters $`R_n`$ and $`\delta _n`$ is contained only in the amplitude; it is to be expected that this dependence is linear in $`\delta _n`$ and quadratic in $`R_n`$, since the central temperature fluctuation is dominated by the Sachs-Wolfe effect. The expression $`b_n=\delta _n(R_n/20H^1)^2`$ is indeed an accurate fit in the range we are interested (see Amendola et al. 1998). The profile $`f_n(\widehat{\gamma }_n,\alpha )`$ contains the full effect of the acoustic oscillations and the adiabatic fluctuations, and has been obtained numerically in Baccigalupi & Perrotta (1999).
We expand $`\mathrm{\Delta }_V`$ in spherical harmonics and obtain the multipole coefficients
$$a_l^m=\frac{4\pi }{2l+1}\underset{n}{}b_nF_n^lY_l^m^{}(\widehat{\gamma }_n),$$
(5)
where $`F_n^l`$ is the Legendre trasform of the intensity profile,
$$F_n^l=\frac{2l+1}{2}𝑑\mathrm{\Omega }_\alpha f_n(\widehat{\gamma }_n,\alpha )P_l(\mathrm{cos}\alpha ).$$
(6)
Inserting (5) in (2) the collapsed function reduces to:
$`C_3^V(\alpha )`$ $`=`$ $`4\pi {\displaystyle \underset{l_1,l_2,l_3}{}}P_{l_1}(\mathrm{cos}\alpha )𝒲^{l_1}𝒲^{l_2}𝒲^{l_3}𝐉^{l_1l_2l_3}`$ (7)
$`\times {\displaystyle \underset{n_1,n_2,n_3}{}}b_{n_1}b_{n_2}b_{n_3}F_{n_1}^{l_1}F_{n_2}^{l_2}F_{n_3}^{l_3},`$
where $`𝐉^{l_1l_2l_3}`$ represents
$$𝐉^{l_1l_2l_3}=\left(\genfrac{}{}{0pt}{}{l_1l_2l_3}{000}\right)^2$$
(8)
We take the window function of COBE to be $`e^{l(l+1)\sigma ^2/2}`$, with $`\sigma =3.2^{}`$. We assume now that there are $`N`$ identical voids on the CMB sky. Developing the sum on $`n_1`$, $`n_2`$ and $`n_3`$ we obtain three terms that represent the contribution to the $`C_3^V(\alpha ),`$ when the bubble signals are not correlated and when are correlated two by two or three by three and etc. We take into account the correlation at lowest order, in other words we consider just the first two terms. Then the mean value of (7) for a Poissonian bubble distribution on the sky is obtained substituting the sum on the bubble index with an integral over the whole sky. In fact, the number of bubbles in a circular ring centered on a single bubble is proportional to angular extension of the ring, therefore we have:
$`C_3^V(\alpha )`$ $`=`$ $`4\pi \delta ^3\left({\displaystyle \frac{R}{20H^1}}\right)^6N{\displaystyle \underset{l_1,l_2,l_3}{}}P_{l_1}(\mathrm{cos}\alpha )`$ (9)
$`\times 𝒲^{l_1}𝒲^{l_2}𝒲^{l_3}𝐉^{l_1l_2l_3}I^{l_1l_2l_3},`$
where
$$I^{l_1l_2l_3}=F^{l_1}F^{l_2}F^{l_3}+\frac{3}{2}F^{l_1}F^{l_2}F^{l_3}(\theta )d(\mathrm{cos}\theta ),$$
(10)
and
$$F^{l_3}(\theta )=f(\theta +\alpha )P_{l_3}(\mathrm{cos}\alpha )d(\mathrm{cos}\alpha ).$$
(11)
Using the same approch, after a tedious calculation, we have found an analitycal expression for the variance $`\sigma _V^2(\alpha )=C_3^V(\alpha )^2C_3^V(\alpha )^2`$, that we do not report for shortness. We compare the experimental data with the behaviour of the $`C_3^V(\alpha )`$ for differents values of the parameters $`\delta `$ and $`X`$. When $`C_3^V(\alpha )\pm \sigma _V(\alpha )`$ is larger than COBE data plus the noise and cosmic variance, we have some constraints on the parameters of our model.
## 3 Results
The COBE data has been taken from Kogut et al. (1996). We assume a fraction of bubbles corresponding to $`0.31<X<0.54`$, consistent with da Costa et al. (1997). We have computed the $`C_3^V(\alpha )`$ for $`0.001<\delta <0.0026`$, without dipole and quadrupole contribution, $`l_{\mathrm{min}}=4`$. In the figures we report the behaviour of the $`C_3^V(\alpha )`$ for two values of $`\delta =0.002,0.0012`$. The oscillating behaviour of the plots is due to the sum of the Legendre polynomials in (9). In the plots the errorbars are the $`\sigma (\alpha )`$’s. The level of the cosmic variance $`\sigma (\alpha )`$ generated from the model is very high for $`\alpha <40^{}`$, while it is small on the large angular scales, $`\alpha >45^{}`$, where the contribution of the lowest multipoles is small. In figure (1) we have the model with $`\delta =0.002`$: we may note that for $`X=0.54`$ the signal is larger than cosmic variance and the observed data points, while $`X=0.31`$, the plot is marginally consistent with the experimental data.
In figure (2) we report the $`C_3^V(\alpha )T_0^3`$ for $`\delta =0.0012`$: it fits the COBE data very well.
Notice that in the range $`\alpha >50^{}`$ the behaviour of the collapsed function seems to follow the trend of the COBE measures. Values of $`\delta <0.0012`$ produce a $`C_3^V(\alpha )`$ within the cosmic variance band and smaller than the COBE data. In this case the observations do not impose constraints and we may obtain only an upper limit on the value of $`\delta `$. We have applied a $`\chi ^2`$ analysis to our models. In figure (3) we report the confidence regions with a confidence level set to 99.9% (grey region) and to 99.5% (black region).
We may note that all models with $`\delta 0.0017`$ are ruled out by the experimental data. Then we may conclude that although the bubbles produce a non-Gaussian signal on the CMB, this is in agreement with the present observation provided that the density contrast $`\delta 0.0017`$ or $`X54\%`$. So we obtain a constraint stronger than that found in Amendola et al. (1998), where the bubble power spectrum was compared to the measures of the CAT experiment. The next high resolution experiments, like MAP and Planck, and the recent observations of Boomerang and Maxima should be able to detect the voids signal on the CMB. In fact these missions can probe the multipoles $`l>100`$, where the contribution of the bubbles is important, and the effects on $`C_3(\alpha )`$ may be large.
## 4 Conclusion
Several galaxy surveys found huge spherical voids in the matter distribution, the galaxies lying in the surrounding shells: these structures may be generated in inflationary models with first order phase transitions. These bubbles produce a non-Gaussian signal on the CMB. We analyse this signal developing an analytical expression for the three-point collapsed function of a bubble distribution, using the formalism of Magueijo (1995). Our free parameters are the density contrast and volume fraction of the bubbles, while the radius $`R`$ is fixed to a value consistent with the galaxy surveys. We compare the behaviour of the three point collapsed function for the bubble model with the COBE data. We obtain a constraint on the value of $`\delta `$: in fact, the existence of the voids at decoupling is not in contrast with the measures of the COBE three-point collapsed function, provided $`\delta 0.0017`$ or $`X0.54`$. This still leaves plenty of room for the bubbles to cooperate efficiently to structure formation, both via the central voids and via the possibility of shocking on the outher shell: in fact a central density contrast of $`0.001`$ can still evolve linearly in an empty void by today. More information will be obtained comparing the results of the future high resolution experiments.
###### Acknowledgement 1
L.A. and F. O. acknowledge financial support from the Italian Ministry of University Research and Scientific Technology
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# I Introduction
## I Introduction
Recently, Randall and Sundrum made an proposal that we may live in a four dimensional section of 5 dimensional universe . The key point of the first one is that gravitational warp factor can generate a small number to solve the hierarchy problem. The idea of the second one is that due to the the attraction of the brane energy, the metric fluctuation around the domain wall admits a bound state and the background cosmological constant makes this ’bound state’ a zero-mode to be identified as a graviton. The brane tension is fine tuned so that the effective cosmological constant on the domain wall is zero and the brane is stationary. If the fine tuning is relaxed, bulk metric depends on time and the brane inflates. These solutions can be used to discuss the cosmology of the RS-model.
More recently, Kraus pointed out that these solutions can be interpreted as motion of a domain wall in a stationary background . Ref. discussed the motion of a ’probe brane’ in fixed background using the Dirac-Born-Infeld action. The idea of is that for the given cosmological constant we can glue two pieces of AdS-Schwartzshild (AdSS) solutions along the moving domain wall using the Israel matching conditions to construct the new space-time containing the domain wall. The matching condition determines the motion of the domain wall in terms of the domain wall tension and the bulk cosmological constant.
The moving brane solutions are interesting since those can lead us scenarios for the graceful exit from the inflation in the Randall-Sundrum type cosmology. For example, suppose we live in a brane whose brane tension is fine tuned to have flat 4 dimensional geometry. Now if there is a moving brane with finite thickness, the collision of our brane with such a moving brane provide inflation in our brane world for the finte collision time, as we will argue. Having explicit analytic solution for the two moving branes with ’finite thickness’ is not an easy task. So finding solutions for two moving thin branes in a fixed background is the first step in this direction and this is one of goals of this paper. Generalization of the RS1 model to time dependent case has been discussed in wherein either the distant between the two branes is fixed or the bulk metric itself has explicit time dependence.
We first examine the case with orbifold symmetry, which is the generalization of the Randall-Sundrum to the moving brane case. We try to solve the full 5 dimensional bulk Einstein solution with the warp factor corresponding to the moving branes. It turns out that the Einstein equation dictates that the information of the brane location must disappear from the bulk metric. However the motion of the brane can not be arbitrary: it is fixed by the tension of the brane and the bulk cosmological constant. We will get a solution with two branes; one is stationary and the other is moving. In the presence of the reflection symmetry, the tensions of the two branes must have opposite sign. To get the solution where both of the branes have positive tension, we have to abandon the orbifold symmetry. We use Kraus’ method. We show that the 4-dimensional effective cosmological constant on the brane world in the absence of the reflection symmetry, is not well defined. We find a condition for a brane to be stationary. Finally we consider the brane scattering using these solutions and suggest a scenario for the inflation during finite duration.
## II Summary of RS-model
We start with following action:
$$S=S_{gravity}+\mathrm{\Sigma }_iS_{i(\mathrm{wall})},$$
(1)
where $`S_{gravity}`$ is given by
$$S_{gravity}=d^5x\sqrt{g}\left[\frac{1}{2\kappa ^2}\mathrm{\Lambda }\right],$$
(2)
with $`\mathrm{\Lambda }`$ being a cosmological constant and $`g_{\mu \nu }`$ a metric of five-dimensional space-time. The domain wall action can be written as
$$S_{i(wall)}=d^5x\sqrt{g}_{i(wall)}\delta (yR_i(t)),$$
(3)
where $`y`$ is a coordinate of a transverse direction. Here the delta function implies that the domain wall lies at the position $`y=R_i(t)`$ and can possibly move.
Assuming that all excitation modes of the matter on the wall are absent, the action Eq. (1) is reduced to the
$$S=d^5x\sqrt{g}\left[\frac{1}{2\kappa ^2}\mathrm{\Lambda }\right]\underset{i}{}\sigma _id^5x\sqrt{g}\delta (yR_i(t)),$$
(4)
where $`\sigma _i`$ is a tension of the $`i`$-th domain wall. Since we are interest to the AdS space-time in the bulk, only the case $`\mathrm{\Lambda }<0`$ will be considered. From Eq. (4), Einstein equations become
$$_{}^{\mu }{}_{\nu }{}^{}\frac{1}{2}\delta _{}^{\mu }{}_{\nu }{}^{}=\kappa ^2T_{}^{\mu }{}_{\nu }{}^{}.$$
(5)
Here, the energy-momentum tensor $`T_{}^{\mu }{}_{\nu }{}^{}`$ is given by
$`T_{}^{\mu }{}_{\nu }{}^{}=|\mathrm{\Lambda }|\mathrm{diag}(1,1,1,1,1)+\mathrm{\Sigma }_iT_{i(wall)}^{}{}_{}{}^{\mu }{}_{\nu }{}^{},`$ (6)
$`T_{i(wall)}^{}{}_{}{}^{\mu }{}_{\nu }{}^{}=\sigma _i\mathrm{diag}(1,1,1,1,0)\delta (yR_i(t)).`$ (7)
where $`T_{i(wall)}^{}{}_{}{}^{\mu }{}_{\nu }{}^{}`$ is an energy-momentum tensor on the i-th wall.
The vacuum solution of the Einstein equation in the presence of the negative cosmological constant $`\mathrm{\Lambda }`$ is AdS space: The Randall-Sundrum model is constructed by joining two AdS
$$ds^2=e^{2ky}(dt^2+\delta _{ij}dx^idx^j)+dy^2,$$
(8)
where
$$k^2=\frac{\kappa ^2}{6}|\mathrm{\Lambda }|,$$
(9)
The Randall-Sundrum solution is that domain wall solution can be obtained by joining two ’inner-part’ of the AdS spaces along the hyper surface at $`y=0`$ :
$$ds^2=e^{2ky}(dt^2+\delta _{ij}dx^idx^j)+dy^2.$$
(10)
The solution has manifest $`𝐙_\mathrm{𝟐}`$ summetry $`yy`$. The domain wall at $`y=0`$ has tension given by
$$\sigma =6k/\kappa ^2.$$
(11)
Now, for later purpose, let’s ask what should be the solution if we put the domain wall along $`y=R`$. Due to the $`𝐙_\mathrm{𝟐}`$ symmetry, the metric should be as follows:
$`ds^2`$ $`=`$ $`e^{2kyR}(dt^2+\delta _{ij}dx^idx^j)+dy^2\mathrm{for}\text{near y=R},`$ (12)
$`=`$ $`e^{2ky+R}(dt^2+\delta _{ij}dx^idx^j)+dy^2\mathrm{for}\text{near y=-R},`$ (13)
The sign in the warp factors is chosen such that as $`y\mathrm{}`$ we are going deep inside the AdS spaces. Physically, there is no difference whether we have domain wall at $`y=0`$ or $`y=R`$ apart from the over all scale $`e^{2kR}`$. This can be verified by observing that above metric can be written as
$$ds^2=e^{2k|y|+2kR}(dt^2+\delta _{ij}dx^idx^j)+dy^2,\text{for }|y|R.$$
(14)
Notice that $`y=R`$ and $`y=R`$ are identified and there is no domain wall at $`y=0`$ since the region $`RyR`$ is not in the universe in this construction. This is the 1-wall solution with positive tension (RS2) .
However, in the region $`RyR`$, the same expression (13) can be written as
$$ds^2=e^{2k|y|2kR}(dt^2+\delta _{ij}dx^idx^j)+dy^2.$$
(15)
This means that if $`k`$ is positive, there is a domain wall with negative tension at $`y=0`$, as well as positive tension domain wall at $`y=R`$. Therefore (15) represent solution with two domain walls (RS1). Therefore, according to the region we are looking at, the solution (13) can gives either RS2 or RS1. The RS solution is not a geometry for one ads space with a brane but a patching two ads along a brane. See figure 1(a,b). Notice that the figure 1.(b) also tells us that as $`R\mathrm{}`$, RS1 is reduced to RS2. Later on, we will see that replacing $`RR(t)`$ is still solution for specific $`R(t)`$.
If a constant $`R`$ is replaced with time dependent $`R(t)`$, we would be considering a moving domain wall. In the next section we are looking for a solution which has a warp factor corresponding moving domain walls.
## III Moving Domain Walls
### A RS1-type model
In this section, we look for a solution with moving domain walls located at $`y=R(t)`$ and $`y=0`$. If we impose $`𝐙_\mathrm{𝟐}`$ symmetry, there must be a mirror image of the moving wall at $`y=R(t)`$. So we are looking for the moving version of the solution (13). Therefore we start from the following ansatz for the region $`R(t)yR(t)`$:
$$ds^2=e^{2k|y|+2kR(t)}(dt^2+H(t,y)^2\delta _{ij}dx^idx^j)+b(t,y)^2dy^2,$$
(16)
where $`k`$ in Eq. (15) is replaced with $`k`$ and this metric is RS1-type metric which has two walls at $`y=0`$ and $`y=R(t)`$. Since we impose $`𝐙_\mathrm{𝟐}`$ symmetry, the above metric must be invariant under $`yy`$. To determine $`H(t,y)`$ and $`b(t,y)`$, we have to solve five-dimensional Einstein equations. If we write the metric as
$$ds_{5}^{}{}_{}{}^{2}=n(t,y)^2dt^2+a(t,y)^2\delta _{ij}dx^idx^j+b(t,y)^2dy^2.$$
(17)
Einstein equations are given by
$`G_{00}`$ $`=`$ $`3\left[{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{b}}{b}}\right){\displaystyle \frac{n^2}{b^2}}\left\{{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{b^{}}{b}}\right)\right\}\right]=\kappa ^2T_{00},`$ (18)
$`G_{ii}`$ $`=`$ $`{\displaystyle \frac{a^2}{b^2}}\left[{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+2{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{b^{}}{b}}\left({\displaystyle \frac{n^{}}{n}}+2{\displaystyle \frac{a^{}}{a}}\right)+2{\displaystyle \frac{a^{\prime \prime }}{a}}+{\displaystyle \frac{n^{\prime \prime }}{n}}\right]`$ (20)
$`+{\displaystyle \frac{a^2}{n^2}}\left[{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}+2{\displaystyle \frac{\dot{n}}{n}}\right)+{\displaystyle \frac{\dot{b}}{b}}\left(2{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{\dot{n}}{n}}\right)2{\displaystyle \frac{\ddot{a}}{a}}{\displaystyle \frac{\ddot{b}}{b}}\right]=\kappa ^2T_{ii},`$
$`G_{55}`$ $`=`$ $`3\left[{\displaystyle \frac{a^{}}{a}}\left({\displaystyle \frac{a^{}}{a}}+{\displaystyle \frac{n^{}}{n}}\right){\displaystyle \frac{b^2}{n^2}}\left\{{\displaystyle \frac{\dot{a}}{a}}\left({\displaystyle \frac{\dot{a}}{a}}{\displaystyle \frac{\dot{n}}{n}}\right){\displaystyle \frac{\ddot{a}}{a}}\right\}\right]=\kappa ^2T_{55},`$ (21)
$`G_{05}`$ $`=`$ $`3\left[{\displaystyle \frac{n^{}}{n}}{\displaystyle \frac{\dot{a}}{a}}+{\displaystyle \frac{a^{}}{a}}{\displaystyle \frac{\dot{b}}{b}}{\displaystyle \frac{\dot{a}^{}}{a}}\right]=0,`$ (22)
where we use dot and prime to describe a derivative with respect to $`t`$ and $`y`$ respectively. Due to $`𝐙_\mathrm{𝟐}`$ symmetry, it is sufficient that we pay attention to the right-hand side of the domain wall only. Our ansatz means that, $`n(t,y)`$ and $`a(t,y)`$ are given by
$`n(t,y)`$ $`=`$ $`e^{ky+kR(t)},`$ (23)
$`a(t,y)`$ $`=`$ $`n(t,y)H(t,y)=e^{ky+kR(t)}H(t,y).`$ (24)
in the region $`0yR(t)`$. Using these equations, $`G_{05}`$ in the bulk reads
$$\left[\left(\frac{H^{}}{H}\right)k\right]\frac{\dot{b}}{b}=\left(\frac{\dot{H}^{}}{H}\right)+k\dot{R}\left(\frac{H^{}}{H}\right).$$
(25)
Now we take the gauge in which $`b(y,t)=1`$. Then Eq. (25) can be solved by
$$H(t,y)=e^{kR(t)}G(y).$$
(26)
for any $`G(y)`$. Inserting Eq. (26) into Einstein equations, the other Einstein equations can be rewritten as
$`G_{00}`$ $`=`$ $`3n^2\left[{\displaystyle \frac{G^{\prime \prime }}{G}}+\left({\displaystyle \frac{G^{}}{G}}\right)^24k{\displaystyle \frac{G^{}}{G}}+2k^2\right]=n^2\kappa ^2|\mathrm{\Lambda }|,`$ (27)
$`G_{ii}`$ $`=`$ $`n^2H^2\left[2{\displaystyle \frac{G^{\prime \prime }}{G}}+\left({\displaystyle \frac{G^{}}{G}}\right)^28k{\displaystyle \frac{G^{}}{G}}+6k^2\right]=n^2H^2\kappa ^2|\mathrm{\Lambda }|,`$ (28)
$`G_{55}`$ $`=`$ $`3\left[\left({\displaystyle \frac{G^{}}{G}}\right)^23k{\displaystyle \frac{G^{}}{G}}+2k^2\right]=\kappa ^2|\mathrm{\Lambda }|.`$ (29)
If we set $`G(y)=e^{mky}`$, only the case of $`m=0`$ satisfies all equations in Eq. (29) with following relation:
$$k^2=\kappa ^2|\mathrm{\Lambda }|/6.$$
(30)
Finally, the metric in the region $`0yR(t)`$ is
$$ds^2=e^{2ky+2kR(t)}(dt^2+e^{2kR(t)}\delta _{ij}dx^idx^j)+dy^2.$$
(31)
Due to $`𝐙_\mathrm{𝟐}`$ symmetry the metric in the region $`R(t)y0`$ is
$$ds^2=e^{2ky+2kR(t)}(dt^2+e^{2kR(t)}\delta _{ij}dx^idx^j)+dy^2.$$
(32)
Therefore the metric in $`R(t)yR(t)`$ can be written simply as
$$ds^2=e^{2k|y|}(e^{2kR(t)}dt^2+\delta _{ij}dx^idx^j)+dy^2,$$
(33)
which is a AdS metric in the bulk by redefining the time. At $`y=0`$, there is a static wall. By integrating the Einstein equation across the static wall domain wall, we get the relation of the tension and $`k`$:
$$\sigma =\frac{6k}{\kappa ^2}:=\sigma _c.$$
(34)
This, together with Eq. (30), gives the exact value of the tension of the static wall in terms of the bulk cosmological constant:
$$\sigma =\pm \sqrt{\frac{6|\mathrm{\Lambda }|}{\kappa ^2}}.$$
(35)
where $`\pm =k/|k|`$.
Note that the bulk Einstein equations do not determine $`R(t)`$, the position of the moving wall: for arbitrary function $`R(t)`$, (33) is a solution of the bulk Einstein equations. As is well known, $`R(t)`$ is determined by the Israel junction equation. In terms of the extrinsic curvature of the wall world-volume
$$K_{\mu \nu }=_\mu n_\nu ,$$
(36)
with $`n_\nu `$ the unit normal vector on the domain wall world-volume, the Junction condition reads
$$\mathrm{\Delta }K_{mn}=\kappa ^2\left(T_{mn}\frac{1}{3}T_{}^{l}{}_{l}{}^{}g_{mn}\right),$$
(37)
where $`\mathrm{\Delta }K_{mn}:=lim_{y+0}K_{mn}lim_{y0}=K_{}^{+}{}_{mn}{}^{}K_{}^{}{}_{mn}{}^{}`$. Here $`g_{mn}`$ is a metric on the wall and $`T_{mn}`$ is a energy-momentum tensor of the wall. If we impose $`𝐙_\mathrm{𝟐}`$ symmetry, $`\mathrm{\Delta }K_{mn}`$ is equal to $`2K_{}^{+}{}_{mn}{}^{}`$ due to $`K_{}^{+}{}_{mn}{}^{}=K_{}^{}{}_{mn}{}^{}`$. For ‘extremal’ wall, the case where $`T_{mn}=\sigma g_{mn}`$, the Israel junction equations become
$$2K_{}^{+}{}_{mn}{}^{}=\frac{\kappa ^2}{3}\sigma g_{mn}.$$
(38)
To be explicit, let’s write the metric as
$$ds^2=d\tau ^2=f(t,y)dt^2+g(y)\delta _{ij}dx^idx^j+dy^2,$$
(39)
and let $`u^\mu `$ be the tangent velocity vector of the wall satisfying $`u^\mu u_\mu =1`$. Because the wall moves in $`y`$ direction, $`u^\mu `$ is given by
$$u^\mu =(\sqrt{\frac{1+(_\tau y)^2}{f}},0,0,0,_\tau y).$$
(40)
Since the unit normal vector $`n^\mu `$ satisfies $`n^\mu u_\mu =0`$,
$$n^\mu =(\frac{_\tau y}{\sqrt{f}},0,0,0,\sqrt{1+(_\tau y)^2}).$$
(41)
Insert $`n_\mu `$ into Eq. (36), we obtain two junction equations: one is
$$2K_{}^{+}{}_{ii}{}^{}=\sqrt{1+(_\tau y)^2}_yg=\frac{\kappa ^2}{3}\sigma g,$$
(42)
and the other $`K_{00}`$ equation does not give an independent one. In our case, $`f(t,y)`$ and $`g(t,y)`$ are given by $`f(t,y)=e^{2ky+2kR(t)}`$ and $`g(t,y)=e^{2ky}`$. The $`ii`$ component of Israel junction equations is reduced to
$`\sqrt{1+(_\tau R(t))^2}={\displaystyle \frac{\sigma }{\sigma _c}}.`$ (43)
Notice that if $`k`$ is positive, the tension of the moving wall must be negative and $`R(t)`$ is determined as
$$R(t)=v\tau ,$$
(44)
where $`v=\pm \sqrt{|\sigma /\sigma _c|^21}`$. Using $`u^0=dt/d\tau =\sqrt{\frac{1+(_\tau y)^2}{f}}`$,
$$\tau (t)=\frac{1}{\sqrt{1+v^2}}t.$$
(45)
So $`R(t)`$ is given by
$$R(t)=\frac{v}{\sqrt{1+v^2}}t.$$
(46)
From the bulk metric, the four-dimensional metric on the static wall at $`y=0`$ can be reduced to
$$ds_{4}^{}{}_{}{}^{2}=d\stackrel{~}{t}^2+\delta _{ij}dx^idx^j,$$
(47)
where $`d\stackrel{~}{t}=e^{kR(t)}dt`$. This metric is a Mikowskian one as is the Randall-Sundrum case. The metric on the moving wall at $`y=R(t)`$ is
$$ds_{4}^{}{}_{}{}^{2}=d\tau ^2+e^{2kv\tau }\delta _{ij}dx^idx^j.$$
(48)
For $`k>0`$, the static wall has a positive tension at $`y=0`$ and the moving wall has a negative one at $`y=R(t)`$. For $`k<0`$, the situation is opposite. If $`kv>0`$, the scale factor on the moving wall is exponentially contracted as $`\tau `$ runs. This is because the wall is moving toward the inside the AdS space, where the scale factor decreases. If $`kv<0`$, the wall is moving toward the boundary: the scale factor increases and this is described as an inflation by the the observer at the domain wall.
When we consider the region $`|y|R`$, we have similar situation but with just one wall at $`y=R(t)`$. Exactly parallel discussion can be done. however this is the case which is already discussed in ref.. So we so not treat it here.
### B Moving domain walls with all positive tension
In the previous case, we imposed $`𝐙_\mathrm{𝟐}`$ symmetry so that we glued two vacuum AdS solutions bounded by one static and one moving brane whose velocity is determined by its tension. We also noticed that one of the brane tension must be negative. If we want to have solution with two or more branes having the same sign, we can not in general impose $`𝐙_\mathrm{𝟐}`$ symmetry and at least one side of the wall must be AdS-Schwarzschild solution (AdSS). The situation is just like the case of spherical shells in the Minkowski space: Inside the shell, it is vacuum solution while it must be Schwarzschild solution outside. The motion of each shell can be is described by the Israel junction condition. More explicitly, Let there be N walls at $`y=y_i(t)`$ $`i=1,2,\mathrm{},N`$. Now suppose that the metric in the region $`y_i<y<y_{i+1}`$ is described by the AdSS:
$$ds^2=e^{2ky}\left[(1\mu _ie^{4ky})dt^2+d\stackrel{}{x}^2\right]+\frac{dy^2}{1\mu _ie^{4ky}}.$$
(49)
Then the equation of motion for the wall is given by
$$\dot{y}_i^2+1=\left(\frac{\sigma }{\sigma _c}\right)^2+\frac{1}{4}(\mu _i+\mu _{i1})e^{4ky_i}+\frac{1}{32}\left(\frac{\sigma _c}{\sigma }\right)^2(\mu _i\mu _{i1})^2e^{8ky_i}.$$
(50)
What is important for us at this moment is just the existence of such solution describing the multi-domain walls with all positive tension.
### C Cosmology with Moving Domain wall
Now let’s try to realize a situation where there is a static domain wall at $`y=y_1=0`$ and there is a moving brane at $`y=y_2(t)>0`$. In the leftmost region $`y<0`$ we have a Ads vacuum. In between the two walls, we have a AdSS solution with non-extremal parameter $`\mu _1`$. In the right most region $`y>y_2(t)`$, we have another AdSS solution with parameter $`\mu _2`$. The motion for the first wall is described by
$$\dot{y}_1^2+1=\left(\frac{\sigma }{\sigma _c}\right)^2+\frac{1}{4}(\mu _1)e^{4ky_1}+\frac{1}{32}\left(\frac{\sigma _c}{\sigma }\right)^2(\mu _1)^2e^{8ky_1}.$$
(51)
If we want to have a static wall at $`y=0`$,
$$\sigma =\sigma _c\sqrt{\frac{1}{2}(1\mu _1/4)+\frac{1}{2}\sqrt{1\frac{5}{8}\mu _1+\frac{1}{16}\mu _1^2}}$$
(52)
This is the condition for the static wall.
Another fact one should notice is that in the absence of the orbifold symmetry, the effective 4 dimensional cosmological constant on the brane is not well defined. Decomposing the 5 dimensional Einstein equation to the paralell and vertical to the brane, we get 4 dimensional equation which contains term proportional to the induced metric to be identified as the 4 dimensional cosmological constant. It is a sum of bulk cosmological constant and a quadratic form of the extrinsic curvature:
$$\mathrm{\Lambda }_4=\frac{\kappa ^2}{2}\mathrm{\Lambda }_5+\frac{1}{2}((K_\mu ^\mu )^2K^{\mu \nu }K_{\mu \nu })+\frac{1}{4}(K_m^\sigma K_{n\sigma }g^{mn}K_\mu ^\mu K_{mn}g^{mn}),$$
(53)
where $`K`$ is the extrinsic curvature tensor and greek indices are 5 dimensional indices while roman ones are 4 dimensional. In the presence of the orbifold symmetry, the extrinsic curvatures across the brane, $`K_{\mu \nu }^+`$ and $`K_{\mu \nu }^{}`$, differ only in sign, which is indifferent in the quadratic form. In the absence of the symmetry, there is no way to define the effective 4 dimensional constant naturally. The best thing one can do is do define by hand such that the stationary brane has zero effective cosmological constant. We will not push this idea further here. The motion of the second wall is described by (50) with two parameter $`\mu _1,\mu _2`$.
What happens to the stationary brane(our universe) if the other brane comes and overlaps? Even in the absence of the interaction between the branes, tensions will be added therefore we expect that the stationary brane must move and inflate. If both branes are infinitely thin, the impact is instantaneous and the consequence is expected to be small if we neglect the interaction of the two branes. If the impact brane is thick, then overlapping will change the effective bulk cosmological constant near the thin brane therefore cause a motion of the stationary brane. (Remember that the static condition is the balance of the bulk cosmological constant and the brane tension both with and without the orbifold symmetry.) During the scattering, the thin brane moves hence inflates. After the thick brane pass away, the condition of the static brane is restored and it must stop moving. Therefore the process of thin and thick brane scattering gives a mechanism of inflation for finite time. In order words, there is a natural graceful exit out of the inflation. Furthermore if the moving brane is thick enough, the scattering process is smooth enough so that the density fluctuation caused by the inflation may fit the present observation.
## IV Conclusion
We considered moving-brane-solutions in AdS type background. In the first Randall-Sundrum configuration, there are two branes at fixed points of the orbifold symmetry. We proved that if one brane is fixed and the other brane is moving, the configuration is still a solution provided that the motion of the brane has a specific velocity determined by its tension and the bulk cosmological constant. In the absence of the $`𝐙_\mathrm{𝟐}`$ symmetry, we could construct multi-brane configurations patching AdS-Schwarzshild solutions. In this case, effective four dimensional cosmological constant on the brane is ambiguous. Instead, we found a condition for a brane to be stationary. Finally, we suggest a scenario where we may have inflation on the brane universe for finite duration of time, i.e, a graceful exit of inflation. Our description of the last point was very qualitative and we did not study the physics of the RS1 due to the brane motion. We wish to comeback to treatment of these aspects in future publication.
Acknowledgments
SJS want to thank S. Nam and P. Yi for discussion. He also want to express thank KIAS for its hospitality during his visit. This work has been supported by KOSEF 1999-2-112-003-5 and BK21.
Note added
After typing of this manuscript was finished, we saw the appearance of the paper hep-th/0004206 by Horowitz, Low and Zee, where a wave solution were constructed in a similar fashion.
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# Indecomposable Higher Chow cycles on Jacobians
## 0. Introduction
The aim of this paper is to construct some natural cycles in the higher Chow groups of Jacobians of smooth projective curves. Most of the paper is devoted to the first higher Chow groups $`CH^k(J(C),1)`$, especially the case $`k=g=genus(C)`$, but in the last section we also construct elements of $`CH^k(J(C),n)`$, $`n>1`$, for curves of low genus.
For a smooth projective variety $`X`$, the subgroup of $`CH^k(X,1)`$ of decomposable cycles, $`CH_{\mathrm{dec}}^k(X,1)`$, is defined to be the image of $`CH^{k1}(X)^{}`$, where we use the isomorphism $`CH^1(X,1)^{}`$. We let $`CH_{\mathrm{ind}}^k(X,1):=CH^k(X,1)/CH_{\mathrm{dec}}^k(X,1)`$ be the quotient group of indecomposable cycles. In , a natural element $`K`$ was constructed in $`CH^g(J(C),1)`$, $`C`$ hyperelliptic, and using the regulator map $`K`$ was shown to give a non-torsion element of $`CH_{\mathrm{ind}}^g(J(C),1)`$ when $`C`$ is generic hyperelliptic. The Pontryagin product of $`K`$ with zero cycles of degree zero, gives elements of $`CH^g(J(C),1)`$ which lie in the kernel of the regulator map. Our first result, Theorem 1.2, shows that such cycles give uncountably many elements of $`CH_{\mathrm{ind}}^g(J(C),1)_{}`$ if $`g3`$. We prove a more precise statement using the decomposition of the higher Chow groups due to Beauville and Deninger-Murre . Applying the motivic hard Lefschetz theorem of Künnemann we obtain uncountably many elements of $`CH_{\mathrm{ind}}^k(J(C),1)_{}`$, $`3kg`$, lying in the kernel of the regulator map. The main technical tool used in the proof is a Hodge theoretic criterion due to J. Lewis . This has been used earlier by Gordon and Lewis to construct indecomposable cycles with similar properties in products of generic elliptic curves.
Our main goal is to construct indecomposable cycles in $`CH^g(J(C),1)`$ for more general curves. A result of Nori \[14, 7.5\] implies that upto torsion the regulator image of $`CH^g(J(C),1)`$ is the same as that of $`CH_{\mathrm{dec}}^g(J(C),1)`$ for a generic curve of genus $`g=3`$, and it seems likely that this should also be true for higher $`g`$. This makes it difficult to use Lewis’ criterion to prove indecomposability; instead we employ a specialization argument. To this end, we first prove that the specialization of a decomposable cycle is decomposable (Theorem 2.1). The difficulty here is that a cycle on a family of smooth projective varieties which restricts to a decomposable cycle on a generic fibre need not be decomposable on the entire family. We circumvent this problem by considering an auxiliary family which is constructed by gluing a product family along a special fibre.
Given a divisor $`D=(a_1+a_2)(b_1+b_2)`$ on a curve $`C`$ of genus $`g2`$ with $`2[D]=0`$ in $`J(C)`$, we construct an element $`Z_D`$ of $`CH^g(J(C),1)`$, called the $`4`$-configuration since it is supported on $`4`$ copies of $`C`$ in $`J(C)`$. When $`C`$ and $`D`$ are generic and $`g3`$, we show (Theorem 3.1 that this element is indecomposable: By using a suitable Hurwitz scheme, we show that we may specialize $`C`$ to a hyperelliptic curve $`C^{}`$ in such a way that the $`4`$-configuration specializes to $`KK_t`$, with $`t`$ a generic point of $`C^{}`$. Indecomposability follows by applying Theorem 1.2 and 2.1. We then deduce that $`CH_{\mathrm{ind}}^g(J(C),1)_{}`$ is uncountable for generic curves of genus $`3`$ and $`4`$.
To construct elements of $`CH^g(J(C),1)`$ for arbitrary curves, we consider divisors $`D=_{i=1}^na_i_{i=1}^nb_i`$ with $`2[D]=0`$ in $`J(C)`$ and $`n>0`$. Associated to such a divisor there is a natural subspace of $`CH^{n+1}(C^{n+1},1)_{}`$ and we describe a method, generalizing that used for the $`4`$-configuration, which we believe should enable one to prove indecomposability of the general such element with $`C`$ and $`D`$ also generic and $`gn+1`$. Unfortunately, due to the combinatorial complexity of the cycles involved, we have not been able to complete the proof. However, we have checked indecomposability using a simple computer program for $`n6`$. In particular, we see that $`CH_{\mathrm{ind}}^g(J(C),1)_{}`$ is uncountable for a generic curve with $`3genus(C)12`$.
We conclude the paper with a construction of elements in $`CH^3(J(C),4g)`$, with $`g2`$. These may be viewed as successive degenerations of the $`4`$-configuration on a genus $`3`$ curve, as the curves acquire nodes. We expect, but do not prove, that these elements are indecomposable in a sense stronger than that of Lewis. We do prove however, using Lewis’s criterion, that for $`B`$ a bielliptic genus $`2`$ curve $`CH^3(J(B),2)/Im(K_2()CH^1(J(C)))`$ is non-trivial modulo torsion.
###### Conventions.
All varieties will be over the complex numbers $``$ and all points will be closed points. We shall say that a condition holds for a *general* point of a variety if it holds for all points in a Zariski open subset and it holds for a *generic* point if it holds for all points outside a countable union of proper subvarieties.
We denote by $``$ the Pontryagin product on the higher Chow groups of an abelian variety.
## 1. Lewis’ conditions, Pontryagin products and hyperelliptic Jacobians
The first higher Chow group $`CH^k(X,1)H^{k1}(X,𝒦_k)`$ of a non singular variety $`X`$ is generated by higher cycles of the form $`Z=_iZ_if_i`$, where the $`Z_i`$ are irreducible subvarieties of codimension $`(k1)`$ and the rational functions $`f_ik(Z_i)^\times `$ obey the rule $`_idiv(f_i)=0`$ as a cycle on $`X`$. Consider the subgroup of *decomposable* cycles
(1.1)
$$CH_{\mathrm{dec}}^k(X,1):=Im\left\{CH^1(X,1)CH^{k1}(X)CH^k(X,1)\right\},$$
and the related quotient of *indecomposable* cycles
$$CH_{\mathrm{ind}}^k(X,1):=CH^k(X,1)/CH_{\mathrm{dec}}^k(X,1).$$
Recall that if $`X`$ is projective, then $`CH^1(X,1)=^{}`$.
Let $`I`$ be the zero cycles of degree zero on an abelian variety $`A`$. Bloch has shown that $`I^n`$ is non-zero for $`1ng`$ whereas $`I^{(g+1)}`$ is always zero, $`g=dim(A)`$. Given $`Z`$ an indecomposable element in $`CH^k(A,1)`$, a natural question to ask is whether $`I^nZ`$ contains indecomposable cycles for $`n1`$. We remark that $`I^nZ`$ is in the kernel of the regulator map for any $`Z`$, since translation acts trivially on the Deligne cohomology $`H_𝒟^{2k1}(X,(k))`$.
Using as a basic tool a condition of Hodge type due to Lewis, for $`Z`$ a real regulator indecomposable element of $`CH^g(J(C),1)`$ we give a criterion in terms of the primitive cohomology of $`J(C)`$ for $`I^nZ`$ to contain indecomposable cycles. A significant instance of this situation is the case of generic hyperelliptic Jacobians, where we can use as $`Z`$ the basic cycle in $`CH_{\mathrm{ind}}^g(J(C),1)`$ found in . We show that $`I^nZ`$ contains indecomposable cycles for $`1ng2`$, whereas all elements of $`I^{(g1)}Z`$ are decomposable (with $``$ coefficients).
### 1.1. Preliminaries
We recall some notation and definitions and then we state a theorem of Lewis. Our aim is to have a concrete reference at hand, for more details the reader should consult either the original paper or the survey .
For $`X`$ projective and nonsingular,
$$\frac{H^{i1}(X,)}{F^jH^{i1}(X,)+H^{i1}(X,(j))}\frac{H^{i1}(X,(j1))}{\pi _{j1}(F^jH^{i1}(X,))}$$
and therefore one has the identification
$`H_𝒟^{2k1}(X,(k))`$ $`H^{2k2}(X,(k1))F^{k1}H^{2k2}(X,)`$
$`=:H^{k1,k1}(X,(k1)).`$
According to Beilinson , the real regulator image of a cycle $`Z=Z_if_i`$ in $`CH^k(X,1)`$ is the element
$$R_{k,1}(Z)H_𝒟^{2k1}(X,(k))H^{k1,k1}(X,(k1)),$$
determined by the class of the current
$$R_{k,1}(Z):\omega (2\pi \sqrt{1})^{k1d}\underset{i}{}_{Z_iZ_i^{\mathrm{sing}}}\omega \mathrm{log}|f_i|.$$
###### Definition.
A higher Chow cycle $`ZCH^k(X,1)`$ is said to be (real) regulator indecomposable if there exists a differential form
$$\omega (Hdg^{k1}(X))^{}H^{dk+1,dk+1}(X,(dk+1))$$
such that the pairing $`[R_{k,1}(Z),\omega ]0`$.
###### Lemma 1.1.
If $`ZCH^k(X,1)`$ is regulator indecomposable then $`Z`$ is indecomposable.
Recall that the coniveau filtration on $`H^i(X,)`$ is
$$N^jH^i(X,):=\mathrm{ker}\left(H^i(X,)lim_{codim{}_{X}{}^{}Yj}H^i(XY,)\right),$$
where the direct limit is over closed subvarieties $`YX`$. The complex subspace generated by the Hodge projected image of the coniveau filtration is
$$H_N^{kl,km}(X):=Im(N^{kl}H^{2klm}(X,)H^{kl,km}(X)).$$
Lewis constructs certain complex subspaces
$$H^{\{k,l,m\}}(X)H^{kl,km}(X),$$
such that for $`m=0`$ one has
(1.2)
$$H^{\{k,l,0\}}(X)H_N^{kl,k}(X).$$
The spaces $`H^{\{k,l,m\}}(X)`$ are obtained by a process of Künneth projection of the Hodge components of the real regulator classes of the elements in $`CH^k(X\times S,m)`$, with varying smooth projective varieties $`S`$.
We give an abridged version of the main result from , since we only need it for $`CH_{\mathrm{ind}}^k(X,1)`$.
###### Theorem 1.1 (Lewis).
Let $`X`$ be a non singular projective variety:
1. $`H^{\{k1,l1,0\}}(X)H^{\{k,l,1\}}(X)`$.
2. If $`H^{\{k,l,1\}}(X)/H^{\{k1,l1,0\}}(X)0`$ for some $`l`$, $`2lk`$, then $`CH_{\mathrm{ind}}^k(X,1)_{}`$ is uncountable.
Note that Lewis’ proof shows that for a cycle $`\xi CH^k(X\times S,1)`$ which gives rise to a nonzero element of $`H^{\{k,l,1\}}(X)/H^{\{k1,l1,0\}}(X)`$, the cycles $`\xi _s`$, $`sS`$, form an uncountable subset of $`CH_{\mathrm{ind}}^k(X,1)_{}`$.
Our results will show that (ii) holds for a generic hyperelliptic Jacobian of genus $`g3`$, for $`k=g`$ and all $`l`$ such that $`2lg1`$.
### 1.2. Lewis’ condition holds for the Pontryagin families $`Z(m)`$
Let $`C`$ be a smooth projective curve of genus $`g`$ and $`J=J(C)`$ be its Jacobian. Let $`ZCH^g(J(C),1)`$ be a regulator indecomposable cycle. Given $`m`$ points $`p_i`$ on $`C`$, we construct higher cycles $`Z(m)`$ in $`CH^g(C^m\times J(C),1)`$, so that the fibre over $`(t_1,\mathrm{},t_m)`$ is $`Z([t_1p_1]e)\mathrm{}([t_mp_m]e)`$, where $`e`$ is the origin in $`J(C)`$. Set $`Z(1):=b_{}(C\times Z)C\times Z`$, where $`b:C\times J(C)C\times J(C)`$ is the twisted isomorphism $`b(t,x)=(t,[tp_1]+x)`$. By iteration this gives $`Z(m)`$ on $`C^m\times J(C)=C\times (C^{m1}\times J(C))`$, where the twisted map is now $`b((t_m,\mathrm{},t_1,x))=(t_m,\mathrm{},t_1,[t_mp_m]+x)`$.
###### Proposition 1.1.
Let $`C`$ be a smooth projective curve of genus $`g`$ and $`Z`$ a regulator indecomposable cycle in $`CH^g(J(C),1)`$. Assume that the primitive cohomology $`P^i(J(C))`$ is an irreducible sub-Hodge structure of $`H^i(J(C),)`$ for some $`mg2`$ and all $`im+2`$. Then the real regulator of $`Z(m)`$ gives rise to a non-zero element of $`H^{\{g,m+1,1\}}(J(C))/H^{\{g1,l,0\}}(J(C))`$, and hence the restriction of $`Z(m)`$ to the fibre over a generic point of $`C^m`$ is indecomposable.
If $`\omega _1^J,\mathrm{},\omega _g^J`$ is a basis for $`H^0(J(C),\mathrm{\Omega }_J^1(C))`$ such that the restriction $`\zeta _i:=\omega _{i|C}^J`$, $`1ig`$, is an orthonormal frame for $`H^0(C,\mathrm{\Omega }_C^1)`$, then the class of the divisor $`\mathrm{\Theta }`$ is determined by the form $`\theta _J(C)=(i/2)_{j=1}^g\omega _j^J\overline{\omega }_j^J`$. We define $`\tau _C=`$ $`\omega _1^J\overline{\omega }_1^J\omega _2^J\overline{\omega }_2^J`$.
###### Lemma 1.2.
Under the assumptions of Proposition 1.1, if Z is real regulator indecomposable then there is a basis as above with $`[R(Z),\tau ]0`$.
###### Proof.
The inner product on $`H^0(C,\mathrm{\Omega }_C^1)`$ allows us to view it as a representation $`V`$ of the unitary group $`U(g)`$. $`H^{1,1}(J(C))`$ is then also a representation of $`U(g)`$, and is isomorphic to a twist of $`VV^{}`$ by a 1-dimensional representation. Hence by Pieri’s formula we see that it decomposes as a direct sum of two irreducible representations $`T`$ and $`U`$, $`T`$ being the subspace corresponding to the class of the $`\mathrm{\Theta }`$ divisor.
Now consider the subspace $`W`$ of $`H^{1,1}(J(C))`$ spanned by all possible $`\tau `$’s as above. Since the unitary group preserves the inner product, it follows that $`W`$ is a subrepresentation of $`H^{1,1}(J(C))`$. Clearly $`W`$ is not equal to $`T`$, hence it must contain $`U`$. If there were no $`\tau `$’s with a non-zero pairing with $`Z`$ then the pairing would be zero on all of $`W`$, contradicting the assumption of real regulator indecomposable. (Note that our assumptions imply that the space spanned by the rational Hodge classes is $`T`$).
Let $`\alpha _m=\overline{\zeta }_{m+2}\mathrm{}\overline{\zeta }_3\tau \omega _3^J\mathrm{}\omega _{m+2}^J`$.
###### Lemma 1.3.
$`[R(Z(m)),\alpha _m]=(1)^m[R(Z),\tau ]`$
###### Proof.
By iteration the proof is the same for all $`m1`$. Say $`m=1`$, then we have $`b^{}(\omega _s^J)=\omega _s^J+\zeta _s`$ and $`b^{}(\zeta _s)=\zeta _s`$, thus
$$b^{}(\overline{\zeta }_3\tau \omega _3^J)=\overline{\zeta }_3b^{}(\tau \omega _3^J)=(\zeta _3\overline{\zeta }_3)\tau +\underset{i}{}\varphi _i^C\psi _i^J.$$
Here $`\varphi _i^C`$ is a form from C and $`\psi _i^J`$ is from $`J(C)`$, and it is either $`\varphi _i^C=\overline{\zeta }_3`$ or else $`\varphi _i^C=\overline{\zeta }_3\zeta _l,l3`$, and therefore it is of volume $`0`$, because of the orthogonality assumption. We have then : $`[R(Z(1)),\overline{\zeta }_3\tau \omega _3^J]=`$ $`[R(C\times Z),b^{}(\overline{\zeta }_3\tau \omega _3^J)][R(C\times Z),\overline{\zeta }_3\tau \omega _3^J]=`$ $`[R(C\times Z),b^{}(\overline{\zeta }_3\tau \omega _3^J)]=[R(Z),\tau ]`$, because $`[R(C\times Z),\varphi ^C\psi ^J]=`$ $`[R(Z),\psi ^J]_C\varphi ^C`$. ∎
###### Proof of Proposition 1.1.
The previous lemma along with equation (1.2) implies that $`R(Z(m))`$ gives a non-zero element of $`H^{\{g,m+1,1\}}(J(C))/H^{\{g1,m,0\}}(J(C))`$ if $`\tau \omega _3^J\mathrm{}\omega _{m+2}^J`$ is orthogonal to $`N^{gm1}H^{2gm2}(J(C))`$. The assumptions on the primitive cohomology imply that $`N^{gm1}H^{2gm2}(J(C))=\mathrm{\Theta }^{gm1}H^m(J(C))`$. Since $`\tau \omega _3^J\mathrm{}\omega _{m+2}^J\mathrm{\Theta }^{gm1}=0`$, it follows that the condition is indeed satisfied. ∎
The next proposition shows that the hypothesis on the primitive cohomology of Proposition 1.1 holds for the Jacobian of a generic hyperelliptic curve. The reader may refer to for the definition and basic properties of the Hodge group.
###### Proposition 1.2.
Let $`C`$ be a generic hyperelliptic curve of genus $`g`$. Then the Hodge group of $`J(C)`$ is isomorphic to $`Sp(2g,)`$, hence $`P^i(J(C))`$ is an irreducible Hodge structure for $`0ig`$.
###### Proof.
We shall use induction on $`g`$, the result for $`g=1`$ and $`2`$ being well known. Assume that $`g3`$ and the result is known for smaller $`g`$. We degenerate $`C`$ to a stable curve $`C_o`$ with $`3`$ smooth irreducible components $`C_1,C_2`$ and $`C_3`$, with $`C_1`$ and $`C_3`$ of genus $`1`$, and $`C_2`$ of genus $`g2`$ intersecting each of $`C_1`$ and $`C_3`$ transversally in a single Weierstrass point. By choosing a path in a suitable parameter space, we can identify $`V:=H^1(C,)`$ as a symplectic vector space with $`V_1V_2V_3`$, where $`V_i=H^1(C_i,)`$, $`i=1,2,3`$. Let $`D_1`$ and $`D_3`$ be generic hyperelliptic curves of genus $`g1`$, with $`D_i`$ specializing to $`C_iC_2`$, $`i=1,2`$. Let $`C^{}=D_1C_3`$ and $`C^{\prime \prime }=C_1D_3`$, the two components of each curve intersecting transversally in a single point. Again, by choosing paths we may identify $`H^1(C^{},)`$ and $`H^1(C^{\prime \prime },)`$ with $`V_1V_2V_3`$ in such a way that $`H^1(D_i,)`$ is identified with $`V_iV_2`$.
Now consider the family of Jacobians. Since $`C`$ is generic, $`G`$, the Hodge group of $`J(C)`$, contains the Hodge groups $`G^{}`$ and $`G^{\prime \prime }`$ of $`J(C^{})`$ and $`J(C^{\prime \prime })`$ respectively. Using induction and the above identifications, we see that $`G`$ contains both $`Sp(V_1V_2)`$ and $`Sp(V_2V_3)`$. One easily checks, by an explicit computation using Lie algebras, that the smallest subgroup of $`GL(V)`$ containing both these two subgroups is $`Sp(V)`$. The Hodge group is always contained in $`Sp(V)`$ so $`G`$ must equal $`Sp(V)`$.
From the representation theory of symplectic groups it follows that $`P^i(J(C))`$ is an irreducible representation of the Hodge group, $`0ig`$. Since the sub-Hodge structures of $`H^i`$ are precisely the subrepresentations of the Hodge group, it follows that the $`P^i(J(C))`$’s are also irreducible as Hodge structures. ∎
### 1.3. On the real regulator image of the basic hyperelliptic cycle
Let $`f:C^1`$ be the double cover associated with a hyperelliptic curve. We fix two ramification points $`w_1`$ and $`w_2`$ on $`C`$ and choose a standard parameterization on $`^1`$ so that $`f(w_1)=0`$ and $`f(w_2)=\mathrm{}`$. The points $`w_1`$ and $`w_2`$ are the distinguished Weierstrass points, and $`ϵ:=[w_1w_2]`$ is the associated element of order two in $`Pic^0(C)`$.
It is convenient to identify $`J(C)=Pic^0(C)`$ with $`Pic^1(C)`$ by adding $`w_1`$. We embed $`C`$ in the natural way in $`Pic^1(C)`$, and for $`tPic^1(C)`$ we let $`C_t`$ be the translate of $`C`$ by $`[tw_1]`$. Now consider $`W_1:=C=C_{w_1}`$ and $`W_2:=C_{w_2}`$, the $`ϵ`$ translate of $`C`$, and fix a point $`tC`$. Observe that the intersection $`C_tW_1W_2`$ is the point $`w_1`$. We shall follow the convention to indicate a rational function on $`C_t`$ by using the same name given to the corresponding function on $`C`$.
Consider $`K:=W_1f+W_2f`$. $`K`$ is the basic hyperelliptic cycle of , where it was provedthat it is a non trivial indecomposable element of $`CH^g(J(C),1)_{}`$ for generic $`C`$. There it was shown that the primitive contribution of the standard regulator image of $`K`$ does not vanish by studying an infinitesimal invariant of Griffiths type associated with the relevant normal function. The following proposition shows that $`K`$ is real regulator indecomposable. Here $`\tau _C`$ is of the form considered in section 1.2.
###### Proposition 1.3.
For $`C`$ a generic hyperelliptic curve of genus $`g2`$ there is $`\tau _C`$ with $`[R(K),\tau _C]0.`$
By definition $`[R(K),\tau ]=2_Clog|f|\tau _{|C}`$; we shall prove that it is not trivial by means of a reduction process to the case of elliptic curves.
Let $`E_\lambda `$ be the elliptic curve with affine equation $`y^2=x(x1)(x\lambda )`$. Define $`f_\lambda :=x`$ as a rational function on $`E_\lambda `$. The form associated with the $`\mathrm{\Theta }`$ divisor is here $`\theta _\lambda =(i/2)\omega _\lambda \overline{\omega }_\lambda `$. It is an invariant form on $`E_\lambda `$ of volume one. We write $`I(\lambda ):=_{E_\lambda }\mathrm{log}|f_\lambda |\theta _\lambda `$.
###### Lemma 1.4.
$`I(\lambda )`$ varies with $`\lambda `$.
###### Proof.
Multiplication of $`x`$ by $`\lambda ^1`$ shows that $`E_\lambda `$ and $`E_{\lambda ^1}`$ are isomorphic models of the same curve $`E`$. On $`E`$ the volume forms coincide, while $`f_\lambda =\lambda f_{\lambda ^1}`$. Thus
$$I(\lambda )=_E\mathrm{log}|\lambda |\theta +I(\lambda ^1)=\mathrm{log}|\lambda |+I(\lambda ^1),$$
hence $`I(\lambda )`$ cannot be constant. ∎
We define $`I(h,\tau ):=_C\mathrm{log}|h|\tau _C`$, for any rational function $`h`$ on $`C`$.
###### Lemma 1.5.
If $`C`$ is a generic curve of genus $`2`$, then $`I(f,\tau )0`$.
###### Proof.
We prove it for a bielliptic curve $`C`$ which is a double cover of $`E_1:=E_{\lambda _1}`$ and of $`E_2:=E_{\lambda _2}`$. Consider the diagram
$$\begin{array}{ccccc}E_2& \stackrel{k_2}{}& C& \stackrel{k_1}{}& E_1\\ f_2& & f& & f_1& & \\ ^1& \stackrel{h}{}& ^1& \stackrel{h}{}& ^1\end{array}$$
Here $`f_i`$ is ramified over $`\{0,1,\mathrm{},\lambda _i\}`$, $`h`$ is the double cover ramified over $`\lambda _1`$ and $`\lambda _2`$, and $`f:C^1`$ is the hyperelliptic cover ramified at $`h^1(\{0,1,\mathrm{}\})`$. On the range of $`h`$ we have already fixed a standard parameter, we choose a standard parameter on the domain of $`h`$ so that $`0`$ maps to $`0`$, and similarly for $`1`$ and for $`\mathrm{}`$. In this manner $`f`$ is a well defined rational function on $`C`$, and we denote by $`\overline{f}`$ its transform under the involution of $`^1`$ associated with $`h`$. Letting $`g:=hf`$, we see that $`f\overline{f}=cg`$, $`c`$ a constant.
One can take $`\tau _C`$ to be the form $`(k_1^{}(\theta _1)k_2^{}(\theta _2))`$ on $`J(C)`$ and thus
$$I(g,\tau _C)=_C\mathrm{log}|g|(k_1^{}(\theta _1)k_2^{}(\theta _2))0.$$
It then follows that $`I(f,\tau _C)0`$ for the general bielliptic curve C because
$$I(f,\tau _C)+I(\overline{f},\tau _C)=I(g,\tau _C)+log|c|_C\tau _C=I(g,\tau _C).\mathit{}$$
###### Proof of Proposition 1.3.
The proof for arbitrary genus is obtained by induction. Starting from a hyperelliptic curve $`G^1`$ and a double cover $`h:^1^1`$ we construct the commutative diagram
$$\begin{array}{ccc}C& \stackrel{\pi }{}& G\\ f_C& & f_G& & \\ ^1& \stackrel{h}{}& ^1\end{array}$$
Here $`C`$ is the normalization of the Cartesian product, hence $`\pi `$ is branched at $`42m`$ points, where $`m`$ is the number of ramification points of $`h`$ which coincide with points of ramification for $`f_G`$. We have $`g(C)=2g(G)+1m`$.
By induction, Proposition 1.3 holds for $`G`$. It then also holds for $`C`$ by using the arguments given for Lemma 1.5, where we now take $`\tau _C`$ to be the form $`\pi ^{}\tau _G`$, the lift to $`J(C)`$ of $`\tau _G`$. ∎
### 1.4. Indecomposable elements on hyperelliptic Jacobians
For a $`g`$ dimensional abelian variety $`A`$, the following decomposition of the higher Chow groups is a consequence of the motivic decomposition of the diagonal due to Beauville , and Deninger and Murre :
(1.3)
$$CH^k(A,m)_{}=\underset{s}{}CH^k(A,m)_s$$
Here $`CH^k(A,m)_s`$ is the subspace of $`CH^k(A,m)_{}`$ on which $`[n]^{}`$ (resp. $`[n]_{}`$) acts by multiplication by $`n^{2ks}`$ (resp. $`n^{2g2k+s}`$). The Fourier transform of Mukai and Beauville induces isomorphisms:
(1.4)
$$:CH^k(A,m)_s\stackrel{}{}CH^{gk+s}(\widehat{A},m)_s$$
where $`\widehat{A}`$ is the dual abelian variety.
For $`\mathrm{\Theta }`$ a symmetric ample divisor, the motivic hard Lefschetz theorem of Künnemann implies that intersecting with powers of $`\mathrm{\Theta }`$ gives isomorphisms:
(1.5)
$$\mathrm{\Theta }^{g+s2k}:CH^k(𝒜,m)_s\stackrel{}{}CH^{g+sk}(𝒜,m)_s,02ksg.$$
It follows from the definitions that the decomposition (1.3) and the isomorphisms (1.4) and (1.5) preserve decomposable cycles as defined in (1.1), hence
(1.6)
$$CH_{\mathrm{ind}}^k(A,1)_{}=\underset{s}{}CH_{\mathrm{ind}}^k(A,1)_s$$
(1.7)
$$:CH_{\mathrm{ind}}^k(A,1)_s\stackrel{}{}CH_{\mathrm{ind}}^{gk+s}(A,1)_s$$
(1.8)
$$\mathrm{\Theta }^{g+s2k}:CH_{\mathrm{ind}}^k(𝒜,1)_s\stackrel{}{}CH_{\mathrm{ind}}^{g+sk}(𝒜,1)_s,02ksg.$$
###### Proposition 1.4.
Let $`A`$ be a g-dimensional abelian variety. Then $`CH_{\mathrm{ind}}^g(A,1)_s=0`$ for $`s<2`$ or $`s>g`$.
###### Proof.
We use equation (1.7) i.e. $`(CH_{\mathrm{ind}}^g(A,1)_s)=CH_{\mathrm{ind}}^s(\widehat{A},1)_s`$. If $`s<1`$, $`CH^s(\widehat{A},1)`$ is itself zero. The action of $`[n]^{}`$ on $`CH^1(\widehat{A},1)=𝐂^{}`$ is trivial hence $`CH^1(\widehat{A},1)_1`$ is also zero. We conclude the proof by observing that for any $`g`$-dimensional variety $`X`$, $`CH^s(X,1)=0`$ for $`s>g+1`$ and $`CH_{\mathrm{ind}}^{g+1}(X,1)=0`$
###### Remark.
A conjecture of C. Voisin says that $`CH_{\mathrm{ind}}^2(X,1)`$ should be countable for any smooth projective variety $`X`$. For an abelian variety $`A`$, the injectivity of the rational regulator on $`CH^k(A,1)_2`$ would imply that $`CH_{\mathrm{ind}}^k(A,1)_2`$ is countable. If $`g=2`$, the proposition shows that then $`CH_{\mathrm{ind}}^2(A,1)_{}`$ would also be countable.
For the rest of this section, $`C`$ will be a hyperelliptic curve of genus $`g`$ and $`K`$ the basic hyperelliptic cycle in $`CH^g(J(C),1)`$.
###### Lemma 1.6.
The component of $`K`$ in $`CH_{\mathrm{ind}}^g(J(C),1)_s`$ is zero for all $`s2`$. Consequently, for any integer $`n`$, $`[n]_{}(K)=n^2K`$ in $`CH_{\mathrm{ind}}^g(J(C),1)_{}`$.
###### Proof.
Consider the following copies of $`C`$ embedded in $`C\times C`$: $`X_1=C\times \{w_1\}`$, $`X_2=\{w_1\}\times C`$, $`X_3=C\times \{w_2\}`$, $`X_4=\{w_2\}\times C`$, $`X_5=\mathrm{\Delta }`$, and $`X_6=\mathrm{\Delta }^{}`$. Here $`w_1`$ and $`w_2`$ are two distinct Weierstrass points, $`\mathrm{\Delta }`$ is the diagonal, and $`\mathrm{\Delta }^{}`$ is the image of $`C`$ via the embedding $`x(x,\sigma (x))`$, with $`\sigma `$ the hyperelliptic involution. If $`f`$ is a Weierstrass function with $`div(f)=2w_12w_2`$ then one easily checks that $`Z=_{i=1}^6X_if`$ is an element of $`CH^2(C\times C,1)`$. If we use $`w_1`$ to embed $`C`$ in $`J(C)`$, then the image of $`Z`$ in $`CH^g(J(C),1)`$ is equal to $`2K(1/2)[2]_{}(K)`$.
The involution $`\tau =(id,\sigma )`$ of $`C\times C`$ preserves $`Z`$, hence $`Z`$ must be the pullback of a cycle from $`C\times C/\tau C\times ^1`$. $`CH^2(C\times ^1,1)`$ is always decomposable for any curve $`C`$, hence $`Z`$ must also be decomposable. Thus $`4K=[2]_{}(K)`$ in $`CH_{\mathrm{ind}}^g(J(C),1)_{}`$. This implies that all the components of $`K`$ except the one in $`CH_{\mathrm{ind}}^g(J(C),1)_2`$ must be zero. ∎
For $`T=(t_1,t_2,\mathrm{},t_m)`$, $`P=(p_1,p_2,\mathrm{},p_m)`$ points on $`C^m`$, let $`K_P(T)=K([t_1p_1]e)([t_2p_2]e)\mathrm{}([t_mp_m]e)`$ $`CH^g(J(C),1)`$, where $`e`$ is the origin in $`J(C)`$. The following theorem is the main result of this section.
###### Theorem 1.2.
If $`C`$ is a generic hyperelliptic curve and $`T`$ a generic point of $`C^m`$, $`1mg2`$, then $`K_P(T)`$ is indecomposable and the set of all such cycles with $`T`$ varying and $`P`$ fixed forms an uncountable subset of $`CH_{\mathrm{ind}}^g(J(C),1)_{}`$. Furthermore, $`CH_{\mathrm{ind}}^k(J(C),1)_s`$ is uncountable for $`3skg`$.
###### Proof.
The first part of the theorem follows directly by combining Propositions 1.1, 1.2 and 1.3.
Let $`I=_{s=1}^gCH^g(J(C))_s`$ be the zero cycles of degree 0 (with $``$ coefficients). Then for $`m1`$, $`I^m=_{s=m}^gCH^g(J(C))_s`$ and if $`TC^m`$, then $`K_P(T)CH^g(J(C),1)_{}I^m`$. Using Proposition 1.4, Lemma 1.6 and the fact that the subscripts are additive under Pontryagin products, we see that if $`m=g2`$ then the only nonzero component of $`K_P(T)`$ in $`CH_{\mathrm{ind}}^g(J(C),1)_{}`$ is the one in $`CH_{\mathrm{ind}}^g(J(C),1)_g`$. The first part of the theorem then implies that the image of $`KCH^g(J(C))_{g2}`$ in $`CH_{\mathrm{ind}}^g(J(C),1)_g`$ is uncountable. We then use that $`CH^g(J(C))_s=CH^g(J(C))_1^s`$, $`s1`$, and descending induction on $`s`$ starting from $`s=g`$ to show that the image of $`KCH^g(J(C))_{s2}`$ in $`CH_{\mathrm{ind}}^g(J(C),1)_s`$ is uncountable for $`3sg`$.
The statement for $`k<g`$ follows by using the Fourier transform and the motivic hard Lefschetz theorem: Letting $`k=s`$, we see by (1.7) that $`CH_{\mathrm{ind}}^s(J(C),1)_s`$ is uncountable for $`3sg`$. By (1.8), intersection with $`\mathrm{\Theta }^{gs}`$ induces an isomorphism from $`CH_{\mathrm{ind}}^s(J(C),1)_s`$ to $`CH_{\mathrm{ind}}^g(J(C),1)_s`$, hence intersection with $`\mathrm{\Theta }^{ks}`$ must induce an injection from $`CH_{\mathrm{ind}}^s(J(C),1)_s`$ to $`CH_{\mathrm{ind}}^k(J(C),1)_s`$, $`3skg`$. Since $`CH_{\mathrm{ind}}^s(J(C),1)_s`$ is uncountable, it follows that $`CH_{\mathrm{ind}}^k(J(C),1)_s`$ is also uncountable. ∎
###### Remark.
Proposition 1.4 shows that the theorem is optimal for $`k=g`$. It should also be optimal for $`k<g`$ — the proof for $`s>k`$ still works but we do not know how to handle the $`s<2`$ case.
## 2. Decomposability specializes
Our aim is to show that decomposability specializes. Consider a flat and projective family $`𝒳A`$, where $`A`$ is a smooth curve, $`𝒳`$ is non singular and so is $`X_0`$, the fibre over $`p_0`$.
###### Theorem 2.1.
If the restriction of an element $`𝒬CH^d(𝒳,1)_{}`$ to the generic fibre is decomposable then restriction of $`𝒬`$ to the central fibre is also decomposable.
### 2.1.
One can replace $`𝒬`$ by a multiple and $`A`$ by an open subset of a finite cover of the original $`A`$ so that now the following holds:
###### Assumption.
Over $`U:=A\{p_0\}`$ the restriction $`𝒬_U`$ is equivalent in $`CH^d(𝒳_U,1)`$ to an element $`𝒲:=𝒵_jf_j`$, where $`𝒵_j`$ are irreducible subvarieties which intersect properly $`X_0`$ and $`f_j`$ are rational functions lifted from $`A`$ and regular on $`U`$.
We meet now the problem that $`𝒲`$ may have a boundary $`B`$ on $`𝒳`$, and then $`B`$ is supported on $`X_0`$. On the other hand $`𝒬`$ is a cycle, and so it has no boundary. Since on $`𝒳_U`$ the classes $`𝒬`$ and $`𝒲`$ coincide by hypothesis, then under the boundary map $`CH^d(𝒳_U,1)CH^{d1}(X_0)`$ the image of $`𝒲`$ is trivial, that is $`B`$ is rationally equivalent to 0 on the central fibre. Let $`R`$ be a relation of rational equivalence on $`X_0`$ which kills $`B`$, then $`𝒲R`$ has no boundary and thus $`𝒲R𝒬`$ is a cycle for $`CH^d(𝒳,1)`$ whose class is represented by a cycle $`M`$ coming from $`CH^{d1}(X_0,1)`$. We will show that $`M`$ and $`R`$ can be moved conveniently for our purposes, but to do this we need to embed $`𝒳`$ in a larger space $`𝒴A`$. Our first result is that we can work on $`𝒴`$ and with a cycle $`𝒲^Y:=𝒵_jf_j`$ as described above. The improvement is that $`𝒲^Y`$ has no boundary now and that it is equivalent to $`𝒬`$ on $`𝒴`$. The second issue is to compute the restriction of $`𝒲^Y`$ to $`X_0`$ so to check that it is indeed of decomposable type. The difficulty that we meet is the fact that each $`f_j`$ may have a boundary on the central fibre, we overcome this trouble by means of some test curves.
### 2.2.
We find useful to work for a while with $`G(T)`$ the $`K`$-theory of coherent sheaves on a quasiprojective scheme $`T`$ and to use the topological filtration $`F_mG_1(T)`$, this is the subgroup generated by the images of $`i_{Z,}:G_1(Z)G_1(T)`$, where $`i_Z:ZT`$ are closed subschemes of $`T`$ of dimension at most $`m`$. The relations with Bloch’s $`CH^p(X,1)`$ are now recalled.
We know from Soulé that the Quillen coniveau spectral sequence
$$E_1^{p,q}=\underset{xX_p}{}K_{pq}k(x)G_{pq}(X)$$
degenerates modulo torsion for $`p+q2`$, in particular we need
$$G_1(X)_{}=\underset{p=0}{\overset{d}{}}E_2^{p,p1}$$
Moreover for $`m=0,1,2`$ the coniveau filtration of the spectral sequence coincides rationally with the $`\gamma `$ filtration. Using Bloch’s isomorphism
$$\underset{i}{}gr_\gamma ^iG_m(X)_{}G_m(X)_{}\underset{i}{}CH^i(X,m)_{}$$
one has then
$$E_2^{p,p1}CH^p(X,1)_{}$$
Let $`i_D:DT`$ be the inclusion of an effective Cartier divisor, then $`i_D^{}G_m(T)G_m(D)`$ is defined. Next proposition, see 2-11 in , implies that if each component of a subscheme $`Z`$ of dimension $`m`$ intersects properly $`D`$, then $`i_D^{}(i_Z(G_1(Z))`$ lands in $`F_{m1}G_1(D)`$.
###### Proposition 2.1 (Quillen).
Consider a cartesian diagram of quasiprojective schemes
$$\begin{array}{ccc}X^{}& \stackrel{g^{}}{}& X\\ f^{}& & f& & \\ Y^{}& \stackrel{g}{}& Y\end{array}$$
Assume that $`f`$ is proper, that $`g`$ is of finite tor-dimension and that $`Y^{}`$ and $`X`$ are tor-independent over $`Y`$ (i.e. $`Tor_i^{O_{Y,y}}(O_{(Y^{},y^{})},O_{(X,x)})=0`$ for $`i1`$)
Then
$$g^{}f_{}=f_{}^{}g_{}^{}{}_{}{}^{}:G_m(X)G_m(Y^{})$$
### 2.3. A larger space
We come back to consider $`𝒳A`$, by gluing it transversally with $`X\times A`$ along the fibre over $`p_0`$ we construct $`𝒴`$, so that $`g:𝒴A`$ has central fibre the scheme $`Y`$, whose reduced support is $`X_0`$.
The diagram
$$\begin{array}{ccccc}X& & 𝒳& & \\ & & & & & & \\ Y& & 𝒴& & :=𝒳_{X\times \{p_0\}}(X\times A)\end{array}$$
satisfies the hypotheses of Proposition 2.1. Indeed one has only to check vanishing of $`Tor_1`$, because $`Y`$ is a Cartier divisor in $`𝒴`$. Let $`t`$ be a local parameter at $`p_0`$, then $`t`$ generates the ideal of $`Y`$, and we see that $`Tor_1=0`$ because $`t`$ is not a 0-divisor in the local rings of $`𝒳`$.
Applying devissage we find then:
$$\begin{array}{ccc}G_q(X)& & G_q(𝒳)\\ & & & & \\ G_q(Y)& & G_q(𝒴)\end{array}$$
We use next the setting above to check that the specialization map $`G_1(𝒳)G_1(X)`$ is compatible with the topological filtration $`F_m`$. The useful fact is that on $`𝒴`$ it is easy to move to general position with respect to $`X`$. This way of moving turns out to be also the trick that we need to show that decomposability specializes.
We work by induction on $`m`$ and consider $`z`$ in the graded quotient $`F_{m+1}/F_mG_1(𝒳)`$, then $`z`$ is represented by a Quillen cycle $`Z_if_i`$, with $`div(f_i)=0`$ where $`Z_i`$ are irreducible subvarieties of dimension $`m+1`$. Now:
###### Lemma 2.1.
On $`X\times A`$ a Quillen chain $`Zf`$ is equivalent to a chain whose support intersects properly $`X\times \{p_0\}`$.
###### Proof.
The only problem is when $`ZX\times \{p_0\}`$. Consider the symbol $`\{g,f\}`$ on $`Z\times A`$, where $`g`$ is a rational function with simple 0 at $`p_0`$ and therefore $`div(g)=p_0+m_iq_i`$. The boundary of $`(Z\times A)\{g,f\}`$ is $`(Z\times \{p_0\})f+(Z\times \{q_i\})f^{m_i}(div(f)\times A)g`$, and thus $`(Z\times \{p_0\})f`$ is equivalent to a chain of the stated kind. ∎
In this way the image of $`z`$ in $`F_{m+1}/F_mG_1(𝒴)`$ is supported on a scheme $`T`$ of dimension $`m+1`$ with the property that each component of $`T`$ meets properly the central divisor $`X`$, therefore the restriction of $`z`$ belongs to $`F_mG_1(X)`$ because it is supported on the intersection $`TX`$.
The same proof yields that the chain $`M+R`$ which was discussed above is equivalent on $`𝒴`$ to a sum of Quillen chains of the preceding type $`(div(f)\times A)g`$. Notice that here we may have had to replace $`A`$ by a smaller open set, with this proviso we have then proved:
###### Proposition 2.2.
The restriction of $`Q`$ to $`X`$ coincides with the restriction from $`𝒴`$ to $`Y`$ of a cycle $`𝒲^Y:=𝒵_jf_j`$, where $`𝒵_j`$ are irreducible varieties in $`𝒴`$ and where $`f_j`$ is pull back of a rational function of the same name from $`A`$ under the flat projection $`𝒵_jA`$. The functions $`f_j`$ are regular away from $`p_0`$.
### 2.4. Computing the restriction of $`𝒲^Y`$
Let $`(_jZ_j)X_0=_lH_l`$ be the decomposition in irreducible components, then we know from the preceding discussion that the restriction morphism $`CH^d(𝒴,1)CH^d(X,1)`$ sends $`𝒲^Y`$ to a cycle supported on $`_lH_l`$, and therefore to a cycle $`_lH_lh_l`$. Next proposition completes the proof of the theorem.
###### Proposition 2.3.
The rational functions $`h_l`$ are constant functions.
###### Proof.
By taking general linear sections of $`𝒵=𝒵_j`$ one may construct curves $`B=B_j`$ such that $`BA`$ is finite. Our program is to prove that at the points of intersection of $`B`$ with $`H_l`$ the value of the relevant function $`h_l`$ is constant, independent of the chosen section $`B`$.
We may assume that $`B`$ is smooth outside the inverse image of $`p_0`$. Let $`Q`$ be the element of $`G_1(B)`$ determined by $`𝒲^Y`$, and thus $`Q`$ is such that its restriction to the smooth part is given by functions $`f_j`$ on the components $`B_j`$ of $`B`$, each $`f_j`$ being the pull-back of a rational function with the same name on $`A`$. Using devissage we restrict $`Q`$ to $`G_1(p)^{}`$, where $`p`$ is a point in the inverse image of $`p_0`$. Then the claim is that for any such point $`p`$ this restriction is equal to ($`f_{j}^{}{}_{}{}^{n_j})(p_0)`$. Here the $`n_j`$’s may depend on $`p`$ and they are such that the product doesn’t have a pole or zero at $`p_0`$. Note that contravariant functoriality for local complete intersections yields here the equality $`(f_{j}^{}{}_{}{}^{n_j})(p_0)`$ $`=`$ $`(h_l(p))`$, where we set $`h_l(p)=1`$ if $`pH_l`$. This argument shows that at the points of intersection of curves like $`B`$ with $`H_l`$ the values of the functions $`h_l`$ lie in a countable set, and therefore $`h_l`$ must be constant on the irreducible component $`H_l`$.
To prove the claim, we let $`C`$ be the normalization of an irreducible component of the fibre product of all the $`B_i`$’s over $`A`$. $`C`$ is flat over $`A`$, so we can replace $`C`$ by $`A`$ and $`B`$ by the fibre product of $`A`$ and $`B`$ over $`C`$. The advantage is that here we are reduced to the case that there is only one point in the inverse image of (the new) ’$`p`$’ because we now have sections, and so the claim follows by functoriality from the simple case, which we describe next.
The simple case is when each $`B_i`$ is an isomorphic copy of $`A`$ and moreover the scheme $`B:=B_i`$, is constructed by gluing at one point $`q`$ (= $`p_0`$ on $`A`$). Let the projection map be $`\pi :BA`$, then we write $`U:=A\{p_0\}`$, and $`V_i=B_i\{q\}`$, and set $`V:=V_i=B\{q\}`$. The map $`\pi :BA`$ is the identity on each component $`B_i`$. Give rational functions $`f_i`$ on $`B_i`$, ($`A`$ is not necessarily complete, and we assume that the only zero or pole of $`f_i`$ is at $`q`$) and assume that $`div(f_i)=0`$. Then $`\{f_i\}`$ defines an element $`f`$ say in $`G_1(V)`$, which comes from $`G_1(B)`$, because it has boundary 0 in the exact sequence $`G_1(B)G_1(V)G_0(q)`$. Consider now $`p_0`$ as a Cartier divisor on $`A`$ and let $`q^{}`$ be the scheme, which is the pull back of $`p_0`$ on $`B`$. The map $`G_1(B)G_1(q)`$ is well defined because $`q^{}`$ is a local complete intersection in $`B`$. By devissage $`G_1(q^{})=G_1(q)=^{}`$; the question is to understand the value of the image of $`f`$ in $`G_1(q)=^{}`$. Each $`f_i`$ is a function $`f_{iA}`$ on $`A`$. The answer is consider the product $`f_A:=\mathrm{\Pi }f_{iA}`$, then $`f_A`$ is a rational function which is in fact regular at $`p_0`$, because of our assumption, and then we have that the image value is $`f_A(p_0)`$. The reason is once more the commutativity from Proposition 2.1 (it is so below at $`p_0`$ and then it must have been so at $`q`$). ∎
## 3. The 4-configuration
For $`C`$ a generic hyperelliptic curve of genus $`g3`$, we have constructed indecomposable elements in $`CH^g(J(C),1)_{}`$ which are in the kernel of the regulator map. We shall now construct such elements in $`CH^g(J(C),1)_{}`$, for more general curves, in particular for generic curves of genus $`3`$ and $`4`$. The natural parameter space for our construction will be a certain Hurwitz scheme of degree $`4`$ covers of $`^1`$. We shall show that one can specialize to a hyperelliptic curve in such a way that our cycle specializes to $`KK_t`$, allowing us to use the results of the previous sections to prove indecomposability.
### 3.1. Construction of the 4-configuration
Let $`C`$ be a smooth projective curve of genus $`g3`$ and let $`D=(a_1+a_2)(b_1+b_2)`$ be a divisor on $`C`$ such that $`[D]=ϵ`$ is a point of order $`2`$ in $`Pic(C)`$. Let $`f`$ be a rational function on $`C`$ such that $`div(f)=2D`$. To this data we shall associate a natural element of $`CH^g(Pic^3(C),1)`$ supported on four copies of $`C`$.
Let $`i(y,z):CPic^3(C)`$ be the map $`i(y,z)(x)=x+y+z`$, $`C(y,z):=i(y,z)(C)`$, and let $`j:CPic^3(C)`$ be the map $`j(x)=x+2(a_1+a_2)`$, $`G:=j(C)`$. We consider $`C(a_1,a_2)`$, $`C(b_1,b_2)`$ and $`G`$. Translation by $`ϵ`$ maps $`C(a_1,a_2)`$ to $`C(b_1,b_1)`$, and we let $`G_ϵ`$ be the image of $`G`$. We shall use the convention that $`f`$ represents the rational function on each of the preceding curves which maps to $`f`$ under the chosen isomorphism with $`C`$. Thus, we set $`Z_1:=C(a_1,a_2)f`$, $`Z_2:=Gf`$, $`Z_3:=C(b_1,b_2)f`$, $`Z_4:=G_ϵf`$.
###### Proposition 3.1.
$`Z_D:=_{i=1}^4(1)^iZ_i`$ is a higher cycle.
Note that $`Z_D`$ also depends on the choice of the rational function $`f`$, but we can neglect this since multiplying $`f`$ by an element of $`^{}`$ amounts to the addition of a decomposable element to $`Z_D`$.
###### Proof.
The only possible difficulty is to see where the curves intersect. $`C(a_1,a_2)`$ intersects $`G`$ in two points; on both curves the points come from $`a_1`$ and $`a_2`$ under the isomorphism with $`C`$, but the point which comes from $`a_1`$ in $`G`$ comes from $`a_2`$ in $`C(a_1,a_2)`$ and conversely. A similar statement holds for $`C(b_1,b_2)G`$, and then intersections with $`G_ϵ`$ can be recovered by using $`ϵ`$-symmetry. Note that if $`C`$ is not hyperelliptic then $`C(a_1,a_2)C(b_1,b_2)=\mathrm{}`$ and $`GG_ϵ=\mathrm{}`$. ∎
###### Remark.
By translation of $`Z_D`$ by a a zero cycle $`\xi `$ of degree $`3`$ on $`C`$, we obtain an element $`Z_{D,\xi }`$ of $`CH^g(J(C),1)`$. The component of $`Z_{D,\xi }`$ in $`CH^g(J(C),1)_2`$ is always zero, hence it is in the kernel of the (rational) regulator map. To see this, first note that this component is preserved by translation since $`CH^g(J(C),1)_s=0`$ for $`s<2`$ (c.f. proof of Proposition 1.4), hence is independent of $`\xi `$. One easily checks that $`[1]^{}(Z_{D,3a_1})=Z_{D,3a_2}`$, This implies the desired fact, since the action of $`[1]^{}`$ on $`CH^g(J(C),1)_2`$ is trivial.
### 3.2. A Hurwitz scheme
In order to construct the specialization needed to prove indecomposability, we shall use certain Hurwitz schemes parametrizing the ramified coverings of $`^1`$ corresponding to the functions $`f`$ as above. We refer the reader to \[7, Section 1\] for the basic facts about Hurwitz schemes.
Let $`g2`$ be an even integer and $`H_g`$ be the Hurwitz scheme whose points correspond to degree 4 covers of $`^1`$ branched over $`n=2g+4`$ distinct points, such that the inverse image of $`n2`$ of these points consists of three points and the inverse image of each of the remaining points consists of two points, each of ramification degree $`2`$. From the Riemann-Hurwitz formula it follows that such a cover $`C^1`$ is of genus $`g`$.
###### Proposition 3.2.
For each $`g2`$, $`H_g`$ consists of two components.
###### Proof.
Let $`n=2g+4`$. Recall that we have a finite etale map $`\delta :H_g^n\mathrm{\Delta }`$, where $`^n`$ is thought of as $`(^1)^{(n)}`$ and $`\mathrm{\Delta }`$ is the discriminant locus. Let $`P=\{p_1,p_2,\mathrm{},p_n\}`$ be an element of $`^n\mathrm{\Delta }`$ and let $`x^1`$ be distinct from the $`p_i`$’s. We may choose loops $`\sigma _i`$ based at $`x`$ and going around $`p_i`$ in such a way so that $`\sigma _1,\sigma _2,\mathrm{}\sigma _n`$ generate $`G=\pi _1(^1\backslash \{p_1,p_2,\mathrm{},p_n\},x)`$ with the single relation as $`\sigma _1\sigma _2\mathrm{}\sigma _n=1`$. Now degree $`4`$ covers of $`^1`$ branched over $`P`$ correspond to transitive representations of $`G`$ in the symmetric group $`\mathrm{\Sigma }_4`$. For the covers to correspond to points of $`H_g`$, the images of $`n2`$ of the $`\sigma _i`$’s must be transpositions and the images of the other two must be products of $`2`$ disjoint transpositions. Two such representations give isomorphic covers if they differ by an inner automorphism of $`\mathrm{\Sigma }_4`$. Thus $`\delta ^1(P)`$ can be identified with classes of $`n`$-tuples $`(s_1,s_2,\mathrm{},s_n)`$ of elements of $`\mathrm{\Sigma }_4`$ , all but two of the $`s_i`$’s being transpositions, the remaining two being products of two disjoint transpositions, and $`s_1s_2\mathrm{}s_n=1`$. Two such $`n`$-tuples are identified if they differ by coordinatewise conjugation by an element of $`\mathrm{\Sigma }_4`$.
The action of the monodromy on $`\delta ^1(P)`$ contains elements $`\mathrm{\Gamma }_i`$, $`1in1`$, which act on $`n`$-tuples as above by:
$$\mathrm{\Gamma }_i(s_1,s_2,\mathrm{},s_n)=(s_1,\mathrm{},s_{i1},s_is_{i+1}s_i^1,s_i,\mathrm{},s_n).$$
To show that $`H_g`$ has two components we use the $`\mathrm{\Gamma }_i`$’s to prove that the monodromy action on $`\delta ^1(P)`$ has two orbits. The proof is a case by case analysis; we shall describe the main steps and leave some simple verifications to the reader. Let $`t_{i,j}`$, $`ij`$, denote the transposition which switches $`i`$ and $`j`$ and let $`v_1`$, $`v_2`$ and $`v_3`$ be the permutations $`(\mathrm{1\; 2})(\mathrm{3\; 4})`$, $`(\mathrm{1\; 3})(\mathrm{2\; 4})`$ and $`(\mathrm{1\; 4})(\mathrm{2\; 3})`$ respectively. Let $`V=\{v_1,v_2,v_3\}`$. Using the action of the $`\mathrm{\Gamma }_i`$’s one sees that each orbit contains an $`n`$-tuple $`S`$ such that $`s_{n1}`$ and $`s_n`$ are in $`V`$. Upto conjugation, we may assume that $`s_n=v_1`$ and $`s_{n1}=v_1`$ or $`v_2`$. Suppose $`s_{n1}=v_2`$, hence $`s_1s_2\mathrm{}s_{n2}=v_3`$. Let $`K`$ be the subgroup of $`\mathrm{\Sigma }_4`$ generated by $`s_1,s_2,\mathrm{}s_{n2}`$. Using the fact that $`v_3K`$ and $`K`$ is generated by transpositions we analyze the two possibilities for the action of $`K`$ on $`\{1,2,3,4\}`$ (i) transitive and (ii) intransitive. In both cases ((i) requires some computations with the $`\mathrm{\Gamma }_i`$’s) we conclude that we may assume that at least one of $`t_{1,4}`$ or $`t_{2,3}`$ occurs among the $`s_i`$’s (without changing $`s_{n1}`$ and $`s_n`$). Again using the action, we may assume that $`i=n2`$. Then replacing $`S`$, by $`\mathrm{\Gamma }_{n2}(\mathrm{\Gamma }_{n1}(S))`$ we may also assume that $`s_{n1}=v_1`$.
As above, let $`K`$ be the subgroup of $`\mathrm{\Sigma }_4`$ generated by $`s_1,s_2,\mathrm{}s_{n2}`$ and note that now $`s_1s_2\mathrm{}s_{n2}=1`$. We now consider four possibilities for the action of $`K`$ on $`\{1,2,3,4\}`$.
1. The action is transitive.
2. The action has a unique fixed point.
3. The action has two fixed points.
4. The action has no fixed points but fixes two disjoint subsets of two elements each.
(i) Since $`s_1s_2\mathrm{}s_{n2}=1`$, it follows by a result of Clebsch (see \[7, p.547\]) that we may assume
$$(s_1,s_2,\mathrm{}s_{n2})=(t_{1,2},t_{1,2},t_{1,3},t_{1,3},t_{1,4},\mathrm{}t_{1,4}).$$
Note that a priori we may also have to use an inner automorphism of $`\mathrm{\Sigma }_4`$ in order to achieve this, but it is easy to check using the action of the $`\mathrm{\Gamma }_i`$’s that we may choose an automorphism preserving $`v_1`$. Hence we may assume that
$$S=(t_{1,2},t_{1,2},t_{1,3},t_{1,3},t_{1,4},\mathrm{}t_{1,4},v_1,v_1).$$
(ii) Without loss of generality, we may assume that the fixed point is $`4`$. Again, by Clebsch’ result we may assume that
$$S=(t_{1,2},t_{1,2},t_{1,3},\mathrm{}t_{1,3},v_1,v_1).$$
Moreover, it is easy to check that if we replace $`v_1`$ by $`v_2`$ or $`v_3`$, we stay in the same orbit, hence all elements of type (ii) lie in one orbit. Further, note that using the $`\mathrm{\Gamma }_i`$’s , we can switch pairs of adjacent transpositions, i.e.
$$(t_{1,2},t_{1,2},t_{1,3},t_{1,3})(t_{1,3},t_{1,3},t_{1,2},t_{1,2}),$$
so we may replace $`S`$ by $`(t_{1,3},\mathrm{},t_{1,3},t_{1,2},t_{1,2},t_{1,3},t_{1,3},v_1,v_1)`$. One checks that
$$\begin{array}{c}\mathrm{\Gamma }_{n2}(\mathrm{\Gamma }_{n3}(\mathrm{\Gamma }_{n4}(\mathrm{\Gamma }_{n5}(\mathrm{\Gamma }_{n4}(\mathrm{\Gamma }_{n5}(\mathrm{\Gamma }_{n4}(\mathrm{\Gamma }_{n3}(\mathrm{\Gamma }_{n2}(\mathrm{\Gamma }_{n4}(S))))))))))\hfill \\ \hfill =(t_{1,3},\mathrm{},t_{1,3},t_{2,4},t_{1,2},t_{1,2},t_{2,4},v_1,v_1).\end{array}$$
Since we have assumed that $`n8`$, it follows that the $`K`$ corresponding to this element acts transitively on $`\{1,2,3,4\}`$. Hence elements of types (i) and (ii) lie in the same orbit.
(iii) In this case we must have $`S=(t,t,\mathrm{},t,v_1,v_1)`$ where $`t`$ is one of $`t_{1,3}`$, $`t_{2,3}`$, $`t_{1,4}`$, $`t_{2,4}`$, since the representation is assumed to be transitive. It is clear that all choices are conjugate by an element of $`\mathrm{\Sigma }_4`$ which fixes $`v_1`$, hence all elements of type (iii) are in the same orbit.
Assume, without loss of generality that $`t=t_{1,3}`$. Then
$$\mathrm{\Gamma }_{n2}(\mathrm{\Gamma }_{n3}(\mathrm{\Gamma }_{n3}(\mathrm{\Gamma }_{n2}(S))))=(t_{1,3},t_{1,3},\mathrm{},t_{1,3},t_{2,4},t_{2,4},v_1,v_1).$$
Since $`t_{1,3}`$ and $`t_{2,4}`$ commute, $`\mathrm{\Gamma }_i`$, for $`1in2`$, just switches $`s_i`$ and $`s_{i+1}`$. Thus we may repeat the above procedure and deduce that $`S`$ is in the same orbit as elements of the form $`(s_1,s_2,\mathrm{},s_{n2},v_1,v_2)`$ where each $`s_i`$ for $`1in2`$, is either $`t_{1,3}`$ or $`t_{2,4}`$, there being an even number of both kinds.
(iv) Elements of type (iv) are precisely those considered in the previous paragraph, so all such elements lie in the same orbit as elements of type (iii).
We thus see that $`H_g`$ has at most two components. Now observe that covers of type (iii) have an automorphism of order $`2`$ commuting with the covering map, corresponding to the representation of $`G`$ in $`\mathrm{\Sigma }_2`$ induced from the original representation in $`\mathrm{\Sigma }_4`$ by identifying $`1`$ with $`2`$ and $`3`$ with $`4`$. Covers of type (i) have no automorphisms commuting with the covering map, hence $`H_g`$ has precisely two components. ∎
### 3.3. Indecomposability of the generic $`4`$-configuration
Let $`H_g^{}`$ be the component of $`H_g`$ corresponding to covers without automorphisms. Then there exists a universal family of curves $`\psi :𝒞_gH_g^{}`$ and a corresponding universal degree $`4`$ map $`\pi :𝒞_gH_g^{}\times ^1`$.
We continue using the same notation as in the proof of Proposition 3.2, so $`P=\{p_1`$, $`p_2`$, …, $`p_n\}`$, $`p_i^1`$, and let $`p`$ be any other point of $`^1`$. Let $`X`$ be the curve in $`^n=()^{(n)}`$ with points $`\{(1t)p_1+tp,(1t)p_2+tp,p_3,p_4,\mathrm{}p_n\}`$, $`t`$. By the previous proposition, there exists an element of $`\delta ^1(P)`$ such that the monodromy representation of the corresponding cover is given by the $`n`$-tuple $`S=(t_{2,3},t_{2,3},t_{1,2},t_{1,2},\mathrm{},t_{1,2},v_1,v_1)`$ of elements of $`\mathrm{\Sigma }_4`$. Let $`Y^{}`$ be the component of $`\delta ^1(X)`$ containing this point and let $`Y`$ be the normalization of $`X`$ in the function field of $`Y^{}`$. Then $`𝒞=\psi ^1(Y^{})`$ maps to $`X\times ^1`$ by the map $`\pi `$ and we let $`\overline{𝒞}`$ be the normalization of $`X\times ^1`$ in the function field of $`𝒞`$. Clearly $`\psi |_𝒞`$ induces a flat and projective morphism $`\overline{\psi }:\overline{𝒞}Y`$. Let $`y_0`$ be a point of $`Y`$ lying above the point $`\{p,p,p_3,p_4,\mathrm{}p_n\}`$ of $`X`$.
###### Lemma 3.1.
The fibre of the map $`\overline{\psi }:\overline{𝒞}Y`$ over $`y_0`$ has two smooth components, $`C_1`$ and $`C_2`$, meeting transversally in a single point lying over $`p`$. The map from $`C_1`$ to $`^1`$ has degree $`2`$ and is branched over $`p_3`$, $`p_4`$, …, $`p_n`$, while the map from $`C_2`$ to $`^1`$ is branched over $`p_{n1}`$ and $`p_n`$. In particular, $`C_1`$ is hyperelliptic of genus $`g=(n4)/2`$ while $`C_2`$ is of genus $`0`$.
###### Proof.
Let $`B`$ be small topological disc in $`Y`$ containing $`y_0`$ and such that the only point where the map from $`B`$ to $`X`$ is ramified is $`y_0`$. By construction of $`X`$, if $`B`$ is small enough then the ramification locus of the map $`\pi `$ intersected with $`(B\{y_0\})\times ^1`$ consists of $`n`$ disjoint punctured discs, each mapping isomorphically to $`B\{y_0\}`$. The closures of all these discs remain disjoint in $`B\times ^1`$ except for two of them which meet at $`(y_0,a)`$. By construction, the product of the local monodromies around these two discs is trivial, hence the map $`\pi |_{\overline{\psi }^1(B)}:\overline{\psi }^1(B)B\times ^1`$ is unramified over the complement of the closures of all the above discs.
To complete the proof we examine the induced monodromy representation of $`\pi _1(^1\backslash \{a,p_3,p_4,\mathrm{},p_n\},x)`$ in $`\mathrm{\Sigma }_4`$. The local monodromy around $`a`$ is equal to the product of the local monodromies around $`p_1`$ and $`p_2`$ of the original representation used in the construction of $`Y`$ and the local monodromies around the other points are the same as those in the original representation. Then we see that the representation is no longer transitive but breaks up into two representations, each of degree $`2`$. The conclusions of the lemma follow immediately on inspection of these two representations. ∎
For any point $`h`$ of $`H_g^{}`$, the covering map $`C=(𝒞_g)_h^1`$ gives us points $`a_1`$, $`a_2`$, $`b_1`$, $`b_2`$ on $`C`$ such that $`2[a_1+a_2]=2[b_1+b_2]`$. Choosing a function $`f`$ on $`C`$ such that $`div(f)=2(a_1+a_2)2(b_1+b_2)`$ allows us to define a $`4`$-configuration $`Z`$ in $`Pic^3(C)`$. Note that $`Z`$ is well defined upto sign in $`CH_{\mathrm{ind}}^g(Pic^3(C),1)`$.
###### Theorem 3.1.
The $`4`$-configuration corresponding to a generic point of $`H_g^{}`$ is indecomposable for all $`g3`$. Moreover $`CH_{\mathrm{ind}}^k(J(C),1)_s`$ is uncountable for $`4skg`$.
###### Proof.
Let $`\overline{𝒞}Y`$ be the family constructed in the discussion preceding Lemma 3.1. We blow down the genus $`0`$ curve $`C_2`$ in the fibre over $`y_0`$ and call the resulting family of curves $`𝒞^{}`$. By replacing $`Y`$ by a Zariski open subset we may assume that all fibres are smooth and then by replacing $`Y`$ by a finite cover we may also assume that branch locus of the map $`\overline{𝒞}Y\times ^1`$ is a union of sections. Finally, by replacing $`Y`$ by a further open subset (containing $`y_0`$) we may assume that there exist a rational function $`F`$ on $`𝒞^{}`$ such that $`div(F)=2(A_1+A_2)2(B_1+B_2)=2𝒟`$, where $`𝒟`$ restricts to the divisor used to define the $`4`$-configuration on each fibre.
$`F`$ allows us to construct an element $`𝒵`$ of $`CH^g(Pic^3(𝒞^{}/Y),1)`$ which restricts to the $`4`$-configuration on each fibre. By Lemma 3.1, we see that upto labelling the restrictions of $`A_1`$, $`B_1`$, $`A_2`$, $`B_2`$ to $`C_1=𝒞_{}^{}{}_{y_0}{}^{}`$ must be $`w_1`$, $`w_2`$, $`t`$, $`t`$, where $`w_1`$ and $`w_2`$ are distinct Weierstrass points on $`C_1`$ and $`t`$ is a point lying over the point $`a^1`$. It is then clear that $`f_{y_0}=F|_{C_1}`$ has divisor $`2w_12w_1`$ and is therefore a Weierstrass function. One easily checks that the element of $`CH^g(J(C_1),1)`$ obtained by translating $`𝒵_{y_0}`$ by $`3w_1`$ is equal to $`KK_t`$, where $`K`$ is the basic hyperelliptic cycle and the subscript denotes translation. Since $`p_1`$, $`p_2`$, …, $`p_n`$ and $`p`$ are arbitrary points in $`^1`$, it follows that we may assume that $`C_1`$ is a generic hyperelliptic curve of genus $`g`$ and $`t`$ a generic point on $`C_1`$. By Theorem 1.2 it follows that $`𝒵_{y_0}`$ is indecomposable and then by Theorem 2.1 it follows that $`𝒵`$ restricted to a generic fibre is also indecomposable.
The second statement follows by considering the Pontryagin product of the $`4`$-configuration with the zero cycles $`(t_1t_0)(t_2t_0)\mathrm{}(t_rt_0)`$, $`0rg3`$, where $`(t_0,t_1,\mathrm{},t_r)`$ is a generic point of $`C^{r+1}`$. Taking $`r=g3>0`$, we see from the proof of Theorem 1.2 that the components of elements of $`ZI^{(g3)}`$ give uncountably many elements of $`CH_{\mathrm{ind}}^g(J(C),1)_g`$ (note that specialization preserves the decompositions). We then deduce that $`CH_{\mathrm{ind}}^g(J(C),1)_s`$ is uncountable for $`4sg`$, as well as the statement for $`k<g`$, in the same way as in Theorem 1.2. ∎
Consider the natural map $`\tau _g:H_g^{}_g`$. Since the map $`C^{(2)}\times C^{(2)}J(C)`$ given by $`(\{a_1,a_2\},\{b_1,b_2\})[a_1+a_2(b_1+b_2)]`$ is surjective for any curve of genus $`g4`$, by counting parameters we see that $`\tau _g`$ is dominant for $`g=3,4`$. For $`g>4`$ a degeneration argument shows that the image has dimension $`2g+1`$. Note that the (non-empty) fibres of $`\tau _g`$ are at least $`3`$ $`=dim(Aut(^1))`$ dimensional and the generic fibre of $`\tau _3`$ is $`4`$ dimensional.
###### Corollary 3.1.
For $`C`$ a generic curve of genus $`g=3,4`$, $`CH_{\mathrm{ind}}^g(J(C),1)_g`$ is uncountable.
###### Proof.
The statement for $`g=4`$ follows from the theorem. For $`g=3`$, the generic curve has a $`1`$-parameter family of $`4`$-configurations. Consider a family of genus 3 curves with special fibre a hyperelliptic curve as in the theorem. By varying the 4-configuration in the Jacobian of the generic fibre in the $`1`$-parameter family, we obtain as specializations cycles of the form $`KK_t`$ where now $`t`$ varies in a $`1`$-parameter family. By Theorem 1.2, it follows that the set of specializations is an uncountable subset of $`CH_{\mathrm{ind}}^3(J(C_1),1)_3`$. Hence by Theorem 2.1 the $`4`$-configurations also form an uncountable subset of $`CH_{\mathrm{ind}}^3(J(C),1)_3`$. ∎
## 4. Higher Chow cycles on self products of a curve
### 4.1. Construction of the cycles
The indecomposability of the basic hyperelliptic cycle and the 4-configuration suggest that to construct indecomposable elements in $`CH^g(J(C),1)`$ for more general Jacobians, we should consider more general divisors $`D`$ on $`C`$, with $`[D]`$ of order 2 in $`Pic(C)`$. Let $`D`$ be such a divisor, so $`D=_{i=1}^na_i_{i=1}^nb_i`$ for some $`n>0`$ and distinct points $`a_i,b_iC`$, and there exists a rational function $`f`$ on $`C`$ with $`div(f)=2D`$. Note that if $`0<g2n`$ then there always exist such divisors on a general curve, and if $`g2n`$ then there exists a $`2g+2n3`$ dimensional family of curves of genus $`g`$ which have such divisors.
Let $`E`$ be the set of embeddings of $`C`$ in $`C^{n+1}`$ such that the composition of any $`eE`$ with the projection to any factor is either the identity map or is constant with image equal to one of the $`a_i`$’s or $`b_i`$’s. We have the following group actions on the set $`E`$.
1. $`\mathrm{\Sigma }_{n+1}`$ acts by permuting the factors.
2. $`\mathrm{\Sigma }_n`$ acts by permuting the $`a_i`$’s.
3. $`\mathrm{\Sigma }_n`$ acts by permuting the $`b_i`$’s.
4. $`/2`$ acts by switching $`a_i`$ and $`b_i`$ for all $`i`$.
Let $`f_e`$ be the function $`f`$ considered as a function on $`e(C)`$.
###### Proposition 4.1.
The set of higher Chow $`1`$-cycles of the form $`_{eE}m_e(e(C)f_e)`$, $`eE`$, $`m_e`$ invariant under all the above group actions on $`E`$, is a positive dimensional vector space for all $`n>0`$.
###### Proof.
By a partition of a positive integer $`n`$ we shall mean a tuple of non-increasing positive integers $`\alpha =(n_1,n_2,\mathrm{},n_r)`$ such that $`|\alpha |:=_{i=1}^rn_i=n`$ and by a partition of zero we mean the empty tuple $`()`$. If $`\alpha `$ is as above and $`m`$ is a positive integer, we let $`\alpha +m`$ be the partition of $`n+m`$ obtained by reordering $`(n_1,n_2,\mathrm{},n_r,m)`$.
The orbits of the action generated by all the above groups on $`E`$ are in $`1`$-$`1`$ correspondence with the set
$$=\left\{\{\alpha ,\beta \}\right|\alpha \text{ a partition of }i,\beta \text{ a partition of }j,i+jn\}.$$
Here the elements of $``$ (and $`𝒫`$ below) are viewed as unordered pairs.
The boundary of an element $`_{eE}m_e(e(C)f_e)`$ is supported on points of $`C^{n+1}`$, each of whose coordinates is one of the $`a_i`$’s or $`b_i`$’s. The orbits of the set of all such points under the group action are in $`1`$-$`1`$ correspondence with the set
$$𝒫=\left\{\{\alpha ,\beta \}\right|\alpha \text{ a partition of }i,\beta \text{ a partition of }j,i+j=n+1\}\backslash \{\{(1,\mathrm{},1),()\}\}.$$
$`\{(1,\mathrm{},1),()\}`$ is not included because the number of $`a_i`$’s and $`b_i`$’s is $`n`$.
Let $`\times 𝒫`$ be the set
$$\underset{\{\alpha ,\beta \}}{}\{(\{\alpha ,\beta \},\{\alpha ,\beta +(n+1|\alpha ||\beta |)\}),(\{\alpha ,\beta \},\{\alpha +(n+1|\alpha ||\beta |),\beta \})\}.$$
Consider the projection $`p_2:𝒫`$ and observe that $`|p_2^1(\{\alpha ,\beta \})|2`$ for all $`\{\alpha ,\beta \}𝒫`$ except for those of the form (i) $`\alpha =(d,d,\mathrm{},d)`$, $`\beta =()`$, with $`d>1`$ a divisor of $`n+1`$ and (ii) $`\alpha =\beta =(d,d,\mathrm{},d)`$ with $`n`$ odd and $`d`$ a divisor of $`(n+1)/2`$. For both these cases $`|p_2^1(\{\alpha ,\beta \})|=1`$. Using the fact that for any integers $`d,m`$, if $`m/2<d<m`$ then $`d`$ cannot divide $`m`$, one checks that if $`n>3`$, then $`\left|\{\{\alpha ,\beta \}𝒫||p_2^1(\{\alpha ,\beta \})|=1\}\right|n`$.
On the other hand, $`|p_1^1(\{\alpha ,\beta \})|=2`$ for all $`\{\alpha ,\beta \}`$ except for those with $`\alpha =\beta `$ or $`\alpha =(1,1,\mathrm{},1)`$, $`|\alpha |=n`$, and $`\beta =()`$. For both these cases $`|p_1^1(\{\alpha ,\beta \})|=1`$. If $`\mathrm{\Pi }(m)`$ denotes the number of partitions of $`m`$, one sees that the number of such elements is $`1+_{i=0}^{[n/2]}\mathrm{\Pi }(i)`$ which is greater than $`n`$ for all $`n>5`$.
Elements of $`𝒫`$ give us sufficient relations among the $`m_e`$’s for the boundary of $`_{eE}m_e(e(C)f_e)`$ to be zero, but if $`n`$ is odd then elements of the form $`\{\alpha ,\alpha \}`$, $`|\alpha |=(n+1)/2`$, give trivial relations because of symmetry. Using this observation along with explicit computations for $`1n5`$, one sees that the number of relations is always strictly less than the number of variables i.e. $`||`$, hence the space of invariant cycles is positive dimensional. ∎
Let $`V(C,f)`$ be the space of all invariant higher Chow 1-cycles as in the proposition. By abuse of notation we shall also denote this space by $`V(C,D)`$. For $`n=1`$ and $`2`$, $`V(C,D)`$ is of rank $`1`$, but it is of rank $`>1`$ for all $`n>2`$.
### 4.2. Indecomposability of the product cycles via specialization
Proposition 4.1 provides us with elements of $`CH^{n+1}(C^{n+1},1)_{}`$. After pushforward by the natural map $`C^{n+1}Pic^{n+1}(C)`$ and translation by a divisor of degree $`(n+1)`$, we obtain elements of $`CH^g(J(C),1)_{}`$. We now outline a method which should allow one to prove indecomposability of the image of a general element of $`V(C,D)`$ in $`CH^g(J(C),1)_{}`$, for $`C`$ generic among curves of genus $`gn+1>2`$, having such divisors $`D`$.
Let $`_g`$ be the moduli space of genus $`g`$ curves (pointed curves if $`g=1`$) with level $`2m`$ structure for some $`m3`$ and let $`𝒞_g`$ be the universal family of curves over $`_g`$. Consider the map $`\sigma :Sym^n(𝒞_g/_g)\times Sym^n(𝒞_g/_g)J(𝒞_g/_g)`$ given on the fibres over $`_g`$ by $`(\{a_1,a_2,\mathrm{},a_n\},\{b_1,b_2,\mathrm{},b_n\})[_na_i_nb_i]`$ and let $`X_{g,n}=\sigma ^1(S_2)`$, where $`S_2`$ is the union of sections of $`J(𝒞_g/_g)`$ corresponding to points of order $`2`$ on each fibre. Let $`\tau `$ be the composite of $`\sigma `$ with the natural map from $`J(𝒞_g/_g)`$ to $`_g`$.
###### Lemma 4.1.
Let $`n>1`$ and $`g>0`$. Then there exists an irreducible component $`Y_{g,n}`$ of $`X_{g,n}`$ with the following properties:
1. $`dim(Y_{g,n})=2g+2n3`$.
2. If $`g2n`$, then the map $`\tau |_{Y_{g,n}}:Y_{g,n}_g`$ is dominant and if $`g2n`$, then $`dim(\tau (Y_{g,n}))=2g+2n3`$.
3. In the fibre of $`\tau |_{Y_{g,n}}`$ over a point on $`_g`$ corresponding to a generic hyperelliptic curve $`C^{}`$, there exist points $`(\{w_1,t_1,t_2,\mathrm{},t_{n1}\},\{w_2,t_1,t_2,\mathrm{},t_{n1}\})`$, where $`w_1`$ and $`w_2`$ are distinct Weierstrass points on $`C^{}`$, and $`(t_1,t_2,\mathrm{},t_{n1})`$ is a generic point of $`C_{}^{}{}_{}{}^{n1}`$.
###### Outline of proof.
The lemma can be proved by considering a suitable Hurwitz scheme as in section 3. We do not know the number of components if $`n>2`$, but a suitable choice of the monodromy representation allows us to single out a component which gives rise to the desired specializations (c.f. the discussion before Lemma 3.1). For example, if $`n=3`$ we would consider the representation in $`\mathrm{\Sigma }_6`$ corresponding to the tuple
$$((23),(23),(45),(45),(12),(12),\mathrm{},(12),(12)(34)(56),(12)(34)(56)).$$
We let $`Y_{g,n}`$ be the closure in $`X_{g,n}`$ of points of the form $`C`$, $`\{a_1,a_2,\mathrm{},a_n\}`$,
$`\{b_1,b_2,\mathrm{},b_n\}`$, where $`C`$ is a cover of $`^1`$ corresponding to a point on the chosen component of the Hurwitz scheme, and $`\{a_1,a_2,\mathrm{},a_n\}`$, $`\{b_1,b_2,\mathrm{},b_n\}`$ are the inverse images of the two points of $`^1`$ over which the cover is not simply ramified. Since the dimension of each component of the Hurwitz scheme is $`2g+2n`$ and $`dim(Aut(𝒫^1))=3`$, it follows that $`dim(Y_{g,n})=2g+2n3`$.
To prove the statement about the $`\tau (Y_{g,n})`$, we consider a degeneration of the cover such that three of the simply ramified points come together (generically) at a point. For instance, in the above example we would let $`p_5`$, $`p_6`$ and $`p_7`$ come together. If $`g2`$, the special fibre of the stable model of the degenerating family of genus $`g`$ curves then consists of two smooth components, one of them a curve of genus $`g1`$ which is a cover of $`^1`$ of the same type, and the other a generic elliptic curve. If $`g=1`$, we obtain a nonconstant family of elliptic curves and so the statement is true in this case. The statement for $`g=2`$ follows, since we may then assume that both components are generic elliptic curves.
If $`2<g2n`$, we may assume using induction that the point of intersection of the two components is a generic point on the curve of genus $`g1`$ (also generic). Thus $`dim(\tau (Y_{g,n}))dim(\tau (Y_{g1,n}))+3`$, and so $`\tau |_{Y_{g,n}}`$ must be dominant. Finally, if $`g>2n`$ we see that $`dim(\tau (Y_{g,n}))dim(\tau (Y_{g1,n}))+2`$. Since $`dim(Y_{g,n})=2g+2n3`$ and $`dim(\tau (Y_{2n,n}))=4n3`$, it follows that we must have equality for all $`g2n`$.
The basic idea for the proof of indecomposability is now the same as that used for the 4-configuration i.e. we specialize to a hyperelliptic curve and then use the results of sections 1 and 2. In somewhat more detail, the argument is as follows. By Lemma 4.1 we may construct a family of smooth curves $`𝒞S`$ of genus $`g`$ and a divisor $`𝒟=_nA_i_nB_i`$ on $`𝒞`$, with $`S`$ a smooth curve with a distinguished point $`s`$, such that the fibre over the generic point, $`(C,D)`$ corresponds to a generic point of $`Y_{g,n}`$ and the fibre over $`s`$, $`(C^{},D^{})`$, corresponds to the special points of $`Y_{g,n}`$ in Lemma 4.1 (iii).
Let $`W`$ be the subspace of $`CH^g(J(C^{}),1)`$ obtained by mapping $`V(C,D)`$ to $`CH^g(Pic^{n+1}(C),1)_{}`$, specialization to $`CH^g(Pic^{n+1}(C^{}),1)_{}`$, and then translation by $`(n+1)[w_1]`$. Note that specialization is always defined here, since we are free to modify $`f`$ by a nonzero constant.
By the condition on $`D^{}`$ it follows that $`f^{}`$, the specialization of $`f`$ must be a Weierstrass function on $`C^{}`$ corresponding to the divisor $`2w_12w_2`$. The components of the support of elements of $`W`$ are all images of the various diagonals in $`C_{}^{}{}_{}{}^{n+1}`$, and the function on each component is $`f^{}`$. Lemma 1.6 then allows us to assume that modulo decomposable elements each element of $`W`$ is a sum of translates of the basic hyperelliptic cycle $`K`$. It is then clear that the results of sections 1 and 2 imply that if $`gn+1`$, the general element of $`V(C,D)`$ (with $`(C,D)`$ corresponding to a generic point of $`Y_{g,n}`$) is indecomposable, provided that the following hypothesis is satisfied.
###### Hypothesis 4.1.
The image of $`W`$ in $`CH_{\mathrm{ind}}^g(J(C^{}),1)_{}`$ is the vector space spanned by the cycle $`K([t_1w_1]e)([t_2w_1]e)\mathrm{}([t_{n1}w_1]e)`$.
###### Remark.
The difficulty in verifying the hypothesis is purely combinatorial; everything can be computed explicitly for any given $`n`$. We have verified it using a computer for $`2n6`$. (In all these cases, the kernel of the map from $`W`$ to $`CH_{\mathrm{ind}}^g(J(C^{}),1)_{}`$ consists of those cycles which do not contain the (small) diagonal of $`C^{n+1}`$ in their support.) Using the dimension formula from Lemma 4.1, we see that $`CH_{\mathrm{ind}}^g(J(C),1)_{}`$ is nonzero for a generic curve with $`3genus(C)12`$. As in the case of the $`4`$-configuration, Pontryagin product with zero cycles can be used to prove that for these cases $`CH_{\mathrm{ind}}^g(J(C),1)_{}`$ is in fact uncountable.
## 5. Higher analogues in lower genus
In the first part of the section we construct elements $`FCH^3(J(C),4g)`$, where the genus $`g2`$. The cycles $`F`$ are natural generalizations of the 4-configuration and we expect that they should be strongly indecomposable in the sense that we explain below. In the second part we study a certain cycle $`BCH^3(J(C),2)`$ where $`C`$ is a bielliptic curve of genus $`2`$. We show that $`B`$ is (weakly) indecomposable using Lewis’ criterion, which we show to hold by means of the same kind of proof as the one which was given in (1.3).
### 5.1. Higher analogues of the 4-configuration in lower genus
Bloch’s groups $`CH^p(X,n)`$ can be described by means of chains built from those integral subvarieties of codimension $`p`$ in $`X\times (^1\{1\})^n`$ which meet all the cubical faces over $`0`$ or over $`\mathrm{}`$ properly. Consider a semistable degeneration $`B^{}`$ of an abelian variety $`B`$ with trivial extension class, then $`B^{}=A\times ^{}`$ with compactification $`\overline{B}^{}=A\times ^1`$ and thus one may expect that $`CH^m(A,n)`$ should retain memory of the properties of $`CH^m(B,n1)`$. This was the guess that prompted our construction of the 4-configuration, which in a way can be seen as being the memory of cycles studied in . In the same vein we describe now a unified procedure for building a series of cycles $`F`$ in $`CH^3(J(C),4g)`$, $`g=genus(C)`$ $`3`$. The first is the 4-configuration and each element is the memory of the preceding one.
Choose a class $`ϵ`$ of order $`2`$, consider a divisor $`D:=(a_1+a_2)(b_1+b_2)`$ on $`C`$ with $`class(D)=ϵ`$, and let $`f`$ be a rational function with $`div(f)=2D`$. We embed $`CJ(C)`$ using the maps $`\alpha _1(x)=[xa_1]`$, $`\alpha _2(x)=[x+a_2]`$, $`\alpha _3(x)=[xa_1]+ϵ`$, $`\alpha _4(x)=[x+a_2]+ϵ`$, and define $`C_i:=\alpha _i(C)`$. The useful property here is the fact that on $`C_1`$ the point which comes from $`a_i`$ coincides with the point on $`C_2`$ from $`a_{i\pm 1}`$. The same happens for $`b_i`$ with respect to the curves $`C_1`$ and $`C_4`$, and things stay the same for every other curve instead of $`C_1`$.
We start from genus $`3`$, here we take $`\beta _i=CJ(C)\times ^1`$ to be the map $`x(\alpha _i(x),f(x)^{s(i)})`$, where $`s(i)`$ means $`1`$, according to the parity of i.
The curves $`K_i:=\beta _i(C)`$ in $`J(C)\times ^1`$ meet properly the cubical faces, $`_{i=1}^4K_i`$ is indeed a cycle and this is our element
$$F(f):=\underset{i=1}{\overset{4}{}}K_iCH^3(J(C),1).$$
Clearly $`F(f)`$ is equivalent to the the 4-configuration as it was previously defined.
To go higher in Chow groups we need at each step a new and convenient rational function. It can be found by imposing restrictions on $`D`$ by means of conditions on $`f`$ and this is the reason why we need to go down in genus.
In genus $`2`$ we require $`ϵ=[w_1w_2]`$, the distinguished ramification points of the Weierstrass double cover $`h:C^1`$. The condition is
(+)
$$f(w_1)=f(w_2)$$
which is satisfied by a 1- dimensional family of $`a`$’s and $`b`$’s as above. Weil reciprocity and (+) yield
(\*)
$$(h(a_1)h(a_2))^2=(h(b_1)h(b_2))^2.$$
We may change $`h`$ by a multiplicative constant so as to have
(\**)
$$1=(h(a_1)h(a_2))^2=(h(b_1)h(b_2))^2.$$
Writing now $`\beta _i(x):=(\alpha _i(x),f(x)^{s(i)},h(x)^{2s(i)})`$ then $`K_i:=\beta _i(C)`$ is a curve in $`J(C)\times ^1\times ^1`$ which meets properly the cubical faces. It is easy to check that
$$F(f,h):=\underset{i=1}{\overset{4}{}}K_iCH^3(J(C),2),$$
for instance one can see that the boundary is trivial over the point of $`C_1`$ which comes from $`w_1`$ by realizing that it coincides with the point on $`C_3`$ from $`w_2`$, and by using then the condition $`f(w_1)=f(w_2)`$.
In genus $`1`$ having fixed the divisor $`D`$ of class $`ϵ`$ as before we choose further two rational functions $`h_1`$ and $`h_2`$ of degree $`2`$ on $`E`$ both ramified over $`0`$ and $`\mathrm{}`$. Let $`div(h_i)=2(q_i^{}q_i^{\prime \prime })`$. We require our choice to satisfy:
(1) $`[q_1^{}q_1^{\prime \prime }]=[q_2^{}q_2^{\prime \prime }]=ϵ`$
(2) $`h_1(q_2^{})=h_1(q_2^{\prime \prime })andh_2(q_1^{})=h_2(q_1^{\prime \prime }).`$
This can be done, see . Here the conditions on $`f`$ read
(+)
$$f(q_i^{})=f(q_i^{\prime \prime })i=1,\mathrm{\hspace{0.17em}2}.$$
As it was before this implies
(\*)
$$((h_i(a_1)h_i(a_2))^2=(h_i(b_1)h_i(b_2))^2.$$
We normalize the choice of the rational functions $`h_i`$ by the request:
(\**)
$$1=((h_i(a_1)h_i(a_2))^2=(h_i(b_1)h_i(b_2))^2.$$
Using $`\beta _i(x):=(\alpha _i(x),f(x)^{s(i)},h_1(x)^{2s(i)},h_2(x)^{2s(i)})`$, we obtain
$$F(f,h_1,h_2):=\underset{i=1}{\overset{4}{}}K_iCH^3(E,3).$$
The 4-configuration $`F(x,f,h_1,h_2)CH^3(,4)`$ is constructed in the same way as $`F(f,h_1,h_2)`$ was. One thinks of $`^1`$ as having the structure of a degenerate Jacobian, given by the choice of a standard parameter $`x`$. The opposite map from $`J(C)`$ to $`J(C)`$ becomes $`xx^1`$ and the sum of points corresponds to product of the coordinates. Translation by $`ϵ`$ is here multiplication by $`1`$. The condition $`class(D)=ϵ`$ reads $`a_1a_2=b_1b_2`$ and $`f:=(xa_1)^2(xa_2)^2(xb_1)^2(xb_2)^2`$. We may take explicitly: $`h_1=c_1(x1)^2(x+1)^2`$ , $`h_2=c_2(xi)^2(x+i)^2`$. The requirements on $`f`$ are
(+)
$$f(1)=f(1),f(i)=f(i)$$
Our choice yields also $`:f(0)=f(\mathrm{})`$. One has
(\*)
$$(h_i(a_1)h_i(a_2))^2=(h_i(b_1)h_i(b_2))^2,i=1,2.$$
Choose and fix the constants $`c_1`$ and $`c_2`$ so that it is
(\**)
$$1=(h_i(a_1)h_i(a_2))^2=(h_i(b_1)h_i(b_2))^2,i=1,2.$$
With this dictionary the maps $`\alpha _i:^1^1`$ are defined as before (for instance $`\alpha _1(t)=ta_1^1`$ ), here the range of $`\alpha _i`$ should be understood as the replacement of $`J(C)`$. In this way the maps $`\beta :^1(^1)^4`$ are here $`\beta _i(x)=(\alpha _i(x),f(x)^{s(i)},h_1(x)^{2s(i)},h_2(x)^{2s(i)})`$. Setting again $`K_i:=\beta _i(^1)`$ our cycle is then
$$F(x,f,h_1,h_2):=\underset{i=1}{\overset{4}{}}K_iCH^3(^{},4).$$
###### Remark.
One may define an element in $`CH^a(X,b)`$ to be strongly indecomposable if it is not in the image of $`CH^{a1}(X,b1)^\times `$. We think that the higher 4-configurations are good candidates to strong indecomposability.
### 5.2. The B-configuration
Following the terminology of we define the group of (weakly) decomposable cycles in $`CH^p(X,2)H^{p2}(X,𝒦_p)`$ to be :
$$CH_{\mathrm{dec}}^p(X,2):=Im\left\{K_2()CH^{p2}(X)CH^p(X,2)\right\},$$
and thus the indecomposable group is
$$CH_{\mathrm{ind}}^k(X,2):=CH^k(X,2)/CH_{\mathrm{dec}}^k(X,2).$$
It is known that translations on an elliptic curve $`E`$ act trivially on $`CH_{\mathrm{ind}}^2(E,2)`$, see \[8, 3.10\]. We show that on the contrary translations on a genus $`2`$ bielliptic Jacobian $`J(G)`$ operate non trivially on $`CH_{\mathrm{ind}}^3(J(G),2)`$. Our procedure is similar to the one which we have applied above in (1.3). We deal here with the B-configuration which is shown to be indecomposable by checking Lewis’ condition on a cycle $`𝔅`$ of $`CH^3(J(G)\times G,2)`$.
S. Bloch in his seminal memoir constructed certain elements $`S_b\mathrm{\Gamma }(E,𝒦_2))`$ associated with a point $`b`$ of finite order on an elliptic curve $`E`$. He proved that the real regulator image of $`S_b`$ is not trivial for some curves with complex multiplication, and thus it is not trivial in general.
Consider a bielliptic curve $`G`$ of genus $`2`$ with associated map $`\delta _G:GE_1`$, and let $`a:GJ(G)`$ be the Abel Jacobi map. In this way $`Z(b):=a_{}\delta _G^{}(S_b)`$ is a cycle in $`CH^3(J(G),2)`$. Translation by an element $`tPic^0(G)`$ maps $`Z(b)`$ to the cycle $`Z_t(b)`$, our aim is to prove
###### Theorem 5.1.
The B-configuration $`B(t):=Z_t(b)Z(b)`$ is indecomposable for generic $`t`$.
Note that $`B(t)`$ has trivial regulator image.
###### Proof.
Consider the cycle $`G\times Z(b)CH^2(G\times G,2)`$. The straight embedding $`\sigma :=id\times a:G\times GG\times J(G)`$ maps it to $`𝔖:=\sigma _{}(G\times Z(b))`$ in $`CH^3(G\times J(G),2)`$. The twisted embedding $`\tau (t,x):=(t,a(x)+(tw_1)))`$ gives instead $`𝔗:=\tau _{}(G\times Z(b))`$, with section $`𝔗_{}(t)=Z_t(b)`$, and therefore $`B(t)`$ is the section at $`tG`$ of $`𝔅:=𝔗𝔖`$.
We use the same type of notations as we did in part 1.2, in particular the holomorphic form $`\omega _i^J`$ comes from $`E_i`$. We need to consider also the forms $`\nu :=\overline{\omega }_1^J\omega _2^J`$ on $`J(G)`$ and $`\overline{\zeta _2}`$ on $`G`$. The procedure of 1.2 gives here again: (i) $`<R(𝔅),\overline{\zeta _2}\nu >0`$. The Neron-Severi space of divisors with rational coefficients on $`J(G)`$ is isomorphic to the same space on the product of the two associated elliptic curves. On the gneric bielliptic Jacobian $`\nu `$ is orthogonal to the Neron-Severi group, because it is orthogonal to the elliptic curves. A look at the proof of the main theorem of shows that this property of $`\nu `$ and (i) imply that the generic section $`𝔅_{}(t)`$ is indeed weakly indecomposable. ∎
###### Acknowledgements.
The second author would like to thank R. Sreekantan for some useful correspondence.
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# Analysis of the Decays 𝐵→𝜋𝜋 and 𝜋𝐾 with QCD Factorization in the Heavy Quark Limit *footnote **footnote *Supported in part by National Natural Science Foundation of China and State Commission of Science and Technology of China
## Acknowledgement
We thank Dr. Z.T. Wei for helpful discussions.
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# Linking Classical and Quantum Key Agreement: Is There “Bound Information”?
## 1 Introduction
In modern cryptography there are mainly two security paradigms, namely computational and information-theoretic security. The latter is sometimes also called unconditional security. Computational security is based on the assumed hardness of certain computational problems (e.g., the integer-factoring or discrete-logarithm problems). However, since a computationally sufficiently powerful adversary can solve any computational problem, hence break any such system, and because no useful general lower bounds are known in complexity theory, computational security is always conditional and, in addition to this, in danger by progress in the theory of efficient algorithms as well as in hardware engineering (e.g., quantum computing). Information-theoretic security on the other hand is based on probability theory and on the fact that an adversary’s information is limited. Such a limitation can for instance come from noise in communication channels or from the laws of quantum mechanics.
Many different settings in the classical noisy-channel model have been described and analyzed, such as Wyner’s wire-tap channel , Csiszár and Körner’s broadcast channel , or Maurer’s key agreement from joint randomness , .
Quantum cryptography on the other hand lies in the intersection of two of the major scientific achievements of the 20th century, namely quantum physics and information theory. Various protocols for so-called quantum key agreement have been proposed (e.g., , ), and the possibility and impossibility of purification in different settings has been studied by many authors.
The goal of this paper is to derive parallels between classical and quantum key agreement and thus to show that the two paradigms are more closely related than previously recognized. These connections allow for investigating questions and solving open problems of purely classical information theory with quantum-mechanic methods. One of the consequences is that, in contrast to what was previously believed, there exists a classical counterpart to so-called bound entanglement (i.e., entanglement that cannot be purified by any quantum protocol), namely intrinsic information shared by Alice and Bob which they cannot use for generating a secret key by any classical protocol.
The outline of the paper is as follows. In Section 2 we introduce the classical (Section 2.1) and quantum (Section 2.2) models of information-theoretic key agreement and the crucial concepts and quantities, such as secret-key rate and intrinsic information on one side, and measurements, entanglement, and quantum privacy amplification on the other. In Section 3, we show the mentioned links between these two models, more precisely, between entanglement and intrinsic information (Section 3.1) as well as between quantum purification and the secret-key rate (Section 3.4). We illustrate the statements and their consequences with a number of examples (Sections 3.2 and 3.5 and Appendix B). In Section 3.6 we define and characterize the classical counterpart of bound entanglement, called bound intrinsic information. Finally we show that not only problems in classical information theory can be addressed by quantum-mechanical methods, but that the inverse is also true: In Section 3.3 we propose a new measure for entanglement based on the intrinsic information measure.
## 2 Models of Information-Theoretically Secure Key Agreement
### 2.1 Key Agreement from Classical Information: Intrinsic Information and Secret-Key Rate
In this section we describe Maurer’s general model of classical key agreement by public discussion from common information . In this setting, two parties Alice and Bob who are willing to generate a secret key have access to repeated independent realizations of (classical) random variables $`X`$ and $`Y`$, respectively, whereas an adversary Eve learns the outcomes of a random variable $`Z`$. Let $`P_{XYZ}`$ be the joint distribution of the three random variables. In addition, Alice and Bob are connected by a noiseless and authentic but otherwise completely insecure channel (see Figure 1 in Appendix A). In this situation, the secret-key rate $`S(X;Y||Z)`$ has been defined as the maximal rate at which Alice and Bob can generate a secret key that is equal for Alice and Bob with overwhelming probability and about which Eve has only a negligible amount of (Shannon) information. For a detailed discussion of the general scenario and the secret-key rate as well as for various bounds on $`S(X;Y||Z)`$, see , , .
Bound (1) implies that if Bob’s random variable $`Y`$ provides more information about Alice’s $`X`$ than Eve’s $`Z`$ does (or vice versa), then this advantage can be exploited for generating a secret key:
$$S(X;Y||Z)\mathrm{max}\{I(X;Y)I(X;Z),I(Y;X)I(Y;Z)\}.$$
(1)
This is a consequence of a result by Csiszár and Körner . It is somewhat surprising that this bound is not tight, in particular, that secret-key agreement can even be possible when the right-hand side of (1) vanishes or is negative. However, the positivity of the expression on the right-hand side of (1) is a necessary and sufficient condition for the possibility of secret-key agreement by one-way communication: whenever Alice and Bob start in a disadvantageous situation with respect to Eve, feedback is necessary. The corresponding initial phase of the key-agreement protocol is then often called advantage distillation.
The following upper bound on $`S(X;Y||Z)`$ is a generalization of Shannon’s well-known impracticality theorem and quantifies the intuitive fact that no information-theoretically secure key agreement is possible when Bob’s information is independent from Alice’s random variable, given Eve’s information: $`S(X;Y||Z)I(X;Y|Z).`$ However, this bound is not tight. Because it is a possible strategy of the adversary Eve to process $`Z`$, i.e., to send $`Z`$ over some channel characterized by $`P_{\overline{Z}|Z}`$, we have for such a new random variable $`\overline{Z}`$ that $`S(X;Y||Z)I(X;Y|\overline{Z}),`$ and hence
$$S(X;Y||Z)\underset{P_{\overline{Z}|Z}}{\mathrm{min}}\{I(X;Y|\overline{Z})\}=:I(X;YZ)$$
(2)
holds. The quantity $`I(X;YZ)`$ has been called the intrinsic conditional information between $`X`$ and $`Y`$ given $`Z`$ . It was conjectured, and evidence supporting this belief was given, that $`S(X;Y||Z)>0`$ holds if $`I(X;YZ)>0`$ does . Some of the results below strongly suggest that this is true if one of the random variables $`X`$ and $`Y`$ is binary and the other one at most ternary, but false in general.
### 2.2 Quantum Key Agreement: Measurements, Entanglement, Purification
We assume that the reader is familiar with the basic quantum-theoretic concepts and notations. For an introduction, see for example .
In the context of quantum key agreement, the classical scenario $`P_{XYZ}`$ is replaced by a quantum state vector $`\mathrm{\Psi }_A_B_E`$<sup>1</sup><sup>1</sup>1We consider pure states, since it is natural to assume that Eve controls all the environment outside Alice and Bob’s systems., where $`_A`$, $`_B`$, and $`_E`$ are Hilbert spaces describing the systems in Alice’s, Bob’s, and Eve’s hands, respectively. Then, measuring this quantum state by the three parties leads to a classical probability distribution. In the following, we assume that Eve is free to carry out so-called generalized measurements (POVMs) . In other words, the set $`\{|z\}`$ will not be assumed to be an orthonormal basis, but any set generating the Hilbert space $`_E`$ and satisfying the condition $`_z|zz|=\text{1}\text{1}__E`$. Then, if the three parties carry out measurements in certain bases<sup>2</sup><sup>2</sup>2We assume all bases to be orthonormal. $`\{|x\}`$ and $`\{|y\}`$, and in the set $`\{|z\}`$, respectively, they end up with the classical scenario $`P_{XYZ}=|x,y,z|\mathrm{\Psi }|^2`$. Since this distribution depends on the chosen bases and set, a given quantum state $`\mathrm{\Psi }`$ does not uniquely determine a classical scenario: some measurements may lead to scenarios useful for Alice and Bob, whereas for Eve, some others may (see Appendix B).
The analog of Alice and Bob’s marginal distribution $`P_{XY}`$ is the partial state $`\rho _{AB}`$, obtained by tracing over Eve’s Hilbert space $`_E`$. More precisely, let $`\mathrm{\Psi }=_{xyz}c_{xyz}|x,y,z`$, where $`|x,y,z`$ is short for $`|x|y|z`$. We can write $`\mathrm{\Psi }=_z\sqrt{P_Z(z)}\psi _z|z`$, where $`P_Z`$ denotes Eve’s marginal distribution of $`P_{XYZ}`$. Then $`\rho _{AB}=\mathrm{Tr}__E(P_\mathrm{\Psi }):=_zP_Z(z)P_{\psi _z}`$, where $`P_{\psi _z}`$ is the projector to the state vector $`\psi _z`$.
An important property is that $`\rho _{AB}`$ is pure ($`\rho _{AB}^2=\rho _{AB}`$) if and only if the global state $`\mathrm{\Psi }`$ factorizes, i.e., $`\mathrm{\Psi }=\psi _{AB}\psi _E`$, where $`\psi _{AB}_A_B`$ and $`\psi _E_E`$. In this case Alice and Bob are independent of Eve: Eve cannot obtain any information on Alice’s and Bob’s states by measuring her system.
After a measurement, Alice and Bob obtain a classical distribution $`P_{XY}`$. However, in order to obtain a well-defined classical scenario one has to assume that also Eve performs a measurement, i.e., that Eve treats her information on the classical level. Indeed, only then a classical distribution $`P_{XYZ}`$ is defined. But considering that in practice all $`P_{XYZ}`$ result from some physical process, the assumption that Eve performs the measurement one would like her to perform is not founded on basic principles<sup>3</sup><sup>3</sup>3One could argue that if the system in Eve’s hand is classical, then she has no choice for her measurement. But ultimately all systems are quantum mechanical and the apparent lack of choice might purely be a matter of technology.. For example, Eve’s measurement could be done later and depend on the public discussion between Alice and Bob. Consequently, the common approach which starts from $`P_{XYZ}`$ to prove the security of a key agreement protocol hides an assumption about Eve’s measurement. As we shall see, avoiding this hidden assumption and staying in the quantum regime can actually simplify the analysis of the scenario.
When Alice and Bob share many independent systems<sup>4</sup><sup>4</sup>4Here we do not consider the possibility that Eve coherently processes several of her systems. This corresponds to the assumption in the classical scenario that repeated realizations of $`X`$, $`Y`$, and $`Z`$ are independent of each other. $`\rho _{AB}`$, there are basically two possibilities for generating a secret key. Either they first measure their systems and then run a classical protocol (process classical information) secure against all measurements Eve could possibly perform (i.e., against all possible distributions $`P_{XYZ}`$ that can result after Eve’s measurement). Or they first run a quantum protocol (i.e., process the information in the quantum domain) and then perform their measurements. The idea of quantum protocols is to process the systems in state $`\rho _{AB}`$ and to produce fewer systems in a pure state (i.e., to purify $`\rho _{AB}`$), thus to eliminate Eve from the scenario. Moreover, the pure state Alice and Bob end up with should be maximally entangled (i.e., even for some different and incompatible measurements, Alice’s and Bob’s results are perfectly correlated). Finally, Alice and Bob measure their maximally entangled systems and establish a secret key. This way of obtaining a key directly from a quantum state $`\mathrm{\Psi }`$, without any error correction nor classical privacy amplification, is called quantum privacy amplification<sup>5</sup><sup>5</sup>5 The term “quantum privacy amplification” is somewhat unfortunate since it does not correspond to classical privacy amplification, but includes advantage distillation and error correction. (QPA for short) , . Note that the procedure described in and guarantees that Eve’s relative information (relative to the key length) is arbitrarily small, but not that her absolute information is negligible. The analog of this problem in the classical case is discussed in .
The precise conditions under which a general state $`\rho _{AB}`$ can be purified are not known. However, the two following conditions are necessary. First, the state must be entangled or, equivalently, not separable. A state $`\rho _{AB}`$ is separable if and only if it can be written as a mixture of product states, i.e., $`\rho _{AB}=_jp_j\rho _{Aj}\rho _{Bj}`$. Separable states can be generated by purely classical communication, hence it follows from bound (2) that entanglement is a necessary condition. The second condition is more subtle: The matrix $`\rho _{AB}^t`$ obtained from $`\rho _{AB}`$ by partial transposition must have at least one negative eigenvalue , . The partial transposition of the density matrix $`\rho _{AB}`$ is defined as $`(\rho _{AB})_{i,j;\mu ,\nu }^t:=(\rho _{AB})_{i,\nu ;\mu ,j}`$, where the indices $`i`$ and $`\mu `$ \[$`j`$ and $`\nu `$\] run through a basis of $`_A`$ \[$`_B`$\]. Note that this definition is base-dependent. However, the eigenvalues of $`\rho _{AB}^t`$ are not . The second of these conditions implies the first one: negative (i.e., at least one eigenvalue is negative) partial transposition implies entanglement.
In the binary case ($`_A`$ and $`_B`$ both have dimension two), the above two conditions are equivalent and sufficient for the possibility of quantum key agreement: all entangled binary states can be purified. The same even holds if one Hilbert space is of dimension 2 and the other one of dimension 3. However, for larger dimensions there are examples showing that these conditions are not equivalent: There are entangled states whose partial transpose has no negative eigenvalue, hence cannot be purified . Such states are called bound entangled, in contrast to free entangled states, which can be purified. Moreover, it is believed that there even exist entangled states which cannot be purified although they have negative partial transposition .
## 3 Linking Classical and Quantum Key Agreement
In this section we derive a close connection between the possibilities offered by classical and quantum protocols for key agreement. The intuition is as follows. As described in Section 2.2, there is a very natural connection between quantum states $`\mathrm{\Psi }`$ and classical distributions $`P_{XYZ}`$ which can be thought of as arising from $`\mathrm{\Psi }`$ by measuring in a certain basis, e.g., the standard basis<sup>6</sup><sup>6</sup>6A priori, there is no privileged basis. However, physicists often write states like $`\rho _{AB}`$ in a basis which seems to be more natural than others. We refer to this as the standard basis. Somewhat surprisingly, this basis is generally easy to identify, though not precisely defined. One could characterize the standard basis as the basis for which as many coefficients as possible of $`\mathrm{\Psi }`$ are real and positive. We usually represent quantum states with respect to the standard basis.. (Note however that the connection is not unique even for fixed bases: For a given distribution $`P_{XYZ}`$, there are many states $`\mathrm{\Psi }`$ leading to $`P_{XYZ}`$ by carrying out measurements.) When given a state $`\mathrm{\Psi }`$ between three parties Alice, Bob, and Eve, and if $`\rho _{AB}`$ denotes the resulting mixed state after tracing out Eve, then the corresponding classical distribution $`P_{XYZ}`$ has positive intrinsic information if and only if $`\rho _{AB}`$ is entangled. However, this correspondence clearly depends on the measurement bases used by Alice, Bob, and Eve. If for instance $`\rho _{AB}`$ is entangled, but Alice and Bob do very unclever measurements, then the intrinsic information may vanish (see Example 7 in Appendix B). If on the other hand $`\rho _{AB}`$ is separable, Eve may do such bad measurements that the intrinsic information becomes positive, despite the fact that $`\rho _{AB}`$ could have been established by public discussion without any prior correlation (see Example 6 in Appendix B). Consequently, the correspondence between intrinsic information and entanglement must involve some optimization over all possible measurements on all sides.
A similar correspondence on the protocol level is supported by many examples, but not rigorously proven: The distribution $`P_{XYZ}`$ allows for classical key agreement if and only if quantum key agreement is possible starting from the state $`\rho _{AB}`$.
We show how these parallels allow for addressing problems of purely classical information-theoretic nature with the methods of quantum information theory, and vice versa.
### 3.1 Entanglement and Intrinsic Information
Let us first establish the connection between intrinsic information and entanglement. Theorem 1 states that if $`\rho _{AB}`$ is separable, then Eve can “force” the information between Alice’s and Bob’s classical random variables (given Eve’s classical random variable) to be zero (whatever strategy Alice and Bob use). In particular, Eve can prevent classical key agreement.
###### Theorem 1
Let $`\mathrm{\Psi }_A_B_E`$ and $`\rho _{AB}=\mathrm{Tr}__E(P_\mathrm{\Psi })`$. If $`\rho _{AB}`$ is separable, then there exists a generating set $`\{|z\}`$ of $`_{}`$ such that for all bases $`\{|x\}`$ and $`\{|y\}`$ of $`_𝒜`$ and $`_{}`$, respectively, $`I(X;Y|Z)=0`$ holds for $`P_{XYZ}(x,y,z):=|x,y,z|\mathrm{\Psi }|^2`$.
Proof. If $`\rho _{AB}`$ is separable, then there exist vectors $`|\alpha _z`$ and $`|\beta _z`$ such that $`\rho _{AB}=_{z=1}^{n_z}p_zP_{\alpha _z}P_{\beta _z}`$, where $`P_{\alpha _z}`$ denotes the one-dimensional projector onto the subspace spanned by $`|\alpha _z`$.
Let us first assume that $`n_zdim_E`$. Then there exists a basis $`\{|z\}`$ of $`_E`$ such that $`\mathrm{\Psi }=_z\sqrt{p_z}|\alpha _z,\beta _z,z`$ holds , , .
If $`n_z>dim_E`$, then Eve can add an auxiliary system $`_{aux}`$ to hers (usually called an ancilla) and we have $`\mathrm{\Psi }|\gamma _0=_z\sqrt{p_z}|\alpha _z,\beta _z,\gamma _z`$, where $`|\gamma _0_{aux}`$ is the state of Eve’s auxiliary system, and $`\{|\gamma _z\}`$ is a basis of $`_E_{aux}`$. We define the (not necessarily orthonormalized) vectors $`|z`$ by $`|z,\gamma _0=\text{1}\text{1}__EP_{\gamma _0}|\gamma _z`$. These vectors determine a generalized measurement with positive operators $`O_z=|zz|`$. Since $`_zO_zP_{\gamma _0}=_z|z,\gamma _0z,\gamma _0|=_z\text{1}\text{1}__EP_{\gamma _0}|\gamma _z|\gamma _z|\text{1}\text{1}__EP_{\gamma _0}=\text{1}\text{1}__EP_{\gamma _0}`$, the $`O_z`$ satisfy $`_zO_z=\text{1}\text{1}__E`$, as they should in order to define a generalized measurement . Note that the first case ($`n_zdim_E`$) is a special case of the second one, with $`|\gamma _z=|z,\gamma _0`$. If Eve now performs the measurement, then we have $`P_{XYZ}(x,y,z)=|x,y,z|\mathrm{\Psi }|^2=|x,y,\gamma _z|\mathrm{\Psi },\gamma _0|^2`$, and
$$P_{XY|Z}(x,y,z)=|x,y|\alpha _z,\beta _z|^2=|x|\alpha _z|^2|y|\beta _z|^2=P_{X|Z}(x,z)P_{Y|Z}(y,z)$$
holds for all $`|z`$ and for all $`|x,y_𝒜_{}`$. Consequently, $`I(X;Y|Z)=0`$. $`\mathrm{}`$
Theorem 2 states that if $`\rho _{AB}`$ is entangled, then Eve cannot force the intrinsic information to be zero: Whatever she does (i.e., whatever generalized measurements she carries out), there is something Alice and Bob can do such that the intrinsic information is positive. Note that this does not, a priori, imply that secret-key agreement is possible in every case. Indeed, we will provide evidence for the fact that this implication does generally not hold.
###### Theorem 2
Let $`\mathrm{\Psi }_A_B_E`$ and $`\rho _{AB}=\mathrm{Tr}__E(P_\mathrm{\Psi })`$. If $`\rho _{AB}`$ is entangled, then for all generating sets $`\{|z\}`$ of $`_{}`$, there are bases $`\{|x\}`$ and $`\{|y\}`$ of $`_𝒜`$ and $`_{}`$, respectively, such that $`I(X;YZ)>0`$ holds for $`P_{XYZ}(x,y,z):=|x,y,z|\mathrm{\Psi }|^2`$.
Proof. We prove this by contradiction. Assume that there exists a generating set $`\{|z\}`$ of $`_E`$ such that for all bases $`\{|x\}`$ of $`_A`$ and $`\{|y\}`$ of $`_B`$, $`I(X;YZ)=0`$ holds for the resulting distribution. For such a distribution, there exists a channel, characterized by $`P_{\overline{Z}|Z}`$, such that $`I(X;Y|\overline{Z})=0`$ holds, i.e.,
$$P_{XY|\overline{Z}}(x,y,\overline{z})=P_{X|\overline{Z}}(x,\overline{z})P_{Y|\overline{Z}}(y,\overline{z}).$$
(3)
Let now $`\rho _{\overline{z}}:=(1/p_{\overline{z}})_zp_zP_{\overline{Z}|Z}(\overline{z},z)P_{\psi _z}`$ with $`p_z=P_Z(z)`$ and $`p_{\overline{z}}=_zP_{\overline{Z}|Z}(\overline{z},z)p_z,`$ and where $`\psi _z`$ is the state of Alice’s and Bob’s system conditioned on Eve’s result $`z`$: $`\mathrm{\Psi }|\gamma _0=_z\psi _z|\gamma _z`$ (see the proof of Theorem 1).
From (3) we can conclude $`\mathrm{Tr}(P_xP_y\rho _{\overline{z}})=\mathrm{Tr}(P_x\text{1}\text{1}\rho _{\overline{z}})\mathrm{Tr}(\text{1}\text{1}P_y\rho _{\overline{z}})`$ for all one-dimensional projectors $`P_x`$ and $`P_y`$ acting in $`_𝒜`$ and $`_{}`$, respectively. Consequently, the states $`\rho _{\overline{z}}`$ are products, i.e., $`\rho _{\overline{z}}=\rho _{\alpha _{\overline{z}}}\rho _{\beta _{\overline{z}}},`$ and $`\rho _{AB}=_{\overline{z}}p_{\overline{z}}\rho _{\overline{z}}`$ is separable. $`\mathrm{}`$
Theorem 2 can be formulated in a more positive way. Let us first introduce the concept of a set of bases $`\{|x^j,|y^j\}`$, where the $`j`$ label the different bases, as they are used in the 4-state (2 bases) and the 6-state (3 bases) protocols , , . Then if $`\rho _{AB}`$ is entangled there exists a set $`\{|x^j,|y^j\}`$ of $`N`$ bases such that for all generalized measurements $`\{|z\}`$, $`I(X;Y[Z,j])>0`$ holds. The idea is that Alice and Bob randomly choose a basis and, after the transmission, publicly restrict to the (possibly few) cases where they happen to have chosen the same basis. Hence Eve knows $`j`$, and one has $`I(X;Y[Z,j])=(1/N)_{j=1}^NI(X^j;Y^jZ).`$ If the set of bases is large enough, then for all $`\{|z\}`$ there is a basis with positive intrinsic information, hence the mean is also positive. Clearly, this result is stronger if the set of bases is small. Nothing is proven about the achievable size of such sets of bases, but it is conceivable that $`\mathrm{max}\{dim_A,dim_B\}`$ bases are always sufficient.
###### Corollary 3
Let $`\mathrm{\Psi }_A_B_E`$ and $`\rho _{AB}=\mathrm{Tr}__E(P_\mathrm{\Psi })`$. Then the following statements are equivalent:
(i) $`\rho _{AB}`$ is entangled,
(ii) for all generating sets $`\{|z\}`$ of $`_E`$, there exist bases $`\{|x\}`$ of $`_A`$ and $`\{|y\}`$ of $`_B`$ such that the distribution $`P_{XYZ}(x,y,z):=|x,y,z|\mathrm{\Psi }|^2`$ satisfies $`I(X;YZ)>0`$,
(ii) for all generating sets $`\{|z\}`$ of $`_E`$, there exist bases $`\{|x\}`$ of $`_A`$ and $`\{|y\}`$ of $`_B`$ such that the distribution $`P_{XYZ}(x,y,z):=|x,y,z|\mathrm{\Psi }|^2`$ satisfies $`I(X;Y|Z)>0`$.
A first consequence of the fact that Corollary 3 often holds with respect to the standard bases (see below) is that it yields, at least in the binary case, a criterion for $`I(X;YZ)>0`$ that is efficiently verifiable since it is based on the positivity of the eigenvalues of a $`4\times 4`$ matrix (see also Example 5). Previously, the quantity $`I(X;YZ)`$ has been considered to be quite hard to handle.
### 3.2 Examples I
The following examples illustrate the correspondence established in Section 3.1. They show in particular that very often (Examples 1, 2, and 3), but not always (Examples 6 and 7 in Appendix B), the direct connection between entanglement and positive intrinsic information holds with respect to the standard bases (i.e., the bases physicists use by commodity and intuition).
Example 1. Let us consider the so-called 4-state protocol of . The analysis of the 6-state protocol is analogous and leads to similar results . We compare the possibility of quantum and classical key agreement given the quantum state and the corresponding classical distribution, respectively, arising from this protocol. The conclusion is, under the assumption of incoherent eavesdropping, that key agreement in one setting is possible if and only if this is true also for the other.
After carrying out the 4-state protocol, and under the assumption of optimal eavesdropping (in terms of Shannon information), the resulting quantum state is
$$\mathrm{\Psi }=\sqrt{F/2}|0,0\xi _{00}+\sqrt{D/2}|0,1\xi _{01}+\sqrt{D/2}|1,0\xi _{10}+\sqrt{F/2}|1,1\xi _{11}𝐂^2𝐂^2𝐂^4,$$
where $`D`$ (the disturbance) is the probability that $`XY`$ holds if $`X`$ and $`Y`$ are the classical random variables of Alice and Bob, respectively, where $`F=1D`$ (the fidelity), and where the $`\xi _{ij}`$ satisfy $`\xi _{00}|\xi _{11}=\xi _{01}|\xi _{10}=12D`$ and $`\xi _{ii}|\xi _{ij}=0`$ for all $`ij`$. Then the state $`\rho _{AB}`$ is (in the basis $`\{|\mathrm{\hspace{0.17em}00}`$,$`|\mathrm{\hspace{0.17em}01}`$,$`|\mathrm{\hspace{0.17em}10}`$,$`|\mathrm{\hspace{0.17em}11}\}`$)
$$\rho _{AB}=\frac{1}{2}\left(\begin{array}{cccc}D& 0& 0& D\left(12D\right)\\ 0& 1D& \left(1D\right)\left(12D\right)& 0\\ 0& \left(1D\right)\left(12D\right)& 1D& 0\\ D\left(12D\right)& 0& 0& D\end{array}\right),$$
and its partial transpose
$$\rho _{AB}^t=\frac{1}{2}\left(\begin{array}{cccc}D& 0& 0& \left(1D\right)\left(12D\right)\\ 0& 1D& D\left(12D\right)& 0\\ 0& D\left(12D\right)& 1D& 0\\ \left(1D\right)\left(12D\right)& 0& 0& D\end{array}\right)$$
has the eigenvalues $`(1/2)(D\pm (1D)(12D))`$ and $`(1/2)((1D)\pm D(12D))`$, which are all non-negative (i.e., $`\rho _{AB}`$ is separable) if
$$D1\frac{1}{\sqrt{2}}.$$
(4)
From the classical viewpoint, the corresponding distributions (arising from measuring the above quantum system in the standard bases) are as follows. First, $`X`$ and $`Y`$ are both symmetric bits with $`\mathrm{Prob}[XY]=D`$. Eve’s random variable $`Z=[Z_1,Z_2]`$ is composed of 2 bits $`Z_1`$ and $`Z_2`$, where $`Z_1=XY`$, i.e., $`Z_1`$ tells Eve whether Bob received the qubit disturbed ($`Z_1=1`$) or not ($`Z_1=0`$) (this is a consequence of the fact that the $`\xi _{ii}`$ and $`\xi _{ij}`$ ($`ij`$) states generate orthogonal sub-spaces), and where the probability that Eve’s second bit indicates the correct value of Bob’s bit is Prob$`[Z_2=Y]=\delta =(1+\sqrt{1\xi _{00}|\xi _{11}^2})/2=1/2+\sqrt{D(1D)}`$. We now prove that for this distribution, the intrinsic information is zero if and only if
$$\frac{D}{1D}2\sqrt{(1\delta )\delta }=12D$$
(5)
holds. We show that if the condition (5) is satisfied, then $`I(X;YZ)=0`$ holds. The inverse implication follows from the existence of a key-agreement protocol in all other cases (see Example 1 (cont’d) in Section 3.5). If (5) holds, we can construct a random variable $`\overline{Z}`$, that is generated by sending $`Z`$ over a channel characterized by $`P_{\overline{Z}|Z}`$, for which $`I(X;Y|\overline{Z})=0`$ holds. We can restrict ourselves to the case of equality in (5) because Eve can always increase $`\delta `$ by adding noise.
Consider now the channel characterized by the following distribution $`P_{\overline{Z}|Z}`$ (where $`\overline{𝒵}=\{u,v\}`$): $`P_{\overline{Z}|Z}(u,[0,0])=P_{\overline{Z}|Z}(v,[0,1])=1`$, $`P_{\overline{Z}|Z}(l,[1,0])=P_{\overline{Z}|Z}(l,[1,1])=1/2`$ for $`l\{u,v\}`$. (The channel $`P_{\overline{Z}|Z}`$ is illustrated in Figure 2 in Appendix A.) We show $`I(X;Y|\overline{Z})=\mathrm{E}_{\overline{Z}}[I(X;Y|\overline{Z}=\overline{z})]=0`$, i.e., that $`I(X;Y|\overline{Z}=u)=0`$ and $`I(X;Y|\overline{Z}=v)=0`$ hold. By symmetry it is sufficient to show the first equality. For $`a_{ij}:=P_{XY\overline{Z}}(i,j,u)`$, we get
$$a_{00}=(1D)(1\delta )/2,a_{11}=(1D)\delta /2,a_{01}=a_{10}=(D(1\delta )/2+D\delta /2)/2=D/4.$$
From equality in (5) we conclude $`a_{00}a_{11}=a_{01}a_{10}`$, which is equivalent to the fact that $`X`$ and $`Y`$ are independent, given $`\overline{Z}=u`$.
Finally, note that the conditions (4) and (5) are equivalent for $`D[0,1/2]`$. This shows that the bounds of tolerable noise are indeed exactly the same for the quantum and classical scenarios. $`\mathrm{}`$
Example 2. We consider the bound entangled state presented in . This example received quite a lot of attention by the quantum-information community because it was the first known example of bound entanglement (i.e., entanglement without the possibility of quantum key agreement). We show that its classical counterpart seems to have similarly surprising properties. Let $`0<a<1`$ and
$$\mathrm{\Psi }=\sqrt{\frac{3a}{8a+1}}\psi |0+\sqrt{\frac{1}{8a+1}}\varphi _a|1+\sqrt{\frac{a}{8a+1}}(|122+|133+|214+|235+|326),$$
where $`\psi =(|11+|22+|33)/\sqrt{3}`$ and $`\varphi _a=\sqrt{(1+a)/(2)}|31+\sqrt{(1a)/(2)}|33.`$ It has been shown in that the resulting state $`\rho _{AB}`$ is entangled.
The corresponding classical distribution is as follows. The ranges are $`𝒳=𝒴=\{1,2,3\}`$ and $`𝒵=\{0,1,2,3,4,5,6\}`$. We write $`(ijk)=P_{XYZ}(i,j,k)`$. Then we have $`(110)=(220)=(330)=(122)=(133)=(214)=(235)=(326)=2a/(16a+2)`$, $`(311)=(1+a)/(16a+2)`$, and $`(331)=(1a)/(16a+2)`$. We study the special case $`a=1/2`$. Consider the following representation of the resulting distribution (to be normalized). For instance, the entry “$`(0)1,(1)1/2`$” for $`X=Y=3`$ means $`P_{XYZ}(3,3,0)=1/10`$ (normalized), $`P_{XYZ}(3,3,1)=1/20`$, and $`P_{XYZ}(3,3,z)=0`$ for all $`z\{0,1\}`$.
| X | 1 | 2 | 3 |
| --- | --- | --- | --- |
| Y (Z) | | | |
| 1 | (0) 1 | (4) 1 | (1) 3/2 |
| 2 | (2) 1 | (0) 1 | (6) 1 |
| 3 | (3) 1 | (5) 1 | (0) 1 |
| | | | (1) 1/2 |
As we would expect, the intrinsic information is positive in this scenario. This can be seen by contradiction as follows. Assume $`I(X;YZ)=0`$. Hence there exists a discrete channel, characterized by the conditional distribution $`P_{\overline{Z}|Z}`$, such that $`I(X;Y|\overline{Z})=0`$ holds. Let $`\overline{𝒵}𝐍`$, and let $`P_{\overline{Z}|Z}(i,0)=:a_i`$, $`P_{\overline{Z}|Z}(i,1)=:x_i`$, $`P_{\overline{Z}|Z}(i,6)=:s_i`$. Then we must have $`a_i,x_i,s_i[0,1]`$ and $`_ia_i=_ix_i=_is_i=1`$. Using $`I(X;Y|\overline{Z})=0`$, we obtain the following distributions $`P_{XY|\overline{Z}=i}`$ (to be normalized):
| X | 1 | 2 | 3 |
| --- | --- | --- | --- |
| Y | | | |
| 1 | $`a_i`$ | $`\frac{3a_ix_i}{2s_i}`$ | $`\frac{3x_i}{2}`$ |
| 2 | $`\frac{2a_is_i}{3x_i}`$ | $`a_i`$ | $`s_i`$ |
| 3 | $`\frac{2a_i(a_i+x_i/2)}{3x_i}`$ | $`\frac{a_i(a_i+x_i/2)}{s_i}`$ | $`a_i+\frac{x_i}{2}`$ |
By comparing the $`(2,3)`$-entries of the two tables above, we obtain
$$1\underset{i}{}\frac{a_i(a_i+x_i/2)}{s_i}.$$
(6)
We now prove that (6) implies $`s_ia_i`$ (i.e., $`s_i=a_i`$ for all $`i`$) and $`x_i0`$. Clearly, this does not lead to a solution and is hence a contradiction. For instance, $`P_{XY|\overline{Z}=i}(1,2)=2a_is_i/(3x_i)`$ is not even defined in this case if $`a_i>0`$.
It remains to show that (6) implies $`a_is_i`$ and $`x_i0`$. We show that whenever $`_ia_i=_is_i=1`$ and $`a_is_i`$, then $`_ia_i^2/s_i>1.`$ First, note that $`_ia_i^2/s_i=_ia_i=1`$ for $`a_is_i`$. Let now $`s_{i_1}a_{i_1}`$ and $`s_{i_2}a_{i_2}`$. We show that $`a_{i_1}^2/s_{i_1}+a_{i_2}^2/s_{i_2}<a_{i_1}^2/(s_{i_1}\epsilon )+a_{i_2}^2/(s_{i_2}+\epsilon )`$ holds for every $`\epsilon >0`$, which obviously implies the above statement. It is straightforward to see that this is equivalent to $`a_{i_1}^2s_{i_2}(s_{i_2}+\epsilon )>a_{i_2}^2s_{i_1}(s_{i_1}\epsilon ),`$ and holds because of $`a_{i_1}^2s_{i_2}(s_{i_2}+\epsilon )>a_{i_1}^2a_{i_2}^2`$ and $`a_{i_2}^2s_{i_1}(s_{i_1}\epsilon )<a_{i_1}^2a_{i_2}^2`$. This concludes the proof of $`I(X;YZ)>0`$. $`\mathrm{}`$
As mentioned, the interesting point about Example 2 is that the quantum state is bound entangled, and that also classical key agreement seems impossible despite the fact that $`I(X;YZ)>0`$ holds. This is a contradiction to a conjecture stated in . The classical translation of the bound entangled state leads to a classical distribution with very strange properties as well! (See Example 2 (cont’d) in Section 3.4).
In Example 3, another bound entangled state (first proposed in ) is discussed. The example is particularly nice because, depending on the choice of the parameter $`\alpha `$, the quantum state can be made separable, bound entangled, and free entangled.
Example 3. We consider the following distribution (to be normalized). Let $`2\alpha 5`$.
| X | 1 | 2 | 3 |
| --- | --- | --- | --- |
| Y (Z) | | | |
| 1 | (0) 2 | (4) $`5\alpha `$ | (3) $`\alpha `$ |
| 2 | (1) $`\alpha `$ | (0) 2 | (5) $`5\alpha `$ |
| 3 | (6) $`5\alpha `$ | (2) $`\alpha `$ | (0) 2 |
This distribution arises when measuring the following quantum state. Let $`\psi :=(1/\sqrt{3})(|11+|22+|33).`$ Then
$$\mathrm{\Psi }=\sqrt{\frac{2}{7}}\psi |0+\sqrt{\frac{a}{21}}(|12|1+|23|2+|31|3)+\sqrt{\frac{5a}{21}}(|21|4+|32|5+|13|6),$$
$$\text{and }\rho _{AB}=\frac{2}{7}P_\psi +\frac{a}{21}(P_{12}+P_{23}+P_{31})+\frac{5a}{21}(P_{21}+P_{32}+P_{13})$$
is separable if and only if $`\alpha [2,3]`$, bound entangled for $`\alpha (3,4]`$, and free entangled if $`\alpha (4,5]`$ (see Figure 3 in Appendix A).
Let us consider the quantity $`I(X;YZ)`$. First of all, it is clear that $`I(X;YZ)=0`$ holds for $`\alpha [2,3]`$. The reason is that $`\alpha 2`$ and $`5\alpha 2`$ together imply that Eve can “mix” her symbol $`Z=0`$ with the remaining symbols in such a way that when given that $`\overline{Z}`$ takes the “mixed value,” then $`XY`$ is uniformly distributed; in particular, $`X`$ and $`Y`$ are independent. Moreover, it can be shown in analogy to Example 2 that $`I(X;YZ)>0`$ holds for $`\alpha >3`$. $`\mathrm{}`$
Examples 1, 2, and 3 suggest that the correspondence between separability and entanglement on one side and vanishing and non-vanishing intrinsic information on the other always holds with respect to the standard bases or even arbitrary bases. We show in Appendix B that this is not true in general. More precisely, Examples 6 and 7 demonstrate how Eve as well as Alice and Bob can perform bad measurements. Hence the parallelity between the quantum and classical situation must be as it is stated in Theorems 1 and 2.
### 3.3 A Classical Measure for Quantum Entanglement
It is a challenging problem of theoretical quantum physics to find good measures for entanglement . Corollary 3 above suggests the following measure, which is based on classical information theory.
###### Definition 1
Let for a quantum state $`\rho _{AB}`$
$$\mu (\rho _{AB}):=\underset{\{|z\}}{\mathrm{min}}(\underset{\{|x\},\{|y\}}{\mathrm{max}}(I(X;YZ))),$$
where the minimum is taken over all $`\mathrm{\Psi }=_z\sqrt{p_z}\psi _z|z`$ such that $`\rho _{AB}=\mathrm{Tr}__E(P_\mathrm{\Psi })`$ holds and over all bases $`\{|z\}`$ of $`_E`$, the maximum is over all bases $`\{|x\}`$ of $`_A`$ and $`\{|y\}`$ of $`_B`$, and where $`P_{XYZ}(x,y,z):=|x,y,z|\mathrm{\Psi }|^2`$.
Then $`\mu `$ has all the properties required from such a measure. If $`\rho _{AB}`$ is pure, i.e., $`\rho _{AB}=|\psi _{AB}\psi _{AB}|`$, then we have in the Schmidt basis (see for example ) $`\psi _{AB}=_jc_j|x_j,y_j`$, and $`\mu (\rho _{AB})=_j|c_j|^2\mathrm{log}_2(|c_j|^2)=\mathrm{Tr}(\rho _{AB}\mathrm{log}_2\rho _{AB})`$, as it should . It is obvious that $`\mu `$ is convex, i.e., $`\mu (\lambda \rho _1+(1\lambda )\rho _2)\lambda \mu (\rho _1)+(1\lambda )\mu (\rho _2)`$.
Example 4 (based on Werner’s states). Let $`\mathrm{\Psi }=\sqrt{\lambda }\psi ^{()}|0+\sqrt{(1\lambda )/4}|001+012+103+114`$, where $`\psi ^{()}=|1001/\sqrt{2}`$, and $`\rho _{AB}=\lambda P_{\psi ^{()}}+((1\lambda )/4)\text{1}\text{1}`$. It is well-known that $`\rho _{AB}`$ is separable if and only if $`\lambda 1/3`$. Then the classical distribution is $`(010)=(100)=\lambda /2`$ and $`(001)=(012)=(103)=(114)=(1\lambda )/4`$.
If $`\lambda 1/3`$, then consider the channel $`P_{\overline{Z}|Z}(0,0)=P_{\overline{Z}|Z}(2,2)=P_{\overline{Z}|Z}(3,3)=1,P_{\overline{Z}|Z}(0,1)=P_{\overline{Z}|Z}(0,4)=\xi ,P_{\overline{Z}|Z}(1,1)=P_{\overline{Z}|Z}(4,4)=1\xi ,`$ where $`\xi =2\lambda /(1\lambda )1`$. Then $`\mu (\rho _{AB})=I(X;YZ)=I(X;Y|\overline{Z})=0`$ holds, as it should.
If $`\lambda >1/3`$, then consider the (obviously optimal) channel $`P_{\overline{Z}|Z}(0,0)=P_{\overline{Z}|Z}(2,2)=P_{\overline{Z}|Z}(3,3)=P_{\overline{Z}|Z}(0,1)=P_{\overline{Z}|Z}(0,4)=1`$. Then
$$\mu \left(\rho _{AB}\right)=I\left(X;YZ\right)=I(X;Y|\overline{Z})=P_{\overline{Z}}\left(0\right)I\left(X;Y|\overline{Z}=0\right)=\frac{1+\lambda }{2}\left(1q\mathrm{log}_2q\left(1q\right)\mathrm{log}_2\left(1q\right)\right),$$
where $`q=2\lambda /(1+\lambda )`$. $`\mathrm{}`$
### 3.4 Classical Protocols and Quantum Privacy Amplification
It is a natural question whether the analogy between entanglement and intrinsic information (see Section 3.1) carries over to the protocol level. The examples given in Section 3.5 support this belief. A quite interesting and surprising consequence would be that there exists a classical counterpart to bound entanglement, namely intrinsic information that cannot be distilled into a secret key by any classical protocol, if $`|𝒳|+|𝒴|>5`$. In other words, the conjecture in that such information can always be distilled would be proved for $`|𝒳|+|𝒴|5`$, but disproved otherwise.
###### Conjecture 1
Let $`\mathrm{\Psi }_A_B_E`$ and $`\rho _{AB}=\mathrm{Tr}__E(P_\mathrm{\Psi })`$. Assume that for all generating sets $`\{|z\}`$ of $`_E`$ there are bases $`\{|x\}`$ and $`\{|y\}`$ of $`_A`$ and $`_B`$, respectively, such that $`S(X;Y||Z)>0`$ holds for the distribution $`P_{XYZ}(x,y,z):=|x,y,z|\mathrm{\Psi }|^2`$. Then quantum privacy amplification is possible with the state $`\rho _{AB}`$, i.e., $`\rho _{AB}`$ is free entangled.
###### Conjecture 2
Let $`\mathrm{\Psi }_A_B_E`$ and $`\rho _{AB}=\mathrm{Tr}__E(P_\mathrm{\Psi })`$. Assume that there exists a generating set $`\{|z\}`$ of $`_E`$ such that for all bases $`\{|x\}`$ and $`\{|y\}`$ of $`_A`$ and $`_B`$, respectively, $`S(X;Y||Z)=0`$ holds for the distribution $`P_{XYZ}(x,y,z):=|x,y,z|\mathrm{\Psi }|^2`$. Then quantum privacy amplification is impossible with the state $`\rho _{AB}`$, i.e., $`\rho _{AB}`$ is bound entangled or separable.
### 3.5 Examples II
The following examples support Conjectures 1 and 2 and illustrate their consequences. We consider mainly the same distributions as in Section 3.2, but this time under the aspect of the existence of classical and quantum key-agreement protocols.
Example 1 (cont’d). We have shown in Section 3.2 that the resulting quantum state is entangled if and only if the intrinsic information of the corresponding classical situation (with respect to the standard bases) is non-zero. Here, we show that such a correspondence also holds on the protocol level. First of all, it is clear for the quantum state that QPA is possible whenever the state is entangled because both $`_A`$ and $`_B`$ have dimension two. On the other hand, the same is also true for the corresponding classical situation, i.e., secret-key agreement is possible whenever $`D/(1D)<2\sqrt{(1\delta )\delta }`$ holds, i.e., if the intrinsic information is positive. This is shown in Appendix C. There we describe the required protocol, more precisely, the advantage-distillation phase (called repeat-code protocol ), in which Alice and Bob use their advantage given by the authenticity of the public-discussion channel for generating new random variables for which the legitimate partners have an advantage over Eve in terms of the (Shannon) information about each other’s new random variables. For a further discussion of this example, see also . $`\mathrm{}`$
Example 2 (cont’d). The quantum state $`\rho _{AB}`$ in this example is bound entangled, meaning that the entanglement cannot be used for QPA. Interestingly, but not surprisingly given the discussion above, the corresponding classical distribution has the property that $`I(X;YZ)>0`$, but nevertheless, all the known classical advantage-distillation protocols , fail for this distribution! It seems that $`S(X;Y||Z)=0`$ holds (although it is not clear how this fact could be rigorously proven, except by proving Conjecture 1 directly). $`\mathrm{}`$
Example 3 (cont’d). We have seen already that for $`2\alpha 3`$, the quantum state is separable and the corresponding classical distribution (with respect to the standard bases) has vanishing intrinsic information. Moreover, it has been shown that for the quantum situation, $`3<\alpha 4`$ corresponds to bound entanglement, whereas for $`\alpha >4`$, QPA is possible and allows for generating a secret key . We describe a classical protocol here which suggests that the situation for the classical translation of the scenario is totally analogous: The protocol allows classical key agreement exactly for $`\alpha >4`$. However, this does not imply (although it appears very plausible) that no classical protocol exists at all for the case $`\alpha 4`$.
Let $`\alpha >4`$. We consider the following protocol for classical key agreement. First of all, Alice and Bob both restrict their ranges to $`\{1,2\}`$ (i.e., publicly reject a realization unless $`X\{1,2\}`$ and $`Y\{1,2\}`$). The resulting distribution is as follows (to be normalized):
| X | 1 | 2 |
| --- | --- | --- |
| Y (Z) | | |
| 1 | (0) 2 | (4) $`5\alpha `$ |
| 2 | (2) $`\alpha `$ | (0) 2 |
Then, Alice and Bob both send their bits locally over channels $`P_{\overline{X}|X}`$ and $`P_{\overline{Y}|Y}`$, respectively, such that the resulting bits $`\overline{X}`$ and $`\overline{Y}`$ are symmetric. The channel $`P_{\overline{X}|X}`$ \[$`P_{\overline{Y}|Y}`$\] sends $`X=0`$ \[$`Y=1`$\] to $`\overline{X}=1`$ \[$`\overline{Y}=0`$\] with probability $`(2\alpha 5)/(2\alpha +4)`$, and leaves $`X`$ \[$`Y`$\] unchanged otherwise. The distribution $`P_{\overline{X}\overline{Y}Z}`$ is then
| $`\overline{X}`$ | 1 | 2 |
| --- | --- | --- |
| $`\overline{Y}`$ (Z) | | |
| | (0) $`2\frac{9}{2\alpha +4}`$ | (1) $`5\alpha `$ |
| 1 | (2) $`\alpha \frac{9}{2\alpha +4}\frac{2\alpha 5}{2\alpha +4}`$ | (2) $`\alpha \left(\frac{2\alpha 5}{2\alpha +4}\right)^2`$ |
| | | (0) $`22\frac{2\alpha 5}{2\alpha +4}`$ |
| 2 | (2) $`\alpha \left(\frac{9}{2\alpha +4}\right)^2`$ | (0) $`2\frac{9}{2\alpha +4}`$ |
| | | (2) $`\alpha \frac{9}{2\alpha +4}\frac{2\alpha 5}{2\alpha +4}`$ |
It is not difficult to see that for $`\alpha >4`$, we have $`\mathrm{Prob}[\overline{X}=\overline{Y}]>1/2`$ and that, given that $`\overline{X}=\overline{Y}`$ holds, Eve has no information at all about what this bit is. Thus the repeat-code protocol described in Appendix C allows for classical key agreement in this situation. For $`\alpha 4`$ however, classical key agreement, like quantum key agreement, seems impossible. The results of Example 3 are illustrated in Figure 3 in Appendix A. $`\mathrm{}`$
Example 5. The following distribution $`P_{XYZ}`$, with binary $`X`$ and $`Y`$, was discussed and analyzed in as an example of a simple distribution for which the equivalence of $`I(X;YZ)>0`$ and $`S(X;Y||Z)>0`$ could not be shown.
Assume that the random variables $`X`$ and $`Y`$ are distributed according to $`P_{XY}(0,0)=P_{XY}(1,1)=(1\alpha )/2`$, $`P_{XY}(0,1)=P_{XY}(1,0)=\alpha /2`$, and $`Z=[Z_X,Z_Y]`$, where $`Z_X`$ and $`Z_Y`$ are generated by sending $`X`$ and $`Y`$ over two independent binary erasure channels with erasure probabilities $`\delta _X`$ and $`\delta _Y`$, respectively.
If the conjectured parallels between classical and quantum protocols hold, then $`I(X;YZ)`$ $`>0`$ implies $`S(X;Y||Z)>0`$ because both $`X`$ and $`Y`$ are binary. Moreover, due to the proven connection between intrinsic information and entanglement and hence to the eigenvalues of the partial transpose of the density matrix, the condition for $`I(X;YZ)>0`$ can be explicitly given, and is very simple. This is surprising since the determination of $`I(X;YZ)`$, as well as the advantage-distillation protocols for this distribution, turned out to be quite complicated . The condition under which all the eigenvalues of the partial transpose of the density matrix of the corresponding quantum state are non-negative is
$$(\alpha \alpha ^2)\left(\frac{(1\delta _X)(1\delta _Y)}{\delta _X}+2\right)\left(\frac{(1\delta _X)(1\delta _Y)}{\delta _Y}+2\right)1.$$
This bound is compatible with (but stronger than) all the bounds painfully derived, by working purely in the classical information-theoretic world, in . $`\mathrm{}`$
### 3.6 Bound Intrinsic Information
Examples 2 and 3 suggest that, in analogy to bound entanglement of a quantum state, bound classical information exists, i.e., conditional intrinsic information which cannot be used to generate a secret key in the classical scenario. We give a formal definition of bound intrinsic information.
###### Definition 2
Let $`P_{XYZ}`$ be a distribution with $`I(X;YZ)>0`$. Then if $`S(X;Y||Z)>0`$ holds for this distribution, the intrinsic information between $`X`$ and $`Y`$, given $`Z`$, is called free. Otherwise, if $`S(X;Y||Z)=0`$, the information is called bound.
Note that the existence of bound intrinsic information could not be proven so far. However, all known examples of bound entanglement, combined with all known advantage-distillation protocols, do not lead to a contradiction to Conjecture 1! Clearly, it would be very interesting to rigorously prove this conjecture because then, all pessimistic results known for the quantum scenario would immediately carry over to the classical setting (where such results appear to be much harder to prove).
Examples 2 and 3 also illustrate nicely what the nature of such bound information is. Of course, $`I(X;YZ)>0`$ implies both $`I(X;Y)>0`$ and $`I(X;Y|Z)>0`$. However, if $`|𝒳|+|𝒴|>5`$, it is possible that the dependence between $`X`$ and $`Y`$ and the dependence between $`X`$ and $`Y`$, given $`\overline{Z}`$, are “orthogonal.” By the latter we mean that for all fixed (deterministic or probabilistic) functions $`f:𝒳\{0,1\}`$ and $`g:𝒴\{0,1\}`$ for which the correlation of $`f(X)`$ and $`g(Y)`$ is positive, i.e.,
$$P_{f(X)g(Y)}(0,0)P_{f(X)g(Y)}(1,1)>P_{f(X)g(Y)}(0,1)P_{f(X)g(Y)}(1,0),$$
the correlation between the same binary random variables, given $`\overline{Z}=\overline{z}`$, is negative (or “zero”) for all $`\overline{z}\overline{𝒵}`$, where $`\overline{Z}`$ is the random variable generated by sending $`Z`$ over Eve’s optimal channel $`P_{\overline{Z}|Z}`$.
A complete understanding of bound intrinsic information is of interest also because it automatically leads to a better understanding of bound entanglement in quantum information theory.
## 4 Concluding Remarks
We have considered the model of information-theoretic key agreement by public discussion from correlated information. More precisely, we have compared scenarios where the joint information is given by classical random variables and by quantum states (e.g., after execution of a quantum protocol). We proved a close connection between such classical and quantum information, namely between intrinsic information and entanglement. As an application, the derived parallels lead to an efficiently verifiable criterion for the fact that the intrinsic information vanishes. Previously, this quantity was considered to be quite hard to handle.
Furthermore, we have presented examples providing evidence for the fact that the close connections between classical and quantum information extend to the level of the protocols. As a consequence, the powerful tools and statements on the existence or rather non-existence of quantum-privacy-amplification protocols immediately carry over to the classical scenario, where it is often unclear how to show that no protocol exists. In particular, many examples (only some of which are presented above due to space limitations) coming from measuring bound entangled states, and for which none of the known classical secret-key agreement protocols is successful, strongly suggest that bound entanglement has a classical counterpart: intrinsic information which cannot be distilled to a secret key. This stands in sharp contrast to what was previously believed about classical key agreement. We state as an open problem to rigorously prove Conjectures 1 and 2.
Finally, we have proposed a measure for entanglement, based on classical information theory, with all the properties required for such a measure.
## Appendix A: Figures
## Appendix B: Measuring in “bad” Bases
In this appendix we show, by two examples, that the statements of Theorems 1 and 2 do not always hold for the standard bases and, in particular, not for arbitrary bases: Alice and Bob as well as Eve can perform bad measurements and give away an initial advantage. Let us begin with an example where measuring in the standard basis is a bad choice for Eve.
Example 6. Let us consider the quantum states
$$\mathrm{\Psi }=\frac{1}{\sqrt{5}}(|00+01+10|0+|00+11|1),\rho _{AB}=\frac{3}{5}P_{|00+01+10}+\frac{2}{5}P_{|00+11}.$$
If Alice, Bob, and Eve measure in the standard bases, we get the classical distribution (to be normalized)
| X | 0 | 1 |
| --- | --- | --- |
| Y (Z) | | |
| 0 | (0) 1 | (0) 1 |
| | (1) 1 | (1) 0 |
| 1 | (0) 1 | (0) 0 |
| | (1) 0 | (1) 1 |
For this distribution, $`I(X;YZ)>0`$ holds. Indeed, even $`S(X;Y||Z)>0`$ holds. This is not surprising since both $`X`$ and $`Y`$ are binary, and since the described parallels suggest that in this case, positive intrinsic information implies that a secret-key agreement protocol exists.
The proof of $`S(X;Y||Z)>0`$ in this situation is analogous to the proof of this fact in Example 3. The protocol consists of Alice and Bob independently making their bits symmetric. Then the repeat-code protocol can be applied.
However, the partial-transpose condition shows that $`\rho _{AB}`$ is separable. This means that measuring in the standard basis is bad for Eve. Indeed, let us rewrite $`\mathrm{\Psi }`$ and $`\rho _{AB}`$ as
$$\mathrm{\Psi }=\sqrt{\mathrm{\Lambda }}|\stackrel{}{m},\stackrel{}{m}|\stackrel{~}{0}+\sqrt{1\mathrm{\Lambda }}|\stackrel{}{m},\stackrel{}{m}|\stackrel{~}{1},\rho _{AB}=\frac{5+\sqrt{5}}{10}P_{|\stackrel{}{m},\stackrel{}{m}}+\frac{5\sqrt{5}}{10}P_{|\stackrel{}{m},\stackrel{}{m}},$$
where $`\mathrm{\Lambda }=(5+\sqrt{5})/10`$, $`|\stackrel{}{m},\stackrel{}{m}=|\stackrel{}{m}|\stackrel{}{m}`$, $`|\pm \stackrel{}{m}=\sqrt{(1\pm \eta )/2}|0\pm \sqrt{(1\eta )/2}|1`$, and $`\eta =1/\sqrt{5}`$.
In this representation, $`\rho _{AB}`$ is obviously separable. It also means that Eve’s optimal measurement basis is
$$|\stackrel{~}{0}=\sqrt{\mathrm{\Lambda }}|0\frac{1}{\sqrt{5\mathrm{\Lambda }}}|1,|\stackrel{~}{1}=\sqrt{1\mathrm{\Lambda }}|0\frac{1}{\sqrt{5(1\mathrm{\Lambda })}}|1.$$
Then, $`I(X;YZ)=0`$ holds for the resulting classical distribution. $`\mathrm{}`$
Not surprisingly, there also exist examples of distributions for which measuring in the standard bases is bad for Alice and Bob. These states are entangled, but $`I(X;YZ)=0`$ holds.
Example 7. Let the following classical distribution be given:
| X | 0 | 1 |
| --- | --- | --- |
| Y (Z) | | |
| 0 | (0) 0.0082 | (0) 0.0219 |
| | (1) 0.0006 | (1) 0.0202 |
| 1 | (0) 0.0729 | (0) 0.3928 |
| | (1) 0.0905 | (1) 0.3889204545 |
Because of
$$(0.0082+0.0006)(0.03928+0.3889204545)=(0.0219+0.0202)(0.0729+0.0905)$$
we have $`I(X;Y)=0`$, thus $`I(X;YZ)=0`$. On the other hand, the corresponding quantum state, for which the above distribution results by measuring in the standard bases, can be shown to be entangled. $`\mathrm{}`$
## Appendix C: A Protocol for Advantage Distillation
The following protocol for classical advantage distillation is called repeat-code protocol and was first proposed in . Note that there exist more efficient protocols in terms of the amount of extractable secret key. However, since we only want to prove qualitative possibility results, it is sufficient to look at this simpler protocol. Assume the scenario of Example 1.
Let $`N>0`$ be an even integer, and let Alice choose a random bit $`C`$ and send the block
$$X^NC^N:=[X_1C,X_2C,\mathrm{},X_NC]$$
over the classical channel. Here, $`X^N`$ stands for the block $`[X_1,X_2,\mathrm{},X_N]`$ of $`N`$ consecutive realizations of the random variable $`X`$, whereas $`C^N`$ stands for the $`N`$-bit block $`[C,C,\mathrm{},C]`$. Bob then computes $`[(CX_1)Y_1,\mathrm{},(CX_N)Y_N]`$ and (publicly) accepts exactly if this block is equal to either $`[0,0,\mathrm{},0]`$ or $`[1,1,\mathrm{},1]`$. In other words, Alice and Bob make use of a repeat code of length $`N`$ with the only two codewords $`[0,0,\mathrm{},0]`$ and $`[1,1,\mathrm{},1]`$.
Bob’s conditional error probability $`\beta _N`$ when guessing the bit sent by Alice, given that he accepts, is
$$\beta _N=\frac{1}{p_{a,N}}D^N\left(\frac{D}{1D}\right)^N,$$
where $`p_{a,N}=D^N+(1D)^N`$ is the probability that Bob accepts the received block. It is obvious that Eve’s optimal strategy for guessing $`C`$ is to compute the block $`[(CX_1)Z_1,\mathrm{},(CX_N)Z_N]`$ and guess $`C`$ as $`0`$ if at least half of the bits in this block are $`0`$, and as $`1`$ otherwise. Given that Bob correctly accepts, Eve’s error probability when guessing the bit $`C`$ with the optimal strategy is lower bounded by $`1/2`$ times the probability that she decodes to a block with $`N/2`$ bits $`0`$ and the same number of $`1`$’s. Hence we get that
$$\gamma _N\frac{1}{2}\left(\genfrac{}{}{0pt}{}{N}{N/2}\right)(1\delta )^{N/2}\delta ^{N/2}\frac{K}{\sqrt{N}}\left(2\sqrt{(1\delta )\delta }\right)^N$$
holds for some constant $`K`$ and for sufficiently large $`N`$ by using Stirling’s formula. Note that Eve’s error probability given that Bob accepts is asymptotically equal to her error probability given that Bob correctly accepts because Bob accepts erroneously only with asymptotically vanishing probability, given that he accepts.
Although it is not the adversary’s ultimate goal to guess the bits $`C`$ sent by Alice, it has been shown that the fact that $`\beta _N`$ decreases exponentially faster than $`\gamma _N`$ implies that for sufficiently large $`N`$, Bob has more (Shannon) information about the bit $`C`$ than Eve (see for example ). Hence Alice and Bob have managed to generate new random variables with the property that Bob obtains more information about Alice’s random bit than Eve has. Thus $`S(X;Y||Z)>0`$ holds.
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# PAIR PLASMA DOMINANCE IN THE PARSEC-SCALE RELATIVISTIC JET OF 3C345
## 1 Introduction
The study of extragalactic jets on parsec scales is astrophysically interesting in the context of the activities of the central engines of AGN. In particlar, a determination of their matter content would be an important step in the study of jet formation, propagation and emission. The two main candidates are a ‘normal plasma’ consisting of protons and relativistic electrons (for numerical simulations of shock fronts in a VLBI jet, see G$`\stackrel{´}{\mathrm{o}}`$mez et al. 1993, 1994a,b), and a ‘pair plasma’ consisting only of relativistic electrons and positrons (for theoretical studies of two-fluid concept, see Sol, Pelletier & Ass$`\stackrel{´}{\mathrm{e}}`$o 1989; Despringre & Fraix-Burnet 1997). Distinguishing between these possibilities is crucial for understanding the physical processes occurring close to the central ‘engine’ (presumably a supermassive black hole) in the nucleus.
VLBI is uniquely suited to the study of the matter content of pc-scale jets, because other observational techniques cannot image at milliarcsecond resolution and must resort to indirect means of studying the active nucleus. Recently, Reynolds et al. (1996) analyzed historical VLBI data of M87 jet at 5 GHz (Pauliny-Toth et al. 1981) and concluded that the core is probably dominated by an $`e^\pm `$ plasma. In the analysis, they utilized the standard theory of synchrotron self-absorption to constrain the magnetic field, $`B`$ \[G\], and the proper number density of electrons, $`N_\mathrm{e}^{}`$ \[1/cm<sup>3</sup>\] of the jet and derived the following condition for the core to be optically thick for self-absorption: $`N_\mathrm{e}^{}B^2>2\delta _{\mathrm{max}}^2`$, where $`\delta _{\mathrm{max}}`$ refers to the upper limit of the Doppler factor of the bulk motion. This condition is, however, applicable only for the VLBI observations of M87 core at epochs September 1972 and March 1973. Therefore, in order to apply the analogous method to other AGN jets or to M87 at other epochs, we must derive a more general condition.
On these grounds, Hirotani et al. (1999) generalized the condition $`N_\mathrm{e}^{}B^2>2\delta _{\mathrm{max}}^2`$ and applied it to the 3C 279 jet on parsec scales. In that paper, they revealed that core and components C3 and C4, of which spectra are obtained from the literature, are dominated by a pair plasma. It is interesting to note that the same concusion that 3C 279 jet is dominated by a pair plasma is derived from an independent method by Wardle et al. (1998), who studied the circularly polarized radio emission from 3C 279 jet.
In the present paper, we apply the same method to the 3C 345 jet. The quasar 3C345 (redshift z=0.594) is one of a class of core-dominated flat-spectrum radio sources that are believed to emit X-rays via the synchrotron self-Compton (SSC) process. VLBI imaging observations of the “superluminal” quasar 3C345 have been made at 5 GHz every year since 1977 (Unwin & Wehrle 1992) while 10.5 and 22 GHz observations have occurred at more frequent intervals (e.g., Biretta et al. 1986). The apparent speeds of components C2, C3, C4, and C5 increase monotonically with time from $`3c`$ to $`10c`$, consistent with a jet of constant Lorentz factor ($`\mathrm{\Gamma }=10`$) bending away from the line of sight (Zensus, Cohen, & Unwin 1995). Later, Unwin et al. (1997) studied the time evolution of spectral shapes and angular sizes of component C7 at a distance $`0.5`$ mas (2 pc) from the nucleus. Using the physical parameters given in the literature above, and deducing the kinetic luminosity from its core-position offset, we conclude that all the five jet components are likely dominated by an $`e^\pm `$ plasma. In the final section, we discuss the validity of assumptions.
We use a Hubble constant $`H_0=65h`$ km/s/Mpc and $`q_0=0.5`$ throughout this paper. These give a luminosity distance to 3C 345 of $`D_\mathrm{L}=3.06h^1`$ Gpc. An angular size or separation of $`1`$ mas corresponds to $`5.83h^1`$ pc. A proper motion of $`1\text{mas yr}^1`$ translates into a speed of $`\beta _{\mathrm{app}}=30.3h^1`$. Spectral index $`\alpha `$ is defined such that $`S_\nu \nu ^\alpha `$.
## 2 Constraints on Magnetic Flux and Particle Number Densities
We shall distinguish whether a radio-emitting component is dominated by a normal plasma or an $`e^\pm `$ plasma, by imposing two independent constraints on $`N_\mathrm{e}^{}`$. First, in § 2.1, we give the synchrotron self-absorption constraint, which is obtained by extending the work by Reynolds et al. (1996). Secondly in § 2.2, the kinematic luminosity constraint is presented.
### 2.1 Synchrotron Self-absorption Constraint
In this paper, we model a jet component with redshift $`z`$ as homogeneous spheres of angular diameter $`\theta _\mathrm{d}`$, containing a tangled magnetic field $`B`$ \[G\] and relativistic electrons which give a synchrotron spectrum with optically thin index $`\alpha `$ and maximum flux density $`S_\mathrm{m}`$ \[Jy\] at frequency $`\nu _\mathrm{m}`$. We can then compute the magnetic field density as follows (Cohen 1985; Ghisellini et al. 1992):
$$B=10^5b(\alpha )S_\mathrm{m}^2\left(\frac{\nu _\mathrm{m}}{\mathrm{GHz}}\right)^5\left(\frac{\theta _\mathrm{d}}{\mathrm{mas}}\right)^4\frac{\delta }{1+z},$$
(1)
where $`\delta `$ is the beaming factor defined by
$$\delta \frac{1}{\mathrm{\Gamma }(1\beta \mathrm{cos}\phi )},$$
(2)
$`\mathrm{\Gamma }1/\sqrt{1\beta ^2}`$ is the bulk Lorentz factor of the jet component moving with velocity $`\beta c`$, and $`\phi `$ is the orientation of the jet axis to the line of sight. The coefficient $`b(\alpha )`$ is given in Cohen (1985). Both $`\mathrm{\Gamma }`$ and $`\phi `$ can be uniquely computed from $`\delta `$ and $`\beta _{\mathrm{app}}`$ as follows:
$$\mathrm{\Gamma }=\frac{\beta _{\mathrm{app}}^2+\delta ^2+1}{2\delta },$$
(3)
$$\phi =\mathrm{tan}^1\left(\frac{2\beta _{\mathrm{app}}}{\beta _{\mathrm{app}}^2+\delta ^21}\right).$$
(4)
We assume that the electron number density between energies $`E`$ and $`E+dE`$ is expressed by a power law as
$$\frac{dN_\mathrm{e}^{}}{dE}=N_0E^{2\alpha 1}.$$
(5)
Integrating $`dN_\mathrm{e}^{}/dE`$ from $`\gamma _{\mathrm{min}}m_\mathrm{e}c^2`$ to $`\gamma _{\mathrm{max}}m_\mathrm{e}c^2`$, and assuming $`\gamma _{\mathrm{max}}\gamma _{\mathrm{min}}`$ and $`\alpha <0`$, we obtain the electron number density
$$N_\mathrm{e}^{}=\frac{\gamma _{\mathrm{min}}^{2\alpha }}{2\alpha }(m_\mathrm{e}c^2)^{2\alpha }N_0.$$
(6)
Computing the optical depth along the line of sight, Marscher (1983) expressed $`N_0`$ in terms of $`\theta _\mathrm{d}`$, $`S_\mathrm{m}`$, $`\nu _\mathrm{m}`$, and $`\alpha `$. Combining the result with equation (6), we finally obtain (see also Appendix B)
$`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ $`=`$ $`e(\alpha ){\displaystyle \frac{\gamma _{\mathrm{min}}^{2\alpha }}{2\alpha }}{\displaystyle \frac{h(1+z)^2q_0{}_{}{}^{2}\mathrm{sin}\phi }{zq_0+(q_01)(1+\sqrt{2q_0z+1})}}`$ (7)
$`\times `$ $`\left({\displaystyle \frac{\theta _\mathrm{d}}{\mathrm{mas}}}\right)^{4\alpha 7}\left({\displaystyle \frac{\nu _\mathrm{m}}{\mathrm{GHz}}}\right)^{4\alpha 5}S_\mathrm{m}{}_{}{}^{2\alpha +3}\left({\displaystyle \frac{\delta }{1+z}}\right)_{}^{2\alpha 3},`$
where $`e(\alpha )2.39\times 10^{16.77\alpha }`$ ($`0<\alpha <1.25`$). If the component is not resolved enough, this equation gives the lower bound of $`N_\mathrm{e}^{}`$.
### 2.2 Kinetic luminosity constraint
As described in Appendix B in detail, we can infer the kinetic luminosity, $`L_{\mathrm{kin}}`$, from the core-position offset, $`\mathrm{\Omega }_{r\nu }`$, due to synchrotron self-absorption. For the core, we assume a conical geometry with a small half opening angle $`\chi `$. Then $`L_{\mathrm{kin}}`$ measured in the rest frame of the AGN becomes
$`L_{\mathrm{kin}}`$ $``$ $`C_{\mathrm{kin}}K{\displaystyle \frac{r_1^2}{r_{}^3}}\beta \mathrm{\Gamma }(\mathrm{\Gamma }1)\chi ^2\left({\displaystyle \frac{\mathrm{\Omega }_{r\nu }/\nu _{}}{r_1\mathrm{sin}\phi }}\right)^{2(52\alpha )/(72\alpha )}`$ (8)
$`\times `$ $`\left[\pi C(\alpha ){\displaystyle \frac{\chi }{\mathrm{sin}\phi }}{\displaystyle \frac{K}{\gamma _{\mathrm{min}}}}{\displaystyle \frac{r_1}{r_{}}}{\displaystyle \frac{2\alpha }{\gamma _{\mathrm{min}}^{2\alpha }}}\left({\displaystyle \frac{\delta }{1+z}}\right)^{3/2\alpha }\right]^{4/(72\alpha )},`$
where $`K`$ is defined by equation (B13) and becomes $`0.1`$ for $`\alpha =0.5`$ if an energy equipartition holds between the radiating particles and the magnetic field.
For a pure pair plasma, we obtain $`C_{\mathrm{kin}}=\pi ^2\gamma _{}m_\mathrm{e}c^3/\gamma _{\mathrm{min}}`$, where $`\gamma _{}`$ is the averaged Lorentz factor of randomly moving electrons and positrons, which could be computed from equation (5) for a power-law distribution of radiating particles.
For a normal plasma, on the other hand, we obtain $`C_{\mathrm{kin}}=\pi ^2m_\mathrm{p}c^3/(2\gamma _{\mathrm{min}})`$, where $`m_\mathrm{p}`$ refers to the rest mass of a proton. It should be noted that $`\gamma _{\mathrm{min}}`$ takes a different value from a pair plasma.
Once $`L_{\mathrm{kin}}`$ of a stationary jet is obtained, we can deduce $`N_\mathrm{e}^{}`$ at an arbitrary position along the jet, even if the geometry deviates from a cone. When the jet has a perpendicular half width $`R_{}`$ at a certain position, $`L_{\mathrm{kin}}`$ and $`N_\mathrm{e}^{}`$ are related by
$$L_{\mathrm{kin}}=\pi R_{}{}_{}{}^{2}\beta c\mathrm{\Gamma }N_\mathrm{e}^{}(\mathrm{\Gamma }1)\left(\gamma _{}m_\mathrm{e}c^2+\gamma _+m_+c^2\right),$$
(9)
where $`\gamma _{}`$ and $`\gamma _+`$ refer to the averaged Lorentz factors of electrons and positively charged particles, respectively; $`m_+`$ designates the mass of the positive charge. Replacing the angular diameter distance, $`2R_{}/\theta _\mathrm{d}`$, with the luminosity distance divided by $`(1+z)^2`$, we can solve equation (9) for $`N_\mathrm{e}^{}`$ to obtain
$`N_\mathrm{e}^{}^{(\mathrm{kin})}`$ $`=`$ $`{\displaystyle \frac{3.42\times 10^2h^2q_0{}_{}{}^{4}(1+z)_{}^{4}}{\left[zq_0+(q_01)(1+\sqrt{2q_0z+1})\right]^2}}`$ (10)
$`\times `$ $`\left({\displaystyle \frac{\theta _\mathrm{d}}{\mathrm{mas}}}\right)^2{\displaystyle \frac{1}{\beta \mathrm{\Gamma }(\mathrm{\Gamma }1)}}{\displaystyle \frac{L_{46.5}}{\gamma _{}+\gamma _+m_+/m_\mathrm{e}}}\text{ cm}^3,`$
where $`L_{46.5}`$ refers to the kinetic luminosity in the unit of $`10^{46.5}\text{ergs s}^1`$. It should be noted that $`\gamma _{}+\gamma _+m_+/m_\mathrm{e}`$ becomes roughly $`2\gamma _{\mathrm{min}}\mathrm{ln}(\gamma _{\mathrm{min}}/\gamma _{\mathrm{max}})`$ for a pair plasma with $`\alpha 0.5`$, while it becomes $`1836`$ for a normal plasma. As a result, $`N_\mathrm{e}^{}^{(\mathrm{kin})}`$ for a pair plasma becomes about $`100\gamma _{\mathrm{min}}^1`$ times greater than that for a normal plasma. Since $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ is proportional to $`\gamma _{\mathrm{min}}^{2\alpha }`$, the ratio $`N_\mathrm{e}^{}{}_{}{}^{(\mathrm{kin})}/N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ for a pair plasma becomes about $`100\gamma _{\mathrm{min}}^{12\alpha }`$ times greater than that for a normal plasma. For a jet component close to the VLBI core, we may put $`\alpha 0.5`$; therefore, the dependence on $`\gamma _{\mathrm{min}}`$ virtually vanishes.
In short, we can exclude the possibility of a normal plasma dominance if $`1<N_\mathrm{e}^{}{}_{}{}^{(\mathrm{pair})}/N_\mathrm{e}^{}{}_{}{}^{(\mathrm{SSA})}100`$ is satisfied, where $`N_\mathrm{e}^{}^{(\mathrm{pair})}`$ refers to the value of $`N_\mathrm{e}^{}^{(\mathrm{kin})}`$ computed for a pair plasma. On the other hands, $`N_\mathrm{e}^{}{}_{}{}^{(\mathrm{pair})}/N_\mathrm{e}^{}{}_{}{}^{(\mathrm{SSA})}<1`$ implies that $`L_{\mathrm{kin}}`$ is underestimated. The conclusion is invulnerable against the value of $`\gamma _{\mathrm{min}}`$ of electrons and positrons.
## 3 Application to the 3C345 Jet
Let us apply the method described above to the 3C345 jet on parsec scales and investigate the matter content. It is, however, difficult to define $`\alpha `$, $`\nu _\mathrm{m}`$, and $`S_\mathrm{m}`$ of each component well, because the spectral information for an individual component is limited by the frequency coverage and quality of VLBI measurements near a given epoch. Therefore, Zensus et al. (1996) chose self-consistent values that matched the data and gave a reasonable fit to the overall spectrum when the components C2, C3, and C4 (hereafter, C2-C4), and the core are considered together (Table 1). For C2 and C3 they used the highest value for $`\nu _\mathrm{m}`$, while for C4 they used a representative possibility. Subsequently, Unwin et al. (1997) obtained these radio parameters for C5 and C7 by analogous method. We present these parameters together with their errors in Table 2. The jet half opening angle $`\chi 2.4^{}`$ is calculated from measuring the jet size within $`1`$ mas distance from the core (§ 4.3 in Lobanov 1998). We choose $`\alpha =0.65`$ as the spectral index of the core below the turnover frequency at $`700`$ GHz (§5.2 of Zensus et al. 1995).
### 3.1 Kinetic Luminosity
To estimate the kinetic luminosity from equation (8), we have to input $`\mathrm{\Gamma }`$, $`\phi `$, $`\mathrm{\Omega }_{r\nu }`$, and $`\delta `$ for a given $`C_{\mathrm{kin}}`$, $`K`$, $`\chi `$, and $`\alpha `$. Let us first consider $`\mathrm{\Gamma }`$, $`\phi `$, and $`\delta `$. As demonstrated in figure 4 in Unwin et al. (1995), a component (C7) accelerated as it moved away from the core: the Lorentz factor increased from $`\mathrm{\Gamma }5`$ to $`\mathrm{\Gamma }>10`$, and the viewing angle increased from $`\phi 2^{}`$ to $`\phi 10^{}`$. It is inappropriate to consider the case $`\phi \chi `$; therefore, we assume $`\phi 2^{}`$ for the core. In this case, $`\delta 1`$ holds to give $`L_{\mathrm{kin}}\mathrm{\Gamma }(\mathrm{\Gamma }1)/\delta \delta `$. In the case of a newly born component (C7) at 1992.05, Unwin et al. (1995) derived a conservative limit $`\delta >11.7`$, by assuming that C7 was the origin of the observed X-rays. Therefore, it is likely that $`\delta `$ is much greater than $`10`$ for the core, because $`\delta `$ decreased as the component moved away.
The core-position offset of the 3C 345 jet was reported by Lobanov (1998), who derived the reference value $`\mathrm{\Omega }_{r\nu }=10.7\text{pc}\text{Hz}`$. For a pair plasma with $`\alpha 0.5`$, $`\gamma _{}\gamma _{\mathrm{min}}\mathrm{ln}(\gamma _{\mathrm{max}}/\gamma _{\mathrm{min}})`$ holds in the expression of $`C_{\mathrm{kin}}`$; therefore, equation (8) gives
$$L_{\mathrm{kin}}10^{46}\frac{\mathrm{ln}(\gamma _{\mathrm{max}}/\gamma _{\mathrm{min}})}{10}K^{0.5}\left(\frac{\delta }{20}\right)\text{ergs s}^1.$$
(11)
On the other hand, for a normal plasma, equation (8) gives
$$L_{\mathrm{kin}}10^{46}\left(\frac{\gamma _{\mathrm{min}}}{100}\right)^1K^{0.5}\left(\frac{\delta }{20}\right)\text{ergs s}^1.$$
(12)
Unless the particles significantly dominates the magnetic field, $`K^{0.5}`$ does not exceed unity (see eqs. \[B14\] and \[B15\], which hold when an energy equipartition is realized between the radiating particles and the magnetic field). For a normal plasma jet, the energy distribution must cut off at $`\gamma _{\mathrm{min}}100`$ (§ 4; see also Wardle et al. 1998). Since $`\delta >100`$ is unlikely for the 3C 345 jet, we adopt $`L_{\mathrm{kin}}=10^{46.5}\text{ergs s}^1`$ (or equivalently $`L_{46.5}=1`$) as the representative upper bound in this paper. If $`L_{\mathrm{kin}}`$ becomes less than this value, the possibility of normal plasma dominance further decreases.
### 3.2 Equipartition Doppler factor
We estimate the value of $`\delta `$ by assuming an energy equipartition between the magnetic field and the radiating particles. In this case, $`K`$ becomes of the order of unity and $`\delta `$ is given by the so-called “equipartition Doppler factor” (Readhead 1994),
$`\delta `$ $`=`$ $`\delta _{\mathrm{eq}}`$ (13)
$``$ $`\left\{\left[{\displaystyle \frac{10^3F(\alpha )}{(\theta _\mathrm{d}/\mathrm{mas})}}\right]^{34}\left[{\displaystyle \frac{2(h/1.54)}{11/\sqrt{1+z}}}\right]^2(1+z)^{152\alpha }S_\mathrm{m}^{16}\left({\displaystyle \frac{\nu _\mathrm{m}}{\mathrm{MHz}}}\right)^{352\alpha }\right\}^{1/(132\alpha )},`$
where $`F(\alpha )`$ is given in Scott and Readhead (1977).
There is much justice in adopting the equipartition Doppler factor as the representative value. First, as G$`\ddot{\mathrm{u}}`$ijosa & Daly (1996) pointed out, $`\delta _{\mathrm{eq}}`$’s of various AGN jets have a high correlation with $`\delta _{\mathrm{min}}`$, the minimum allowed Doppler factor derived by comparing the predicted and the observed self-Compton flux (Marscher 1983, 1987; Ghisellini et al. 1993). (If a homogeneous moving sphere emits all the observed X-ray flux via synchrotron self-Compton process, then $`\delta `$ equals $`\delta _{\mathrm{min}}`$.) Secondly, the ratio $`\delta _{\mathrm{eq}}/\delta `$ depends weakly on the ratio $`u_\mathrm{p}/u_\mathrm{B}`$, where $`u_\mathrm{p}`$ and $`u_\mathrm{B}`$ refer to the energy densities of radiating particles (i.e., electrons and positrons) and the magnetic field, respectively. For $`\alpha =0.75`$ for instance, we obtain $`\delta _{\mathrm{eq}}/\delta =(u_\mathrm{p}/u_\mathrm{B})^{2/17}`$ (Readhead 1994). It is noteworthy that $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ depends relatively weakly on $`\theta _\mathrm{d}`$, $`\nu _\mathrm{m}`$, and $`\alpha `$, if we adopt $`\delta =\delta _{\mathrm{eq}}`$. For example, we obtain $`N_\mathrm{e}^{}{}_{}{}^{(\mathrm{SSA})}\theta _\mathrm{d}{}_{}{}^{2.9}\nu _{\mathrm{m}}^{}{}_{}{}^{5.5}S_{\mathrm{m}}^{}^{1.5}`$ for $`\alpha =0.75`$. This forms a striking contrast with $`N_\mathrm{e}^{}{}_{}{}^{(\mathrm{SSA})}\theta _\mathrm{d}{}_{}{}^{10}\nu _{\mathrm{m}}^{}{}_{}{}^{8}S_{\mathrm{m}}^{}{}_{}{}^{4.5}\delta _{}^{4.5}`$, which would be obtained from equation (7) without making any assumptions on $`\delta `$.
We present such representative values of $`\delta _{\mathrm{eq}}`$, $`\mathrm{\Gamma }_{\mathrm{eq}}(\beta _{\mathrm{app}}^2+\delta _{\mathrm{eq}}^2+1)/(2\delta _{\mathrm{eq}})`$, $`B`$, and $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ for C2-C4 in Table 1, and those for C5 and C7 in Table 2.
We first compare the values of $`\delta _{\mathrm{eq}}`$ with $`\delta _{\mathrm{min}}`$. It follows from Tables 1 and 2 that $`\delta _{\mathrm{eq}}>\delta _{\mathrm{min}}`$ is satisfied for all the eight cases, as expected. Moreover, the values of $`\delta _{\mathrm{eq}}`$ for C2-C4 at 1982.0 and those for C7 at the four epochs, decrease with increasing projected distance, $`\rho `$ \[mas\], from the core. As a result, the viewing angle computed from $`\beta _{\mathrm{app}}`$ and $`\delta _{\mathrm{eq}}`$ (see eq. ), $`\phi _{\mathrm{eq}}`$, increases with increasing $`\rho `$. (We exclude C5, for which the trajectory appears in a different position angle from those for C2-C4.) The results are qualitatively consistent with Zensus et al. (1995) and Unwin et al. (1997).
Let us next consider $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$. This variable is roughly constant at $`0.2\mathrm{cm}^3`$ for C2-C4, whereas it increases from $`0.5\text{cm}^3`$ at 1992.05 to $`10\text{cm}^3`$ at 1993.55 for C7. We consider that this tendency comes from insufficient angular resolution in particular when a component is close to the core. We can alternatively compute $`N_\mathrm{e}^{}`$ from $`N_\mathrm{e}^{}=(K/\gamma _{\mathrm{min}}m_\mathrm{e}c^2)(B^2/8\pi )`$, the energy equipartition. Reminding $`K0.1`$ for $`\alpha 0.5`$, we find that $`N_\mathrm{e}^{}`$ computed in this way is consistent with $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$.
We can compute $`N_\mathrm{e}^{}^{(\mathrm{pair})}`$, $`N_\mathrm{e}^{}^{(\mathrm{kin})}`$ for a pair plasma, from equation (8). The results of $`N_\mathrm{e}^{}^{(\mathrm{pair})}`$ are presented in Tables 1 and 2, together with the ratio $`N_\mathrm{e}^{}{}_{}{}^{(\mathrm{pair})}/N_\mathrm{e}^{}^{(\mathrm{SSA})}`$. It follows from Table 1 that C2 and C3 are likely dominated by a pair plasma. It is also suggested that C4 is dominated by pair plasma unless $`L_{\mathrm{kin}}`$ exceeds $`10^{46.5}`$ ergs/s. Unfortunately, the errors in $`B`$, $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ and $`N_\mathrm{e}^{}^{(\mathrm{pair})}`$ cannot be calculated, because those in $`\nu _\mathrm{m}`$ and $`S_\mathrm{m}`$ are not presented in Zensus et al. (1995). Furthermore, Table 2 indicates that C5 and C7 at all the four epochs are likely dominated by a pair plasma. Unfortunately, the meaningful errors in $`B`$, $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$, and $`N_\mathrm{e}^{}^{(\mathrm{pair})}`$ for C5 cannot be calculated, because its error in $`\theta _\mathrm{d}`$ (or $`\xi `$ in their notation) is not presented in Unwin et al. (1997). Nevertheless, the results of $`N_\mathrm{e}^{}{}_{}{}^{(\mathrm{pair})}/N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ strongly suggest that the jet components of 3C 345 on parsec scales are dominated by a pair plasma.
## 4 Discussion
In summary, we derive the proper electron number density, $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$, of a homogeneous radio-emitting component of which spectral turnover is due to synchrotron self-absorption. Comparing $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ with the density derived from the kinetic luminosity of the jet, we can investigate whether we can exclude the possibility of normal plasma ($`e^{}`$-p) dominance. Applying this method to the “superluminal” quasar 3C345, using the published spectrum data of C2, C3, C4, C5, and C7, we find that all the five components are likely dominated by a pair plasma.
As demonstrated in the last part of §2, the conclusion is invulnerable against the undetermined value of $`\gamma _{\mathrm{min}}`$ of electrons and positrons. However, if $`\gamma _{\mathrm{min}}`$ for a normal plasma were to be significantly less than $`100`$, then the possibility of a normal plasma dominance could not be ruled out in general. In the case of the 3C 345 jet, equation (12) would give $`L_{\mathrm{kin}}10^{48}\text{ergs s}^1`$ for a normal plasma with $`\gamma _{\mathrm{min}}1`$. In this case, the large kinetic luminosity ($`10^{48}\text{ergs s}^1`$) is carried by protons, because
$$\gamma _{}m_\mathrm{e}c^2\frac{\gamma _{\mathrm{min}}}{K}m_\mathrm{e}c^2m_\mathrm{p}c^2$$
(14)
holds. Nevertheless, we consider that such a jet is unlikely, because the protons carry about two orders of magnitude more energy than is seen to be dissipated as synchrotron radiation ($`10^{46}\text{ergs s}^1`$). Electrons on parsec scales will not be cooled down so rapidly shortly after being heated-up at the shock fronts.
It is interesting to consider the case when $`\delta `$ is estimated by other methods than the energy equipartition. As an example, let us consider a jet motion with a roughly constant Lorentz factor; Zensus et al. (1995) derived that $`\mathrm{\Gamma }10`$ is close to the smallest value that is consistent with all their available kinematic constraints. Such values of $`\delta `$ and $`\phi `$ are denoted by the solid dots in Fig. 12 of their paper and tabulated again in table 3 in the present paper. Using those data, we can compute $`B`$ and $`N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ of each component (table 3). For C2, we adopt $`\mathrm{\Gamma }=13`$ rather than $`10`$, because $`\beta _{\mathrm{app}}=12.9`$ for $`h=1`$ (or equialently $`H_0=65`$) gives $`\mathrm{\Gamma }>\sqrt{1+\beta _{\mathrm{app}}^2}=12.9`$. The results of $`N_\mathrm{e}^{}{}_{}{}^{(\mathrm{pair})}/N_\mathrm{e}^{}^{(\mathrm{SSA})}`$ show again that C2-C4 at 1982.0 are likely dominated by a pair plasma.
## Appendix A Derivation of the Synchrotron Self-absorption Constraints
We assume that the parsec-scale jet close to the core propagates conically with a half opening angle $`\chi `$ in the observer’s frame. Then the optical depth $`\tau `$ for synchrotron self absorption is given by
$$\tau _\nu (R)=\frac{2R\mathrm{sin}\chi }{\mathrm{sin}(\phi +\chi )}\alpha _\nu ,$$
(A1)
where $`R`$ is the distance of the position from the injection point of the jet and $`\alpha _\nu `$\[1/cm\] refers to the absorption coefficient. For a small half opening angle ($`\chi 1`$), this equation can be approximated as
$$\tau _\nu (R)=2R\frac{\chi }{\mathrm{sin}\phi }\alpha _\nu $$
(A2)
Since $`\tau `$ and $`R\chi `$ are Lorentz invariants, we obtain
$$\frac{\alpha _\nu }{\mathrm{sin}\phi }=\frac{\alpha _\nu ^{}}{\mathrm{sin}\phi ^{}},$$
(A3)
where a quantity with an asterisk is measured in the co-moving frame, while that without an asterisk in the observer’s frame. Since $`\nu \alpha _\nu `$ is also Lorentz invariant, equation (A3) gives
$$\frac{\mathrm{sin}\phi ^{}}{\mathrm{sin}\phi }=\frac{\nu }{\nu ^{}}=\frac{\delta }{1+z}.$$
(A4)
Combining equations (A2) and (A4), we obtain
$$\tau _\nu =\frac{1+z}{\delta }\frac{2R\chi }{\mathrm{sin}\phi }\alpha _\nu ^{}=\frac{1+z}{\delta }\frac{1}{\mathrm{sin}\phi }\frac{\theta _\mathrm{d}D_\mathrm{L}}{(1+z)^2}\alpha _\nu ^{},$$
(A5)
where the angular diameter distance of the jet, $`2R\chi /\theta _\mathrm{d}`$, is rewritten with the luminosity distance, $`D_\mathrm{L}`$, divided by $`(1+z)^2`$; here, $`\theta _\mathrm{d}`$ is the angular diameter of the component in the perpendicular direction of the jet propagation. If we observe $`\tau _\nu `$ at the turnover frequency, $`\nu _\mathrm{m}`$, it becomes a function of the optical thin spectral index $`\alpha `$, which is tabulated in Scott and Readhead (1977).
Averaging over pitch angles of the isotropic electron power-law distribution (eq. ), we can write down the absorption coefficient in the co-moving frame as (Le Roux 1961, Ginzburg & Syrovatskii 1965)
$$\alpha _\nu ^{}=C(\alpha )r_{}{}_{}{}^{2}k_{\mathrm{e}}^{}\frac{\nu _{}}{\nu ^{}}\left(\frac{\nu _\mathrm{B}}{\nu ^{}}\right)^{(2\alpha +3)/2},$$
(A6)
where $`\nu _{}c/r_{}c/[e^2/(m_\mathrm{e}c^2)]`$ and $`\nu _\mathrm{B}eB/(2\pi m_\mathrm{e}c)`$. The coefficient $`C(\alpha )`$ is given in Table 1 of Gould (1979).
Substituting equation (A6) into (A5), and assuming $`\alpha <0`$ and $`\gamma _{\mathrm{min}}\gamma _{\mathrm{max}}`$, we obtain with the aid of (5)
$`N_\mathrm{e}^{}B^{\alpha +1.5}`$ $`=`$ $`{\displaystyle \frac{m_\mathrm{e}c}{e^2}}\left({\displaystyle \frac{e}{2\pi m_\mathrm{e}c}}\right)^{1.5+\alpha }{\displaystyle \frac{\tau _\nu (\alpha )}{C(\alpha )}}{\displaystyle \frac{\gamma _{\mathrm{min}}^{2\alpha }}{2\alpha }}`$ (A7)
$`\times {\displaystyle \frac{(1+z)^2}{D_\mathrm{L}}}{\displaystyle \frac{\mathrm{sin}\phi }{\theta _\mathrm{d}}}\left({\displaystyle \frac{1+z}{\delta }}\right)^{\alpha +1.5}\nu ^{\alpha +2.5}.`$
Evaluating $`\nu `$ at the turnover frequency, $`\nu =\nu _\mathrm{m}`$, and combining with equation (1), we obtain $`N_\mathrm{e}^{}`$ presented in equation (7), which equals $`(\gamma _{\mathrm{min}}m_\mathrm{e}c^2){}_{}{}^{2\alpha }/(2\alpha )`$ times $`N_0`$ given in equation (3) in Marscher (1983). It is noteworthy that electron number density in the observer’s frame can be obtained if we multiply $`(1+z)/\delta `$ on $`N_\mathrm{e}^{}`$.
## Appendix B Kinetic luminosity inferred from core-position offset
In this appendix, we deduce the kinetic luminosity of a jet from its core-position offset due to synchrotron self-absorption. This method was originally developed by Lobanov (1988). However, our results somewhat differs from his results; therefore, we explicitly describe the derivation so that the readers can check it.
### B.1 Scaling Law
First, we assume that $`N_\mathrm{e}^{}`$ and $`B`$ scale on $`r`$ in the following manner:
$`N_e^{}=N_1r^n,B=B_1r^m,`$ (B1)
where $`N_1`$ and $`B_1`$ refer to the values of $`N_\mathrm{e}^{}`$ and $`B`$ at $`r_1=1`$ pc, respectively; $`rR/r_1`$. Introducing dimensionless variables
$`x_\mathrm{N}r_1r_{}{}_{}{}^{2}N_{1}^{}`$
$`x_\mathrm{B}\nu _{\mathrm{B}_1}/\nu _{}={\displaystyle \frac{eB_1}{2\pi m_\mathrm{e}c}},`$ (B2)
and utilizing equation (A6), we obtain from the left equality in equation (A5)
$$\tau _\nu =C(\alpha )\frac{2\chi }{\mathrm{sin}\phi }\frac{2\alpha }{\gamma _{\mathrm{min}}^{2\alpha }}\left(\frac{1+z}{\delta }\right)^ϵ\left(\frac{\nu }{\nu _{}}\right)^{1ϵ}r^{1nmϵ}x_\mathrm{N}x_\mathrm{B}{}_{}{}^{ϵ},$$
(B3)
where $`ϵ3/2\alpha `$.
At a given frequency $`\nu `$, the flux density will peak at the position where $`\tau _\nu `$ becomes unity. Thus setting $`\tau =1`$ and solving equation (B3) for $`r`$, we obtain the distance from the VLBI core observed at frequency $`\nu `$ from the central engine as
$$r(\nu )=\left(x_\mathrm{B}{}_{}{}^{k_\mathrm{b}}F\frac{\nu _{}}{\nu }\right)^{1/k_\mathrm{r}}$$
(B4)
where
$$F(\alpha )\left[C(\alpha )\frac{2\chi }{\mathrm{sin}\phi }\frac{2\alpha }{\gamma _{\mathrm{min}}^{2\alpha }}\left(\frac{\delta }{1+z}\right)^ϵx_\mathrm{N}\right]^{1/(ϵ+1)}$$
(B5)
$$k_\mathrm{b}\frac{32\alpha }{52\alpha },$$
(B6)
$$k_\mathrm{r}\frac{(32\alpha )m+2n2}{52\alpha }.$$
(B7)
### B.2 Core-Position Offset
If we mesure $`r(\nu )`$ at two different frequencies (say $`\nu _\mathrm{a}`$ and $`\nu _\mathrm{b}`$), equation (B4) gives the dimensionless, projected distance of $`r(\nu _\mathrm{a})r(\nu _\mathrm{b})`$ as
$$\mathrm{\Delta }r_{\mathrm{proj}}=\left[r(\nu _\mathrm{a})r(\nu _\mathrm{b})\right]\mathrm{sin}\phi =(x_\mathrm{B}{}_{}{}^{k_\mathrm{b}}F\nu _{})^{1/k_\mathrm{r}}\frac{\nu _\mathrm{b}^{1/k_\mathrm{r}}\nu _\mathrm{a}^{1/k_\mathrm{r}}}{\nu _\mathrm{a}^{1/k_\mathrm{r}}\nu _\mathrm{b}^{1/k_\mathrm{r}}}\mathrm{sin}\phi .$$
(B8)
Defining the core-position offset as
$$\mathrm{\Omega }_{r\nu }r_1\mathrm{\Delta }r_{\mathrm{proj}}\frac{\nu _\mathrm{a}^{1/k_\mathrm{r}}\nu _\mathrm{b}^{1/k_\mathrm{r}}}{\nu _\mathrm{b}^{1/k_\mathrm{r}}\nu _\mathrm{a}^{1/k_\mathrm{r}}},$$
(B9)
we obtain
$$\frac{\mathrm{\Omega }_{r\nu }}{r_1}=(x_\mathrm{B}^{k_\mathrm{b}}F\nu _{})^{1/k_\mathrm{r}}\mathrm{sin}\phi $$
(B10)
To express $`x_\mathrm{B}`$ in terms of $`x_\mathrm{N}`$ and $`\mathrm{\Omega }_{r\nu }`$, we can invert equation (B10) as
$$x_\mathrm{B}=\left(\frac{\mathrm{\Omega }_{r\nu }}{r_1\mathrm{sin}\phi }\right)^{k_\mathrm{r}/k_\mathrm{b}}(F\nu _{})^{1/k_\mathrm{b}}.$$
(B11)
Note that $`x_\mathrm{N}`$ is included in $`F=F(\alpha )`$.
Setting $`\nu _\mathrm{b}\mathrm{}`$ in equation (B8), we obtain the absolute distance of the VLBI core measured at $`\nu `$ from the central engine as
$$r_{\mathrm{core}}(\nu )=\frac{\mathrm{\Omega }_{r\nu }}{r_1\mathrm{sin}\phi }\nu ^{1/k_\mathrm{r}}.$$
(B12)
That is, once $`\mathrm{\Omega }_{r\nu }`$ is obtained from multi-frequency VLBI observations, we can deduce the distance of the synchrotron-self-absorbing VLBI core from the central engine, assuming the scaling laws of $`N_\mathrm{e}^{}`$ and $`B`$ as equation (B1).
We next represent $`x_\mathrm{N}`$ and $`x_\mathrm{B}`$ (or equivalently, $`N_1`$ and $`B_1`$) as a function of $`\mathrm{\Omega }_{r\nu }`$. To this end, we relate $`N_\mathrm{e}^{}`$ and $`B`$ as follows:
$$N_\mathrm{e}^{}\gamma _{\mathrm{min}}m_\mathrm{e}c^2=K\frac{B^2}{8\pi }.$$
(B13)
When an energy equipartition between the radiating particles and the magnetic field holds, equation (5) gives for $`\alpha =0.5`$
$$K=\frac{1}{\mathrm{ln}(\gamma _{\mathrm{max}}/\gamma _{\mathrm{min}})}0.1,$$
(B14)
whereas for $`\alpha <0.5`$
$$K=\frac{2\alpha +1}{2\alpha }\frac{\gamma _{\mathrm{max}}{}_{}{}^{2\alpha }\gamma _{\mathrm{min}}^{2\alpha }}{\gamma _{\mathrm{max}}{}_{}{}^{2\alpha +1}\gamma _{\mathrm{min}}^{2\alpha +1}}.$$
(B15)
Substituting $`N_\mathrm{e}^{}=N_1r^2`$ and $`B=B_1r^1`$ into (B13), and replacing $`N_1`$ and $`B_1`$ with $`x_\mathrm{N}`$ and $`x_\mathrm{B}`$, we obtain
$$x_\mathrm{N}=\frac{\pi }{2}\frac{K}{\gamma _{\mathrm{min}}}\frac{r_1}{r_{}}x_\mathrm{B}^2$$
(B16)
It is noteworthy that the assumptions of $`n=2`$ and $`m=1`$, which results in $`k_\mathrm{r}=1`$, guarantees the energy equipartition at an arbitrary distance, $`r`$.
Combining equations (B11) and (B16), we obtain
$`x_\mathrm{B}`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{\Omega }_{r\nu }/\nu _{}}{r_1\mathrm{sin}\phi }}\right)^{(52\alpha )/(72\alpha )}`$ (B17)
$`\times \left[\pi C(\alpha ){\displaystyle \frac{\chi }{\mathrm{sin}\phi }}{\displaystyle \frac{K}{\gamma _{\mathrm{min}}}}{\displaystyle \frac{r_1}{r_{}}}{\displaystyle \frac{2\alpha }{\gamma _{\mathrm{min}}^{2\alpha }}}\left({\displaystyle \frac{\delta }{1+z}}\right)^ϵ\right]^{2/(72\alpha )}.`$
The particle number density, $`x_\mathrm{N}`$, can be readily computed from equation (B16).
### B.3 Kinetic luminosity
We can now relate the kinetic luminosity with the core-position offset. The factor $`N_\mathrm{e}{}_{}{}^{}R_{}^{2}`$ in equation (9) can be expressed in terms of $`x_\mathrm{N}`$ and hence $`x_\mathrm{B}`$ as
$`N_\mathrm{e}^{}R^2`$ $`=`$ $`N_1r_1{}_{}{}^{2}={\displaystyle \frac{r_1}{r_{}^2}}x_\mathrm{N}`$ (B18)
$`=`$ $`{\displaystyle \frac{\pi }{2}}{\displaystyle \frac{K}{\gamma _{\mathrm{min}}}}{\displaystyle \frac{r_1^2}{r_{}^3}}x_\mathrm{B}{}_{}{}^{2}.`$
For a pure pair plasma, we obtain $`\gamma _+=\gamma _{}`$ and $`m_+=m_\mathrm{e}`$. Therefore, for a conical geometry, we can put $`R_{}=R\chi `$ in equation (9) to obtain equation (8), where $`C_{\mathrm{kin}}=\pi ^2\gamma _{}m_\mathrm{e}c^3/\gamma _{\mathrm{min}}`$.
In the same manner, for a normal plasma, we have $`\gamma _+=1`$ and $`m_+=m_\mathrm{p}`$. In this case, we obtain $`C_{\mathrm{kin}}=\pi ^2m_\mathrm{p}c^3/(2\gamma _{\mathrm{min}})`$.
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# DO-TH 00/10 Single Pion Production in Neutrino Reactions and Estimates for Charge-Exchange Effects.
## I INTRODUCTION
There is strong evidence for the mixing of muon neutrinos with another state, being either tau or sterile neutrinos. The evidence comes from atmospheric neutrino experiments which observe a decrease of muon neutrinos in charged current reactions, but no decrease of the corresponding electron neutrino interactions .
In order to obtain better insight into the oscillation which takes place and in order to eliminate dependence on the flux there are proposals and experiments being planned and constructed, which look at the neutral current interactions. These are reactions which will use neutrinos of an average energy of 1 GeV producing the resonances between 1.0 and 1.6 GeV/c<sup>2</sup>. One proposal considers the production of pions directly by the atmospheric neutrinos and the detection of $`\pi ^0`$’s with the help of two ring events. This method is limited by the low flux of atmospheric neutrinos.
More powerful are experiments which use neutrinos from an accelerator with two detectors; the first one nearby the accelerator and a second further away. The nearby detector will be able to detect all pions and check the flux, as well as the cross-sections for these reactions. The detector with the long-baseline (300-400 km) will observe the charged and neutral current reactions.
We classify the reactions. Quasi-elastic scattering $`\nu _\mu +n\mu ^{}+p`$ is well understood. The oscillation of muon neutrinos into other neutrinos, tau or sterile, will produce a reduction of muon events in the far away detector.
In addition to the above reaction there are excitations of resonances and their subsequent decays
$$\nu _\mu +p\mu ^{}+p+\pi ^+$$
(1)
$$\nu _\mu +n\mu ^{}+n+\pi ^+$$
(2)
$$\nu _\mu +n\mu ^{}+p+\pi ^0.$$
(3)
Furthermore, there are the neutral current reactions
$$\nu _\mu +p\nu _\mu +p+\pi ^0$$
(4)
$$\nu _\mu +p\nu _\mu +n+\pi ^+$$
(5)
$$\nu _\mu +n\nu _\mu +n+\pi ^0$$
(6)
$$\nu _\mu +n\nu _\mu +p+\pi ^{}.$$
(7)
The theory for the production of these states is known for thirty years now and there are several calculations available. The charged current reactions have been studied extensively and the production of the $`\mathrm{\Delta }^{++}`$ has been understood theoretically. It has also been measured experimentally with good agreement between theory and experiments. For the other charged current reactions there are few experimental measurements. For this reason the nearby detectors of the experiments should study the reactions using light and heavy nuclei as targets.
Knowledge of the neutral current reactions is even more limited. Two of the latest calculations of charged and neutral current reactions for the production of single-pion cross sections differ by approximately $`20\%`$ . One should make all possible efforts now to reduce the overall uncertainty and measure the various channels experimentally.
A second difficulty arises from the fact that the experimental targets are heavy materials so that the interactions take place on protons and neutrons bound in nuclei like for example $`{}_{8}{}^{}O_{}^{16}`$, $`{}_{18}{}^{}Ar^{40}`$ or $`{}_{26}{}^{}Fe^{56}`$. In the heavy nuclei, the produced pions rescatter before they exit from the nucleus and are subject to two phenomena: (1) the cross-sections are reduced by the Pauli exclusion principle, when the energy of the recoiling nucleon is low and can not occupy a filled level of nucleons, and (2) the pion charge exchange due to rescattering. These phenomena are known and have been the subject of extensive studies .
The expectations of the experiments are the following. For all charge current reactions we expect a reduction of the observed rates in the far away detector because some of the muon neutrinos oscillated into another state. For the neutral current reactions we expect no reduction in rate if the oscillation is to tau neutrino because all neutrinos contribute equally to neutral current reactions. We expect a reduction if the oscillation is to sterile neutrinos.
Now, since the reduction is expected to be approximately $`40\%`$, it is important to understand all possible corrections. An important requirement is that the nearby and the far away detector use the same nuclei as targets. In case this is not possible, then corrections will have to be applied.
Because of the importance of the experiments and the opportunities they present for establishing the charged and neutral current reactions, we have undertaken the task of calculating the cross-section on free protons and neutrons. This way we produce differential, as well as integrated cross-sections. Then we use the obtained results to calculate the corrections which are present in the nuclei.
This paper is organized as follow: Section 2 is devoted to the calculation of the differential and total cross-sections for single-pion production in neutrino-nucleon interactions. In Section 3 we discuss the nuclear effects involved in this process and calculate the energy spectra for charged and neutral pions for a few different materials typically used as targets in the experiments. Finally, results and conclusions are presented in Section 4.
## II Single-Pion Production
### A General Formalism
In this section we present the main equations and the form factors used to evaluate the differential and total cross sections for single-pion production in neutrino-nucleon interactions. For neutrino energies of a few GeV the single-pion production proceeds mainly through the excitation of the lower resonances. The main contribution to the cross section comes from the production and the subsequent decay of the $`\mathrm{\Delta }`$(1232)P<sub>33</sub> resonance. Nevertheless, some of the channels receive a non negligible contribution from the isospin 1/2 resonances as, for example, the N(1440)P<sub>11</sub> and the N(1535)S<sub>11</sub> resonances.
The channels under investigation in this paper are the three charged current and the four neutral current channels listed in the introduction in Eq. 1.1 to 1.7.
Using Clebsh-Gordan coefficient, it is easy to verify that the amplitudes for these seven channels are given by the following equations:
$`A(\mu ^{}+p+\pi ^+)`$ $`=`$ $`A_3^{cc}`$ (8)
$`A(\mu ^{}+n+\pi ^+)`$ $`=`$ $`1/3A_3^{cc}+2\sqrt{2}/3A_1^{cc}`$ (10)
$`A(\mu ^{}+p+\pi ^0)`$ $`=`$ $`\sqrt{2}/3A_3^{cc}+2/3A_1^{cc}`$ (12)
$`A(\nu _\mu +p+\pi ^0)`$ $`=`$ $`\sqrt{2}/3A_3^{nc}+1/3A_1^{nc}+1/3A_1^0`$ (14)
$`A(\nu _\mu +n+\pi ^+)`$ $`=`$ $`1/3A_3^{nc}+\sqrt{2}/3A_1^{nc}+\sqrt{2}/3A_1^0`$ (16)
$`A(\nu _\mu +n+\pi ^0)`$ $`=`$ $`\sqrt{2}/3A_3^{nc}+1/3A_1^{nc}1/3A_1^0`$ (18)
$`A(\nu _\mu +p+\pi ^{})`$ $`=`$ $`1/3A_3^{nc}\sqrt{2}/3A_1^{nc}+\sqrt{2}/3A_1^0`$ (20)
where, in this paper, $`A_3^{cc,nc}`$ corresponds to the amplitude for the production of the P<sub>33</sub> resonance, $`A_1^{cc,nc}`$ is the sum of the amplitudes for the production of the P<sub>11</sub> and S<sub>11</sub> resonances and $`A_1^0`$ is the sum of the isoscalar contributions of the P<sub>11</sub> and S<sub>11</sub> resonances to the cross section.
Notice that, as suggested in Ref., the neutral current amplitudes $`A_3^{nc}`$, $`A_1^{nc}`$ and $`A_1^0`$ can be derived from the corresponding charged current amplitudes $`A_3^{cc}`$ and $`A_1^{cc}`$ by simply rescaling the vector and axial form factors. In the case of the $`A_3^{nc}`$ and $`A_1^{nc}`$ amplitudes the vector and axial charged current form factors need to be multiplied by $`12sin^2\theta _W`$ and by $`1`$ respectively, where $`\theta _W`$ is the weak mixing angle. For the $`A_1^0`$ instead the vector and axial charged current form factors need to be multiplied respectively by $`2/3sin^2\theta _W`$ and $`0`$. Furthermore, we want to point out that, since the $`A_1^0`$ amplitude turned out to be very small compared to $`A_3^{nc}`$ and $`A_1^{nc}`$, we neglected the isoscalar contribution in our evaluation of the cross sections.
#### 1 $`\mathrm{\Delta }`$(1232)P<sub>33</sub>
As it was mentioned in the introduction, the theory for the production of the $`\mathrm{\Delta }`$(1232)P<sub>33</sub> is well known and understood, and several independent calculations have already been published, showing good agreement with the experimental results . Therefore, rather than developing our own formalism for this process, we decided to follow the article of Schreiner and von Hippel in Ref..
The fully differential cross section $`\mathrm{d}\sigma /\mathrm{d}Q^2\mathrm{d}W\mathrm{d}\mathrm{\Omega }_\pi `$ from Ref. has been integrated over the polar angle $`\varphi _\pi `$ and converted into the triple-differential cross section $`\mathrm{d}\sigma /\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi `$ by using the fact that $`E_\pi ^{lab}=\gamma E_\pi ^{CM}+\beta \gamma |\stackrel{}{p}_\pi ^{CM}|cos\theta _\pi `$, where the superscripts $`lab`$ and $`CM`$ denote quantities measured in the laboratory and in the center of mass frame respectively. The total cross section has been obtained by integrating over the allowed range of values for $`Q^2`$ and $`E_\pi `$, with a cut on the invariant mass $`W`$ range at 1.6 GeV.
The axial and vector form factors used in this calculation are the ones given by Alvarez-Ruso et al. in Eq. 12, 13 and 18 of Ref.. Notice that, since these form factors have been derived from photo- and electro-production experiments in which a $`\mathrm{\Delta }^+`$ or a $`\mathrm{\Delta }^0`$ was produced, in order to obtain the correct cross section for the $`\mathrm{\Delta }^{++}`$ production, all the form factors need to be multiplied by $`\sqrt{3}`$ due to the fact that $`<\mathrm{\Delta }^{++}\mathrm{V}_\alpha \mathrm{p}>=\sqrt{3}<\mathrm{\Delta }^+\mathrm{V}_\alpha ^{\mathrm{em}}\mathrm{p}>`$.
#### 2 N(1440)P<sub>11</sub> and N(1535)S<sub>11</sub>
As shown in Ref., the triple-differential cross section $`\mathrm{d}\sigma /\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi `$ for the production of the P<sub>11</sub> and S<sub>11</sub> resonances is given by the following equation:
$$\frac{\mathrm{d}\sigma }{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }=\frac{W}{4M_N}\frac{G_F^2Q^2}{(2\pi )^4\nu ^2}\frac{1}{\sqrt{\nu ^2+Q^2}}\left(\left(1\frac{\nu }{E_\nu }\right)|M_s|^2+\frac{1}{2}\left(1\frac{\nu }{E_\nu }\right)^2|M_r|^2+\frac{1}{2}|M_l|^2\right)$$
(21)
where $`M_N`$ is the nucleon mass, $`G_F`$ is the Fermi constant and $`\nu `$ is the difference between the energies of the incoming and the outgoing lepton. Remembering that nearby the resonance only the s-channel is essential, the three matrix elements $`M_i`$ with $`i=r,l,s`$ can be defined as follow:
$$M_i=f_R\overline{u}(p)\gamma _5(\text{/}p+\text{/}q+M_R)(g_V\text{/}ϵ^ig_A\text{/}ϵ^i\gamma _5)u(p)f(W)$$
(22)
where $`f_R`$ is the coupling constant of the pion to the nucleon and the resonance, $`M_R`$ is the resonance mass, $`p`$ and $`p`$ are the four-momenta of the initial and final nucleon respectively, q is the four-momentum transferred from the leptons to the hadrons and $`g_{V,A}`$ are the vector and axial form factors. The values for $`f_R`$ are given in Eq. B7 and B10 in Appendix B of Ref..
The Breit-Wigner factor $`f(W)`$ in Eq.2.9 can be written in the following way:
$$f(W)=\frac{1}{(W^2M_R^2)+iM_R\mathrm{\Gamma }_R}$$
(23)
where $`\mathrm{\Gamma }_R(W)=\mathrm{\Gamma }_R^0q_\pi (W)/q_\pi (M_R)`$ with $`q_\pi (W)=\sqrt{(W^2M_N^2M_\pi ^2)^24M_N^2M_\pi ^2}/2W`$ and $`\mathrm{\Gamma }_R^0`$ the width of the resonance.
Finally, the polarization vectors $`ϵ^i`$ with $`i=r,l,s`$ used in Eq.2.9 are defined as:
$$ϵ^s=\frac{1}{Q}(\sqrt{\nu ^2+Q^2},0,0,\nu ),ϵ^{r,l}=\frac{1}{\sqrt{2}}(0,1,\pm i,0).$$
(24)
For the form factors we use the expressions obtained by Fogli and Nardulli in Ref..
As in the case of the P<sub>33</sub> resonance, the total cross section has been obtained by integrating over the allowed range of values for $`Q^2`$ and $`E_\pi `$, with a cut on the invariant mass $`W`$ range at 1.6 GeV.
### B Results
In this section we present our results for the total cross section for the seven channels under examination and, where possible, we compare these results with experimental data.
In Fig.1 the total cross section for the $`\nu _\mu +p\mu ^{}+p+\pi ^+`$ process has been plotted versus the incoming neutrino energy. The data points have been taken from Ref. (solid circles) and from Ref. (empty circles). As it can be seen, the agreement between the theoretical curve and the experimental data is quite good.
Fig.2 and Fig.3 display the total cross sections for the $`\nu _\mu +n\mu ^{}+p+\pi ^0`$ and $`\nu _\mu +n\mu ^{}+n+\pi ^+`$ processes respectively, again plotted versus the incoming neutrino energy. In this case the data points have been taken from Ref. (solid circles), from Ref. (empty circles) and from Ref. (crosses). Also in this case the agreement between the theoretical results and the data points is reasonably good.
The difference between the theoretical and the experimental results can be partially explained by taking into account the fact that, while the theoretical curves have been estimated imposing a cut on the invariant mass $`W`$ at 1.6 GeV, the experimental points have been obtained without any cut. Notice also that we didn’t include any non-resonant background in our evaluation of the cross sections.
In the case of the neutral current interactions, the experimental results are presented in the form of ratios between each of the neutral current channels and one of the charged current channels. For these reason, Fig.4, Fig.5, Fig.6 and Fig.7, which display respectively the total cross sections of the $`\nu _\mu +p\nu _\mu +p+\pi ^0`$, $`\nu _\mu +p\nu _\mu +n+\pi ^+`$, $`\nu _\mu +n\nu _\mu +n+\pi ^0`$ and $`\nu _\mu +n\nu _\mu +p+\pi ^{}`$ processes plotted versus the incoming energy, have no data points. Nevertheless, we compared our results with the experimental ratios from Ref. and found that there is a reasonable agreement, even if, in some cases, the ratios measured by the different experiments differ a lot one from the other.
## III Nuclear Effects
In Section 2 we discussed the reaction $`\nu +Nl+N^{}+\pi ^{\pm ,0}`$, where $`N`$ is a free nucleon (proton or neutron). In order to investigate the nuclear effects taking place in the experimental targets (for example, $`{}_{8}{}^{}O_{}^{16}`$, $`{}_{18}{}^{}Ar^{40}`$ or $`{}_{26}{}^{}Fe^{56}`$), we need to study the modifications necessary for the reaction $`\nu +Tl+T^{}+\pi ^{\pm ,0}`$, where $`T`$ is the nuclear target and $`T^{}`$ is an unobserved final nuclear state.
We visualize the reaction as a two step process with the neutrino interacting with individual nucleons producing single pions and excited nuclei. The production process is influenced by the Pauli principle and the Fermi-motion of individual nucleons. The subsequent journey of the pions is a “random-walk” of multiple scattering until the pion escape from the nucleus. In the multiple scattering the pions can exchange their charge. These phenomena have been studied in the past and we shall adopt the formalism in order to calculate the energy spectra of the pions. We give enough details so that the reader has an overview of the model but for more details he or she should consult the original article.
### A Charge Density Distributions
Following Ref., we treat the target as a collection of independent nucleons which are distributed in space accordingly to a density profile determined through electron-nucleus scattering experiments. For the charge density profile of $`{}_{8}{}^{}O_{}^{16}`$ we adopt the harmonic oscillator model in which the density is given by:
$$\rho (r)=\rho (0)\mathrm{exp}(r^2/R^2)\left(1+C\frac{r^2}{R^2}+C_1\left(\frac{r^2}{R^2}\right)^2\right)$$
(25)
where $`R=a/K`$ with $`K=\sqrt{3(2+5C)/2(2+3C)}`$ and $`a`$ the root mean square radius.
For $`{}_{18}{}^{}Ar^{40}`$ and $`{}_{26}{}^{}Fe^{56}`$ we use the two parameters Fermi model and write the charge density in the following way:
$$\rho (r)=\rho (0)\left(1+\mathrm{exp}((rC)/C_1)\right)^1.$$
(26)
The different parameters used in Eq.\[3.1,3.2\] are given in Ref. and are summarized in Table 1.
### B Charge Exchange
In the first step the neutrinos interact with the bound protons and neutrons with these reactions allowed provided that the energy of the recoiling nucleon is above the Fermi sea. This brings a correction factor calculated in . In the second step the pions rescatter several times until they reach the surface of the nucleus and escape. It is important to notice that the process taking place during the rescattering depend only on the properties of the target nucleus and are independent of the leptons involved in the first step. The differential cross sections for leptonic pion production on nuclear and on free nucleon targets are related to each other through the so called charge exchange matrix $`M`$ in the following way:
$$\left(\begin{array}{c}\frac{\mathrm{d}\sigma (_ZT^A;+)}{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }\\ \frac{\mathrm{d}\sigma (_ZT^A;0)}{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }\\ \frac{\mathrm{d}\sigma (_ZT^A;)}{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }\end{array}\right)=M\left(\begin{array}{c}\frac{\mathrm{d}\sigma (N_T;+)}{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }\\ \frac{\mathrm{d}\sigma (N_T;0)}{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }\\ \frac{\mathrm{d}\sigma (N_T;)}{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }\end{array}\right)$$
(27)
where
$$\frac{\mathrm{d}\sigma (N_T;\pm 0)}{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }=Z\frac{\mathrm{d}\sigma (p;\pm 0)}{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }+(AZ)\frac{\mathrm{d}\sigma (n;\pm 0)}{\mathrm{d}Q^2\mathrm{d}W\mathrm{d}E_\pi }.$$
(28)
Its eigenvalues define beams of pions of specific charge combination, which offer a scattering they produce a beam, which is decreased by the appropriate eigenvalue. The complete scattering phenomenon is characterized by three transition probabilities f($`\lambda `$) corresponding to the three eigenvalues. They describe the probabilities of beams with eigenvalues $`\lambda `$ = 1, $`\frac{5}{6}`$ and $`\frac{1}{2}`$ to survive and exit the nucleus. The “random-walk” of the pions is a stochastic process and several solutions were found in Ref.. We collect here the main formulas for the calculation. When we transpose the final pion to the initial state, we obtain the system $`\pi _i+\overline{\pi }_f`$ whose total isospin can be 0,1 and 2. We use this property to parameterize the charge exchange matrix in terms of three functions A, c and d
$$M=A\left(\begin{array}{ccc}1cd& d& c\\ d& 12d& d\\ c& d& 1cd\end{array}\right).$$
(29)
The overall factor A is given by
$$A=g(W,Q^2)f(1)$$
(30)
with g(W,$`Q^2`$) being the Pauli suppression factor and f(1) the transmission coefficient for the state with eigenvalue 1. The others two functions are
$$c=\frac{1}{3}\frac{1}{2}\frac{f(5/6)}{f(1)}+\frac{1}{6}\frac{f(1/2)}{f(1)}$$
(31)
$$d=\frac{1}{3}\left(1\frac{f(1/2)}{f(1)}\right).$$
(32)
As mentioned already f($`\lambda `$) contains the dynamics of the multiple scattering for the $`\lambda `$ eigenvalues. Both the Pauli factor and several models f($`\lambda `$) were presented in Ref.. It was also shown that approximating the multiple-scattering with scatterings in the forward and backward directions provides a very accurate approximation. Thus what is important is the effective profile of the nucleus that the pions sea. This allows one to write
$$f(\lambda )=\frac{_0^{\mathrm{}}dbbL(b)f(\lambda ,L(b))}{_0^{\mathrm{}}dbbL(b)}$$
(33)
where $`b`$ is the impact parameter and the effective length $`L(b)`$ is given by:
$$L(b)=\frac{1}{\rho (0)}_{\mathrm{}}^+\mathrm{}dz\rho \sqrt{(z^2+b^2)}.$$
(34)
In the case of $`{}_{8}{}^{}O_{}^{16}`$, the effective length $`L(b)`$ is given by:
$$L(b)=R\sqrt{\pi }\mathrm{exp}b^2/R^2\left(1+C\left(\frac{1}{2}+\frac{b^2}{R^2}\right)\right)$$
(35)
while for $`{}_{18}{}^{}Ar^{40}`$ and $`{}_{26}{}^{}Fe^{56}`$ $`L(b)`$ is written as:
$$L(b)=\frac{2}{3}\left(\frac{C}{C_1}\right)^3\left(\pi ^2\left(\frac{C}{C_1}\right)+\left(\frac{C}{C_1}\right)^36\underset{n=1}{\overset{\mathrm{}}{}}\frac{\left(\mathrm{exp}(x)\right)^n}{n^3}\right).$$
(36)
The appropriate expression for the function $`f(\lambda ,L(b))`$ has been derived in Ref. both for the case of a one-dimensional multiple scattering problem and for the case of a spherical one. The two solutions have then been compared showing excellent agreement over the entire range of parameters. Therefore, in this paper, we adopted for $`f(\lambda ,L(b))`$ the approximate expression obtained in Ref. for the one-dimensional problem.
### C Averaging Approximation
It is important to notice that, while the Pauli production factor depends on both $`W`$ and $`Q^2`$, the function $`f(\lambda )`$ depends only on W and this dependence is very weak. Therefore, as it has been already verified in Ref., it is reasonable to average the charge exchange parameters over the leading $`W`$-dependence by defining an averaged function $`\overline{f}(\lambda )`$ in the following way:
$$\overline{f}(\lambda )=\frac{dWq(W)^1\sigma _{3,3}(W)f(\lambda )}{dWq(W)^1\sigma _{3,3}(W)}$$
(37)
where $`\sigma _{3,3}(W)`$ is the pion-nucleon scattering cross section and $`q(W)`$ is the pion momentum. For the definitions of $`\sigma _{3,3}(W)`$ and $`q(W)`$ see Appendix C of Ref..
### D Results
Using the model outlined in the previous subsections, we evaluated the nuclear corrections for leptonic pion production on three different nuclei: oxygen, argon and iron.
The values obtained for $`f(\lambda )`$ and for $`g(W,Q^2)`$ for the three different targets under consideration are listed in Table II and Table III, respectively. Notice that, as Table III shows, the reduction of the cross section due to Pauli exclusion principle is larger for smaller values of $`Q^2`$ and it does not depend on the material.
The charge exchange matrices M for oxygen, argon and iron are given by:
$$M(_8O^{16})=A\left(\begin{array}{ccc}0.782& 0.161& 0.057\\ 0.161& 0.677& 0.161\\ 0.057& 0.161& 0.782\end{array}\right)$$
(38)
$$M(_{18}Ar^{40})=A\left(\begin{array}{ccc}0.731& 0.187& 0.082\\ 0.187& 0.625& 0.187\\ 0.082& 0.187& 0.731\end{array}\right)$$
(39)
$$M(_{26}Fe^{56})=A\left(\begin{array}{ccc}0.718& 0.194& 0.088\\ 0.194& 0.612& 0.194\\ 0.088& 0.194& 0.718\end{array}\right).$$
(40)
The Pauli factor and the charge exchange matrix M for oxygen have been compared with the corresponding quantities previously evaluated in Ref. and have found to be in good agreement with each other. Unfortunately, no comparison with previous calculation or experimental data is possible for argon and iron.
The differential cross sections evaluated in Sect. 2 for free nucleon targets have been used here together with the charge exchange matrix $`M`$ to obtain the differential cross sections for nuclear targets. These cross sections have been integrated over $`W`$ and $`Q^2`$ keeping the neutrino energy fix at 1 GeV in order to obtain the pion energy spectra appearing in Fig.\[8-16\].
Fig.\[8-10\] display respectively the pion energy distributions for positive, neutral and negative pions produced on oxygen targets. In each figure the solid line represent the initial distribution without any nuclear correction, the dashed line represents the same distribution after the application of the Pauli factor in the production, and the dotted line represents the final distribution after applying all the nuclear corrections discussed in the previous subsections. Similarly, Fig.\[11-13\] and Fig.\[14-16\] display the corresponding pion distributions produced on argon and iron targets, respectively.
From these figures it is clear that, while the reduction of the cross section due to the Pauli production factor is the same for all the processes investigated in this paper, the nuclear corrections related to the pion charge exchange and pion absorption are larger for neutral pions than for the positive or negative ones. Furthermore, these corrections turn out to be larger for heavier nuclei. Finally, the magnitude of the nuclear corrections decreases with increasing pion energy.
## IV Conclusions
At least three long-baseline neutrino experiments plan to study low energy neutrino reactions. Their main aim is the observation and better understanding of the neutrino-oscillations, but a necessary input is the understanding of these reactions in free protons and neutrons, as well as the modifications brought about when the nucleons are bound in relatively heavy nuclei.
In order to work in a coherent framework we calculated the cross sections on free protons and neutrons. The theory for the production of the $`\mathrm{\Delta }`$(1232)P<sub>33</sub> resonance is well understood and our results for the total cross section agree with the experimental data. The same holds for other channels where I=1/2 resonances also contribute. The comparisons appear in Figures 1-3, where it is evident that the accuracy of the measurements is subject to large improvements. Thus it is highly desirable that the new experiments use the nearby detector in order to measure the various cross sections. This refers to charged and neutral currents interactions on free protons and neutrons. The main uncertainties on this part of the paper are the functional form and parameters of the form factors and interference between I=3/2 and I=1/2 resonances. We expect that the effects from these uncertainties are small.
More important are changes which are brought about in the scattering of neutrinos in heavy nuclei. It is very likely that the far away detectors will use heavy materials as targets in order to enhance their counting rates. The heavy materials bring in corrections comparable to oscillations. In this article we used an old model for nuclear corrections and calculated the effects on the produced $`\pi ^{\pm ,0}`$. In Section 3 we reviewed the main features of the model so that the interested reader can reproduce the results.
We decided that an interesting and important parameter in the experiments is the energy of the emerging pion. We calculated in Figures 8-16 the pion spectra as function of their energy. We found that the largest correction appears in the spectrum of the $`\pi ^0`$’s. The reduction of the signal for neutral pions with energies approximately 200 MeV is substantial: of the order of 40$`\%`$. Processes with nuclear corrections as large as the ones found in this article require special attention. Several strategies suggest themselves.
One is to use the same material for the front and the far away detector and study the spectra as a function of $`E_\pi `$. Then compare the results from the two detectors and with quasi elastic scattering. In case that the experiments are forced to use different materials detailed calculations for the two materials will point to similarities and possible differences between the two targets.
###### Acknowledgements.
In the progress of this work we profited from the expertise of our colleagues. We wish to thank S. L. Adler, D. Rein, M. H. Reno, J. F. J. Salgado and L. Sehgal for helpful discussions. One of us (EAP) thanks S. L. Adler and the Institute for Advanced Study for its hospitality where part part of this work was performed. The work of LP is supported by the Deutsche Forschungsgemeinschaft (DFG) under contract GRK 54/3. The work of JYY is supported by the German Federal Ministry of Science (BMBF) under contract 05HT9PEA5.
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# A Fermi liquid model for the overdoped and optimally doped cuprate superconductors: scattering rate, susceptibility, spin resonance peak and superconducting transition
## Abstract
We present a Fermi liquid model for the overdoped and optimally doped cuprate superconductors. For the normal state, we provide an analytic demonstration, backed by self-consistent Baym-Kadanoff (BK) numerical calculations, of the linear in temperature resistivity and linear in 1/energy optical conductivity, provided the interacting Fermi liquid has strong peaks in its density of states (van-Hove singularities in 2 dimensions) near the chemical potential $`\mu `$. Recent ARPES expts. by Valla et al., Science 285, 2110 (1999), and e-print cond-mat/0003407, directly support the linearity of the one-particle scattering rate everywhere in the Brillouin zone hereto obtained. We show that the origin of this linearity is the linear in energy term of the imaginary part of the carrier susceptibility. Moreover, we verify that the interactions tend to pin the van-Hove singularities close to $`\mu `$. We show that the low energy dependence of the susceptibility can have a purely fermionic origin. We introduce an ansatz for the susceptibility of the carriers, which we postulate to be enhanced in an additive manner due to the weak antiferromagnetic order of the CuO<sub>2</sub> planes. Inter alia, this ansatz may explain the appearance of the spin resonance peak (observed in neutron scattering) in the normal state of the cuprates. Further, we obtain particularly high transition temperatures $`T_c`$ from our BK-Eliashberg scheme by using this ansatz: we have a $`d_{x^2y^2}`$ gap with $`T_c>120^o`$K for nearest neighbour hopping $`t=250meV`$.
I. Introduction
The nature of the many-body state of the cuprate superconductors is a central question for the understanding of these materials . E.g. one long standing puzzle has been the elucidation of the origin of the linear in temperature in-plane resistivity and linear in 1/energy optical conductivity observed in the optimal doping regime and to a good extent in the overdoped regime. Vice-versa, the answer to this question should shed light on the character of the carriers and, subsequently, on the superconducting transition. Here we address these issues based on a minimum unconventional Fermi liquid model . Our model comprises strong peaks in its density of states (van-Hove singularities in 2 dimensions) near the chemical potential. We show that it accounts in a natural, comprehensive and internally consistent manner for several normal state characteristics. The introduction of an ansatz for the susceptibility of the carriers further allows us both to propose an explanation for the origin of the spin resonance peak and to obtain particularly high $`d`$-wave transition temperatures $`T_c`$. Overall, our results make a strong case for a Fermi liquid approach to the optimally doped and overdoped cuprates.
We perform a combination of analytical and numerical many-body calculations in the context of our model. The rest of this paper is organized as follows. In Section II we write our many-body approximation, which we used in our numerical calculations. We emphasize both that our treatment is relevant for overdoped and optimally doped cuprates (see Section VI for the underdoped regime) and that the results presented in Sections III-V depend only quantitavely and not qualitatively on the specific Hamiltonian and approximation thereof etc. In Section III we discuss our analytical and numerical results for the linear in max($`T,ϵ`$) scattering rate of the carriers, in connection with the existence of van Hove singularities (vHs) close to the chemical potential. We also show that a conductivity linear in $`(1/T,1/ϵ)`$ follows. In Section IV we discuss both the fermionic origin of the energy dependence of the Millis-Monien-Pines susceptibility and our ansatz of eq. (26) for the susceptibility of the carriers. We show that this ansatz may explain the appearance of the so called spin resonance peak, seen in neutron scattering experiments, in the normal state of the cuprates. In Section V we discuss the superconducting transition in the frame of the ansatz. Finally Section VI contains a summary of our results. In the Appendix we examine the role of doping-induced disorder on the carrier susceptibility.
II. General framework
We assume that we deal with a Fermi liquid, albeit an unconventional one, as will become apparent from our discussion of the scattering rate of the carriers below. We choose the 2-dimensional Hubbard Hamiltonian as a specific model for our numerical calculations (c.f. the last paragraph of Section I on this).
$$H=\underset{k,\sigma }{}ϵ_kc_{k,\sigma }^{}c_{k,\sigma }+\frac{U}{2N^2}\underset{k,k^{},q,\sigma }{}c_{k+q,\sigma }^{}c_{k^{}q,\sigma }^{}c_{k^{},\sigma }c_{k,\sigma }.$$
(1)
$`c_{k,\sigma }^{}`$ is an electron creation operator and $`ϵ_k`$ is the electronic tight-binding dispersion suggested by angle-resolved photoemission (ARPES) experiments - e.g. see \- and LDA calculations
$$ϵ_k=2t(\mathrm{cos}k_x+\mathrm{cos}k_y)4t^{}\mathrm{cos}k_x\mathrm{cos}k_y2t^{\prime \prime }(\mathrm{cos}2k_x+\mathrm{cos}2k_y).$$
(2)
It is assumed here that the lattice constant is equal to unity and is the same along the two crystal axis $`a,b`$ in the planes - hence $`k_x,k_y[\pi ,\pi ]`$. $`N\times N`$ is the discretization of the Brillouin zone.
We consider the fluctuation-exchange diagrammatic approximation (FLEX) of Bickers, Scalapino and White for the Hamiltonian, which consists in summing bubble and ladder diagrams. FLEX is a Baym-Kadanoff conserving approximation , meaning that there is a free energy functional $`\mathrm{\Phi }[G]`$ of the Green’ function $`G`$, such that the self-energy $`\mathrm{\Sigma }`$ is given by the relation $`\mathrm{\Sigma }=\delta \mathrm{\Phi }[G]/\delta G`$. We thus obtain a set of self-consistent equations for $`G(k,ϵ_n)`$ and $`\mathrm{\Sigma }(k,ϵ_n)`$ :
$$G(k,ϵ_n)^1=G_o(k,ϵ_n)^1\mathrm{\Sigma }(k,ϵ_n),$$
(3)
$$G_o(k,ϵ_n)=\frac{1}{iϵ_n+\mu ϵ_k},$$
(4)
$$\mathrm{\Sigma }(k,ϵ_n)=\frac{T}{N^2}\underset{q,\omega _m}{}V(q,\omega _m)G(kq,ϵ_n\omega _m).$$
(5)
The potential $`V(q,\omega _m)`$ is given by
$$V(q,\omega _m)=V_{ex}(q,\omega _m)V_H(q,\omega _m),$$
(6)
$$V_{ex}(q,\omega _m)=\frac{U^2\chi _o(q,\omega _m)}{1U^2\chi _o^2(q,\omega _m)},$$
(7)
$$V_H(q,\omega _m)=\frac{U^3\chi _o^2(q,\omega _m)}{1U\chi _o(q,\omega _m)},$$
(8)
The susceptibility $`\chi _o(q,\omega _m)`$ is given by
$$\chi _o(q,\omega _m)=(T/N^2)\underset{ϵ_n,k}{}G(k+q,ϵ_n+\omega _m)G(k,ϵ_n).$$
(9)
$`\mu `$ is the chemical potential and the Matsubara frequencies are $`ϵ_n=(2n+1)\pi T`$ and $`\omega _m=2m\pi T`$ for fermions and bosons, respectively. We solve numerically this self-consistent set of equations, working with a given number $`M`$ of Matsubara frequencies and discretization of the Brillouin zone ($`M=256480`$ and $`N64`$).
There have been a number of similar numerical calculations on the normal-phase and the superconducting transition of the cuprates , with FLEX being a particularly popular approach.
All the convolution operations are done by using the Fast Fourier Transform (FFT), in order to cut down calculation time. We use Padé approximants to analytically continue our results to the real frequency axis.
III. On the scattering rate of the cuprates
We have analytically obtained a scattering rate linear in the maximum of the temperature $`T`$ or energy $`ϵ`$, for a Fermi liquid with strong density of states peaks - van-Hove singularities (vHs) in 2-d - located at an energy $`ϵ_{vH}`$ close to the chemical potential $`\mu `$.
The derivation relies on the relation (5) for the self-energy, which is valid quite generally in the frame of a BK approximation, irrespectively of the specific Hamiltonian and approximation thereof. See also the discussion following eq. (20). It can easily be shown that $`Im\mathrm{\Sigma }(k,ϵ)`$ is given by the following formula at finite temperature :
$$Im\mathrm{\Sigma }^R(k,ϵ)=\underset{q,\omega }{}ImG^R(q,ϵ\omega )ImV^R(kq,\omega )\{\mathrm{coth}(\omega /2T)+\mathrm{tanh}((ϵ\omega )/2T)\}.$$
(10)
Taking
$$ImG^R(k,ϵ)=\pi \delta (s_{k,ϵ}),s_{k,ϵ}=ϵ+\mu ϵ_kRe\mathrm{\Sigma }(k,ϵ),$$
(11)
we obtain
$$Im\mathrm{\Sigma }^R(k,ϵ)=\pi \underset{q}{}ImV^R(kq,s_{q,ϵ})\{\mathrm{coth}(s_{q,ϵ}/2T)+\mathrm{tanh}((ϵs_{q,ϵ})/2T)\}.$$
(12)
Setting $`ImG^R(k,ϵ)`$ equal to a delta function is a reasonable approximation for this purpose, in view of the typical sharp spike feature of $`ImG^R(k,ϵ)`$ shown in fig. 1 - also see figs. 2 and 3, all of which are representative of our numerical solution of eqs. (3)-(9). Further, numerically $`ImG^R(k,ϵ)`$ is very small compared to the band energy for small couplings, and the difference of $`Im\mathrm{\Sigma }(k,ϵ)`$, as seen in our numerical calculation, for small and large coupling constants is mostly quantitative rather than qualitative.
We write
$$ImV^R(q,x)=\underset{n=0}{\overset{\mathrm{}}{}}\frac{V_q^{(2n+1)}(0)x^{2n+1}}{(2n+1)!},$$
(13)
where $`V_q^{(n)}(0)`$ is the $`n`$th derivative of $`ImV^R(q,\omega =0)`$ with respect to $`\omega `$. This is true for an electronically mediated interaction, with a polarization which is a regular function of $`\omega `$ (see also eq. (20) below). There are only odd powers of $`\omega `$ in the series because the imaginary part of the susceptibility is an odd function of energy - e.g. c.f. eq. (2.63) of Pines and Nozières . One possible exception to this is given by Gonzalez, Guinea and Vozmediano . The authors showed that for the underdoped LSCO-type Fermi surface (FS) and for momenta $`K`$ connecting two inflection points of the FS, the imaginary part of the susceptibility goes like $`|\omega |^{1/4}`$ in 2 dimensions. However, this fact will influence the final result for the scattering rate only for a small range of momenta $`k`$ satisfying $`k=K+q_o`$ \- c.f. eqs. (14) and (15), and $`q_o`$ is given in the paragraph below. Moreover, we have already emphasized that our picture is not valid in the underdoped regime.
First we consider the low $`T`$ limit. The sum over $`q`$ is dominated by the van-Hove singularities at the points $`q_o`$. Assuming that $`ϵ_{vH}=ϵ_{q_o}+<Re\mathrm{\Sigma }(q_o,ϵ)>`$ \[this relation is misprinted in the journal version of the paper\] is close to $`\mu `$, the tanh has a vanishing contribution at the vicinity of $`ϵ_q+Re\mathrm{\Sigma }(q,ϵ)\mu `$ (note that for $`ϵ_q+Re\mathrm{\Sigma }(q,ϵ)<\mu `$ and $`ϵ_q+Re\mathrm{\Sigma }(q,ϵ)>\mu +ϵ`$ the contributions of tanh and coth annihilate each other in the low $`T`$ limit). Hence
$$Im\mathrm{\Sigma }^R(k,ϵ)\pi \underset{qq_o}{}\underset{n=0}{\overset{\mathrm{}}{}}\frac{V_{kq}^{(2n+1)}(0)(s_{q,ϵ})^{2n+1}}{(2n+1)!},$$
(14)
For sufficiently small $`V_q^{(n)}(0),n>1`$, we obtain
$$a_k=\pi \underset{qq_o}{}V_{kq}^{(1)}(0)\pi \underset{qq_o}{}\underset{n=1}{\overset{\mathrm{}}{}}\frac{V_{kq}^{(2n+1)}(0)(ϵ+c)^{2n}}{(2n+1)!},$$
(15)
where $`c=\mu ϵ_{vH}`$. This relation is valid for $`ϵ+c<ϵ_c`$, where the latter is the characteristic energy beyond which the infinite sum on the right becomes comparable to the $`V_q^{(1)}`$ term. Also, for energies beyond the bandwidth $`W`$ ($`W=8t`$ for the non-interacting system), $`Im\chi _o`$, and hence $`ImV`$ (see below), decay to zero. These considerations yield the two energy crossovers
$$ϵ_1=|\mu ϵ_{vH}|,ϵ_2=\mathrm{min}\{ϵ_c+ϵ_{vH}\mu ,W+ϵ_{vH}\mu \},$$
(16)
while the assumption above for $`a_k`$ leads to
$$Im\mathrm{\Sigma }^R(k,ϵ)a_k(ϵ+c),ϵ_1<ϵ<ϵ_2.$$
(17)
For $`ϵ>ϵ_2`$ $`Im\mathrm{\Sigma }`$ gradually decreases, due to the finite bandwidth of the system. Finally, we note that if the Fermi surface approaches a van-Hove singularity at $`q_o`$, $`V_{kq}^{(1)}(0)`$ should become bigger, being proportional to $`1/(\stackrel{}{}ϵ_{k_F}\stackrel{}{k}_F`$) (as implied by the standard Fermi liquid result for the imaginary part of the susceptibility \- see the discussion below on the susceptibility of the cuprates).
We consider now the high temperature limit $`T>(\mu ϵ_{vH})/4`$ . We see immediately that
$$Im\mathrm{\Sigma }^R(k,ϵ)=\pi \underset{q}{}ImV^R(kq,s_{q,ϵ})\{2T/s_{q,ϵ}+O(s_{q,ϵ}/2T)\}.$$
(18)
(Note that the term of order $`T`$ of this sum is reminiscent of the left-hand side of the sum rule - c.f. Pines and Nozières \- $`lim_{q0}_0^{\mathrm{}}𝑑\omega Im\chi _o(q,\omega )|\epsilon (q,\omega )|^2/\omega =N\pi /mc_s^2`$, with $`N`$ being the total particle number, $`c_s`$ the speed of sound, $`m`$ the effective mass, and $`\epsilon (q,\omega )`$ the dielectric function.) The sum is dominated by the van-Hove singularities at the points $`q_o`$, thus yielding
$$Im\mathrm{\Sigma }^R(k,ϵ)2T\pi \underset{qq_o}{}V_{kq}^{(1)}(0)=2a_kT.$$
(19)
Here we made use of the condition above for $`a_k`$. In addition, it is straightforward to see from our analytic treatment that $`Im\mathrm{\Sigma }^Rx^2`$, $`x=`$max$`\{T,ϵ\}`$, when both $`T,ϵ0`$. In all respects we have a genuine Fermi liquid.
Note added: ARPES expts. by Valla et al., Science 285, 2110 (1999), and preprint cond-mat/0003407, have very recently shown that in optimally doped Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+y</sub> the one-particle scattering rate is linear in max{$`T,ϵ`$} over most of the Fermi surface, in support of our picture.
A brief comment here. It has been known long ago - see e.g. \- that the scattering rate becomes linear in $`T`$ for $`T>\omega _B/4`$, with $`\omega _B`$ being the characteristic boson frequency mediating the carrier interaction. Our treatment shows that $`\omega _B`$ here is nothing else but the fermionic energy $`\mu ϵ_{vH}`$.
The prefactors in the r.h.s. of eqs. (17) and (19) differ by a factor of 2. This is in agreement with experiments in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> and Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+y</sub> , where the factor is found to be in the range 2.1 - 2.2 . We note that the ”marginal Fermi liquid” phenomenology of Varma, Littlewood, Schmitt-Rink, Abrahams and Ruckenstein gives a factor of $`\pi `$ instead.
We emphasize that the $`T`$ and $`ϵ`$ dependence of the result are independent of $`k`$ \- thus leading necessarily to a linear in $`T`$ resistivity and a linear in $`1/ϵ`$ optical conductivity, even with inclusion of vertex corrections in the calculation. The reason for this being that $`\{T,ϵ\}`$ are obtained as overall prefactors, for the relevant $`T`$ and $`ϵ`$ regimes, in such calculations. E.g. the Kubo formula yields $`\sigma (\omega )=(e^2/\omega )_{k,ϵ}v_k^2G(k,ϵ+\omega )G(k,ϵ)[1+S(k,ϵ)][f(ϵ+\omega )f(ϵ)]`$, where $`v_k`$ is the group velocity, $`S`$ includes vertex corrections from the Ward identity, and $`f`$ is the Fermi occupation factor. The main $`T`$ and $`\omega `$ dependence in the integrand is in the one-particle self-energy in $`G`$ and in $`f`$. Doing the $`k`$ sum, we get the dominant contribution from the poles of $`G`$. Now, the one-particle scattering rate is linear in max($`T,ϵ`$) everywhere in the Brillouin zone, and this linear dependence appears in the denominator of $`\sigma `$. Here we assumed that the vertex corrections do not have a strong temperature dependence over a substantial part of the Brillouin zone. Indeed, Kontani, Kanki and Ueda have recently shown numerically, in the frame of the FLEX approximation, that vertex corrections are small for the resistivity, and do not change its $`T`$ dependence.
Hlubina and Rice considered analytically a model of interacting fermions with a vHs close to $`\mu `$. However, they find a scattering rate similar to ours only close to the vH region, and different otherwise. As a result, their resistivity goes like $`T^2ln^2(1/T)`$. In their ’hot’ and ’cold’ spots scenario, relying on strong scattering off antiferromagnetic fluctuations, they obtain an average scattering rate similar to ours, and numerically a linear in $`T`$ resistivity (however, they seem to assume that the group velocity is finite along the whole Fermi surface - cf. between eqs. (2.5) and (2.6) of ). Similar results are also obtained in the antiferromagnetic scenario of Pines and Stojkovic .
A note on phonons. As they form a - presumably small- part of the effective potential $`V`$, they provide necessary momentum dissipation, yielding a finite resistivity. However, the linear $`T`$ dependence of the latter is not specifically influenced by phonons in our model.
Returning to the derivation of the scattering rate above, we observe that the overall behavior of $`ImV`$ closely follows $`Im\chi _o(q,\omega )`$, as
$$ImV(q,\omega )=Im\chi _o(q,\omega )|\epsilon (q,\omega )|^2,$$
(20)
$`Im\chi _o`$ is odd in $`\omega `$, while $`|\epsilon |^2`$ is even. Eq. (20) follows from any screened interaction between the carriers. Hence the argument for the linear in energy and temperature behavior of $`\tau ^1(T,ϵ)`$ is equally generic. It relies essentially on a large coefficient for the linear in energy term of $`ImV`$ \- i.e. of $`Im\chi _o`$ \- and the presence of van-Hove singularities near the Fermi surface. The result holds regardless of the dimensionality of the system. However, it is important that a significant part of the spectral weight be included in the strong peaks of the density of states lying close to $`\mu `$.
What is more, in our self-consistent numerical solution we observe that the energy $`ϵ_{vH}`$ of the singularities is pushed by the interactions close to the chemical potential \- see fig. 4. This result is especially pronounced when we use the ansatz for the susceptibility of the carriers of eq. (26) below. Then we find typically for $`n0.870.95`$ and for a broad range of $`t^{},t^{\prime \prime },U`$
$$\mu ϵ_{vH}t/20.$$
(21)
The shape of the self-energy $`\mathrm{\Sigma }(k,ϵ)`$ of the interacting system is responsible for the modification of the density of states $`N(ϵ)`$, through the relation $`N(ϵ)=TrImG(k,ϵ)/\pi `$. A trend for the transfer of the spectral weight is indicated by the fact that $`Im\mathrm{\Sigma }(k,ϵ)`$ has a peak below $`\mu `$ and a dip above it. The numerical result concerning the approachment between $`ϵ_{vH}`$ and $`\mu `$ has been known for some years. Si and Levin and Newns, Pattnaik and Tsuei observed the pinning of the vHs close to $`\mu `$ by using a $`U\mathrm{}`$ mean field slave boson approximation of a model with Cu 3d and O 2p orbitals. Recently, Gonzalez, Guinea and Vozmediano were able to obtain analytically the essential part of the approachment between the vHs and $`\mu `$ with a first order renormalization group treatment in the context of the Hubbard model. A review of related work in the frame of the so-called van-Hove scenario has been given by Markiewicz . This pinning of the vHs close to $`\mu `$ seems to be a plausible explanation for the common characteristic of a good many cuprates whose van-Hove singularities are located between 10-30 $`meV`$ below the Fermi surface (see also the next section).
It is interesting that the electron doped Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4+δ</sub> which has a van-Hove singularity much below the Fermi surface, i.e. at approximately $`\mu `$-350 meV, as shown by ARPES, has a usual Fermi liquid $`\tau ^1(T)=const.T^2ln(T)`$ . This lends support to the picture described above. Along the same line, the resistivity of Tl<sub>2</sub>Ba<sub>2</sub>CuO<sub>6+δ</sub> (Tl-2201) switches over from linear to quadratic with increasing doping from the optimal to the overdoped regime , which we suspect to be an indication of the vHs moving well away from $`\mu `$.
Finally, in further support of the relevance of the vHs in the transport properties of the cuprates, Newns et. al and McIntosh and Kaiser have shown that the thermal conductivity of the cuprates can be well accounted for if the vHs are located very close to the Fermi surface, as discussed above.
The numerical solution of the many-body system always corroborates our analytical result for the self energy at finite temperature. $`Im\mathrm{\Sigma }(k,ϵ)`$ turns out to be essentially linear in energy in the interval $`ϵ_1<ϵ<ϵ_2`$. A linear dependence of $`Im\mathrm{\Sigma }(k,ϵ)`$ as a function of either $`T`$ or $`ϵ`$ was also obtained in the numerical work of Beere and Annett , Kontani et al. and Si and Levin . Note that in figs. 2 and 3 we show the self-energy for the set of the system parameters which yields the highest transition temperature $`T_c`$, if use of the ansatz of eq. (26) is made. The linearity of $`Im\mathrm{\Sigma }(k,ϵ)`$ with $`ϵ`$ is even more pronounced for other combinations of $`t^{},t^{\prime \prime }`$ and $`n`$. $`Im\mathrm{\Sigma }(k,ϵ)`$ has always the correct parabolic Fermi liquid bevahior for $`ϵ0`$. Furthermore, the energy interval of linear behavior expands as the energy $`ϵ_{vH}`$ of the (extended) van-Hove singularities at the (vicinity of the) points $`q_o=(\pm \pi ,0),(0,\pm \pi )`$ approaches $`\mu `$.
Another feature of the density of states as seen in our treatment \- c.f. fig. 4 - is the following. The non-interacting density of states has two minor peaks at the bottom and top of the spectrum respectively (in fig. 4 the top one is a vHs). As the strength of the interaction increases, these two peaks are washed out, as a result of the self-energy which becomes substantial in magnitude for energies away from $`\mu `$ \- c.f. figs. 2 and 3.
At the moment it is not clear whether the present mechanism of the linear scattering rate can explain the experimentally observed $`T^3`$ dependence of the Hall resistivity of the cuprates. A way to explain it has been found by Stojkovic and Pines , using an electron interaction peaked at $`Q=(\pm \pi ,\pm \pi )`$. Their argument can be slightly modified, so that it works for our form of the electron potential $`V`$ \- given by eq. (6) - but with a modified effective $`\chi _o`$ peaked at Q, as we propose in the next section \- c.f. eq. (26) and below. Kontani et al. have shown that vertex corrections in the frame of FLEX have a drastic influence on the $`T`$ dependence of the Hall resistivity, in marked contrast to the case of the longitudinal resistivity.
IV. On the susceptibility of the cuprates
The low energy dependence of the susceptibility of the cuprates. The Millis-Monien-Pines susceptibility
$$\chi _{MMP}(q,\omega )=\frac{X_1\xi ^2}{1+\xi ^2(qQ)^2i\omega /\omega _{SF}},$$
(22)
has been used to fit the low energy part of the susceptibility of the cuprates in both NMR rate and inelastic neutron scattering (INS) experiments. Here $`Q=(\pm \pi ,\pm \pi )`$. The short range antiferromagnetic (AF) order, a remnant of the parent antiferromagnetic materials, with correlation length $`\xi `$, is responsible for the peak of the susceptibility for $`q`$ near $`Q`$. Typically $`\xi `$ is of the order of the lattice constant ($`\xi `$ decreases as the doping increases, and e.g. $`\xi 2`$ for optimally doped YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>), while $`\omega _{SF}1040meV`$.
The origin of the small magnitude of $`\omega _{SF}`$ has remained elusive thus far. E.g. Sachdev, Chubukov and Sokol have interpreted it as a damped spin wave mode . Spin waves are clearly observable in underdoped cuprates. However to date there is no experimental proof that they are strong enough in the normal phase of the optimally doped and overdoped regimes. The proximity of the system to an antiferromagnetic instability, i.e. $`\overline{V}_Q\chi (Q,\omega )1`$, where $`\overline{V}_q`$ and $`\chi (q,\omega )`$ are some appropriate coupling and susceptibility respectively, can in principle explain the small magnitude of $`\omega _{SF}`$, as pointed out by Millis, Monien and Pines .
Here we propose that an alternative explanation - which may coexist with the latter - is the following fermionic origin for $`\omega _{SF}`$.
First however, let us present a proposal for the susceptibility which has been put forward by Onufrieva and Rossat-Mignod. This starts by viewing the CuO<sub>2</sub> planes as a lattice of plaquettes centered on the copper site with four nearest neighbour oxygen sites. A Hamiltonian $``$, reminiscent of but more comprehensive than the one of the $`tJ`$ model, was introduced in terms of the Hubbard operators. In this formulation, the itinerant carriers which propagate via Cu spin flips are clearly separate objects from the localized Cu spins with short range AF order. A diagrammatic approach was developed in the frame of $``$, leading to the following RPA-type total susceptibility
$$\chi _t(q,\omega )=\frac{\chi _{AF}(q,\omega )+\chi _F(q,\omega )}{1+J_q(\chi _{AF}(q,\omega )+\chi _F(q,\omega ))}.$$
(23)
$`J_q`$ is the effective Cu spin exchange interaction, $`\chi _F`$ is a purely fermionic susceptibility and $`\chi _{AF}(q,\omega )`$ is due to the localized spins. $`\chi _t(q,\omega )`$ encompasses in an appealing way the idea of the entangled carrier-spin dynamics in the cuprates. Furthermore, this approach is able to account to a good extent for the variation of the total susceptibility as a function of doping and temperature, as measured by INS.
Now, we use the result for $`\chi _t`$ above with
$$\chi _{AF}(q,\omega )=\frac{\chi _1\xi ^2}{1+\xi ^2(qQ)^2f(\omega )},\chi _F(q,\omega )=\chi _{Fo}(1+i\omega /\omega _o+O(\omega ^2)),qQ.$$
(24)
Let us suppose that $`f(\omega )=i\omega /\omega _S`$. If $`\omega _S\omega _o`$ and $`J_Q\xi ^2\chi _1<1`$, and taking $`\chi _F(q,\omega )\chi _o(q,\omega )`$ (as given by eq. (9)), we essentially recover $`\chi _{MMP}(q,\omega )`$ \- which is itself an approximate form of the true susceptibility - with
$$\omega _{SF}\overline{\omega }(q)=\frac{\omega _o\omega _S}{\omega _o+\omega _SJ_q\chi _o^2(1+\xi ^2(qQ)^2)/(\xi ^2\chi _1)}.$$
(25)
From the numerical solution of our system, we easily obtain values of $`\omega _o`$ comparable to the experimentally relevant ones, when the van-Hove energy $`ϵ_{vH}`$ is near $`\mu `$, with $`\omega _o`$ scaling quickly towards zero as $`\mu ϵ_{vH}0`$. Hence $`\omega _o`$ can be interpreted as $`\omega _o(\stackrel{}{q}_F)=\stackrel{}{}ϵ_{q_F}\stackrel{}{q}_F`$ \- c.f. the non-interacting Fermi liquid result $`\omega _o(q)=v_Fq`$ . The small difference $`\mu ϵ_{vH}`$ is observed in a good number of cuprates. E.g. in ref. there is a compilation of several cuprates, the van-Hove singularities of which are located between 10 - 30 $`meV`$ below the Fermi level (c.f. the discussion in the previous section). Also, Blumberg, Stojkovic and Klein (BSK) suggested that this characteristic may be true irrespective of the doping, as long as the latter is appropriate for superconductivity. This is based on ARPES experiments on the bilayer YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub>. ARPES remains the best diagnostic probe for the Fermi surface of the cuprates. Yet it has not proved possible to perform measurements on many other compounds, especially the monolayers such as Tl<sub>2</sub>Ba<sub>2</sub>CuO<sub>6+δ</sub> etc. The point here is the following. By fitting the ARPES data BSK show that one of the two effective bands - the anti-bonding one - formed by hybridization of the two layers by interlayer coupling has a chemical potential only some 20 - 50 $`meV`$ above the van-Hove singularity at $`(0,\pi )`$, irrespective of the doping regime. It is then clear that these carriers, with a large density of states, give rise to a small $`\omega _o`$ as discussed above. Hence it is very interesting to know how universal this band-structure characteristic of the cuprates is, as it may explain naturally the magnitude of $`\overline{\omega }`$. Furthermore, it would be interesting to determine experimentally, e.g. by INS, the value of $`\overline{\omega }`$ for Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4+δ</sub>. In that case, $`\omega _o`$ should be enhanced as a result of the van-Hove singularities being far away from the Fermi surface.
In the Appendix we discuss the (non)influence of weak disorder on the value of $`\omega _o`$.
The antiferromagnetic ansatz for the carrier susceptibility. We thereby propose that the effective non-interacting susceptibility (i.e. without interaction lines connecting the particle-hole lines) of the carriers is given by the following ansatz
$$\chi _o(q,\omega )\chi _o^{eff}(q,\omega )=\chi _o(q,\omega )+a\chi _{AF}(q,\omega ).$$
(26)
$`\chi _o`$ is given by eq. (9) above, $`\chi _{AF}`$ is the antiferromagnetic susceptibility of the localized Cu spins and $`0<a<1`$ is a dimensionless weighting factor, which in principle depends on doping, band structure, temperature (presumably a decreasing function of the latter) etc. The intuitive idea is that the carriers should perceive the ordered antiferromagnetic background of the CuO<sub>2</sub> planes, even in the absence of phase separation . In this way, the fermionic susceptibility acquires an antiferromagnetic enhancement, which may then influence the effective carrier potential and pairing, through an electronically mediated interaction - see also the next section. Furthermore, $`\chi _o^{eff}`$ can explain the temperature dependence of the Hall resistivity of the cuprates, as we mentioned at the end of the previous section on the scattering rate. Finally, we note that with this effective $`\chi _o^{eff}`$ our many-body approximation remains conserving, since the relation $`\mathrm{\Sigma }=\delta F/\delta G`$ between the self-energy, the free energy and the Green’s function is still valid, as $`\chi _{AF}`$ does not depend on $`G`$. We note that this is consistent with the work of Onufrieva and Rossat-Mignod mentioned above. The carriers and the localized spins form two distinct, albeit interrelated, systems. The additive form of the ansatz is also compatible in spirit with $`\chi _t(q,\omega )`$ of eq. (23).
We mention here the alternative treatment of the spin and charge susceptibilities of Imada, Fujimori and Tokura .
Taking the ansatz of eq. (26) at face value in the context of our many body scheme means that both the charge and spin susceptibility of the carriers acquire an AF enhancement. Currently we cannot prove this ansatz. We emphasize that our ansatz can be taken to apply only for the spin susceptibility of the carriers. In that case, in the frame of our Baym-Kadanoff scheme, only $`\chi _o(q,\omega )`$ entering the ladder diagrams with opposite particle-hole spins would be replaced by $`\chi _o^{eff}`$. Quantitatively, the differerence between this case and the case in which both charge and spin susceptibilities are enhanced is small (for relevant values of the AF enhancement) - c.f. the discussion on the critical temperature $`T_c`$ in the next section.
In our numerical implementation, we consider two similar forms for the AF susceptibility, namely (A)
$`\chi _{AF}^A(q,\omega _m)=X_o{\displaystyle \underset{i=1}{\overset{4}{}}}\mathrm{\Gamma }_i^1\theta (\omega _c|\omega _m|),`$ (27)
with $`\mathrm{\Gamma }_i=\xi ^2+(qQ_i)^2`$, $`\omega _c`$ being a cut-off, and (B)
$$\chi _{AF}^B(q,\omega _m)=X_oD\underset{i=1}{\overset{4}{}}\frac{\omega _m(2\omega _m/\pi )\mathrm{arctan}(\omega _m/D)\mathrm{\Gamma }_iD+(2\mathrm{\Gamma }_iD/\pi )\mathrm{arctan}(\mathrm{\Gamma }_i)}{\omega _m^2(\mathrm{\Gamma }_iD)^2},$$
(28)
with $`Q_i=(\pm \pi ,\pm \pi )`$ and $`D`$ being a cut-off frequency, above which $`Im\chi _{AF}(q,\omega )=0`$. Form (B) has appeared in . Here the characteristic spin wave frequency obeys $`\omega _S\xi ^z`$, and the $`z=2`$ scaling regime has been assumed, in agreement with the analysis of Sokol and Pines for the optimally doped and overdoped regime of YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> etc.
The spin resonance peak of the cuprates. INS experiments in YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> \- see e.g. and therein - and Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+y</sub> have revealed the existence of a strong peak in the spin triplet channel, centered at $`Q=(\pm \pi ,\pm \pi )`$ and at a characteristic resonance energy $`\omega _R3040`$ meV. Although this peak is usually seen exclusively in the superconducting state, it has been observed up to temperatures $`250^o`$K in YBa<sub>2</sub>(Cu<sub>0.995</sub>Zn$`{}_{0.005}{}^{})_3`$O<sub>7</sub> for $`\omega _R40`$ meV. Interestingly, this fact cannot be accounted for by most of the theoretical models so far available - see e.g. for refs. - as these models require the onset of superconductivity. An exception is the model of Bulut , which, however, differs drastically from ours.
The use of the ansatz of eq. (26) in the spin channel only may account in a natural way for the appearance of the resonance peak in the normal state through a bilayer effect. Both YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7-δ</sub> and Bi<sub>2</sub>Sr<sub>2</sub>CaCu<sub>2</sub>O<sub>8+y</sub> are bilayer materials. Here one can define bonding and antibonding bands $`ϵ_{\pm k}=ϵ_k\pm t_{}(k)`$, $`t_{}(k)`$ being the $`k`$-dependent interlayer hopping element. Further, one defines the susceptibilities $`\chi _{o\pm }(q,\omega )=Tr[G_+(+q,+\omega )G_\pm +G_{}(+q,+\omega )G_{}]`$. When $`ϵ_{vH}\mu `$ is small, $`\chi _o(q,\omega )`$ has a narrow peak at $`q=Q`$ and $`\omega \mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is the bilayer splitting at $`qq_o=(\pi ,0),(0,\pi )`$, i.e. the van Hove neighbourhood. Liechtenstein et al. have shown in the frame of the FLEX approximation that $`\mathrm{\Delta }`$ is drastically reduced for a finite interaction, compared to the non-interacting value $`\mathrm{\Delta }_o=2t_{}(q=q_o)`$, such that it becomes comparable to the experimental $`\omega _R`$. Our ansatz can be taken to apply to $`\chi _{o\pm }`$, yielding $`\chi _{o\pm }\chi _{o\pm }^{eff}=\chi _{o\pm }+a\chi _{AF}`$. As a result $`\chi _{}(q,\omega )=\chi _o^{eff}(q,\omega )/(1U\chi _o^{eff}(q,\omega ))`$ is strongly peaked at $`Q`$ for an energy $`\omega _{cR}\mathrm{\Delta }`$, and may account for the experimental observations. Of course we require $`\chi _o>\chi _{o+}`$ for this to work, which is valid for a range of the parameters $`t^{},t^{\prime \prime },n`$.
One can ask the question: why does the resonance peak not appear in the normal state in general? As demonstrated here, the amplitude of the AF enhancement needs to be sufficiently large for the peak to be visible. Zn is known to enhance AF fluctuations in the CuO<sub>2</sub> planes - e.g. , which if interpreted as yielding a larger ($`a\chi _{AF}`$) contribution in $`\chi _o^{eff}`$, is in agreement with the results above. On the other hand, as we discuss in the next section on the superconducting transition, too strong a factor $`(a\chi _{AF})`$ reduces $`T_c`$ \- c.f. $`T_c=93^o`$K for the pure YBa<sub>2</sub>Cu<sub>3</sub>O<sub>7</sub> versus $`T_c=87^o`$K for the Zn doped material mentioned above (also see the last paragraph of section V).
V. On the superconducting transition of the cuprates
With the solution of the normal system at hand, we solve the gap (Eliashberg) equation
$$\mathrm{\Delta }(k,ϵ_n)=\frac{T}{N^2}\underset{k^{},ϵ_n^{}}{}V_p(kk^{},ϵ_nϵ_n^{})G(k^{},ϵ_n^{})G(k^{},ϵ_n^{})\mathrm{\Delta }(k^{},ϵ_n^{}).$$
(29)
This form of the equation is valid close to the transition temperature $`T_c`$ only. The pairing potential $`V_p`$ is given by
$$V_p(q,\omega _m)=V_x(q,\omega _m)+V_h(q,\omega _m),$$
(30)
$$V_x(q,\omega _m)=\frac{U}{1U^2\chi _o^2(q,\omega _m)},$$
(31)
$$V_h(q,\omega _m)=\frac{U^2\chi _o(q,\omega _m)}{1U\chi _o(q,\omega _m)}.$$
(32)
We have assumed that the gap is an even function of both its momentum and energy arguments in order to write $`V_p(q,\omega _m)`$ in this form. We also assumed that the gap is spin singlet, as Knight shift experiments have shown . The superconducting state of the cuprates is probably a generalized BCS state - see e.g. \- with the transition being due to a momentum anisotropic potential, as envisaged by Kohn and Luttinger .
The symmetry of the gap is determined by the exact shape and sign of $`V_p`$. Based on symmetry grounds as well as on experimental evidence, we expect a $`d_{x^2y^2}`$ or a $`s`$ wave gap - but see also . (In principle, higher order even angular momentum harmonics are also possible.) The highest $`T_c`$’s correspond to a $`d_{x^2y^2}`$ gap, and are obtained by including an antiferromagnetically enhanced susceptibility in the calculation, following our ansatz of eq. (26). In passing, let us note that in general the proximity of $`ϵ_{vH}`$ to $`\mu `$ plays a less significant role than the ansatz in raising $`T_c`$ \- but c.f. below. Nevertheless we also obtain a $`s`$-wave gap under the same conditions, but with a much lower $`T_c^{}`$ . This is consistent with the experimental situation. Most of the cuprates appear to have a $`d_{x^2y^2}`$ gap at $`T_c`$ \- e.g. c.f. . Experiments point to the opening of a secondary order parameter gap at $`T_c^{}T_c`$ .
We note here that the van-Hove singularities at $`q_o`$ tend to suppress a $`s`$-wave gap if the pairing potential $`V_p>`$0 ($`V_p`$ defined above is negative for sufficiently large and negative $`U`$ and/or an appropriate phonon coupling, and positive otherwise). The (likely) $`s`$-wave gap of Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4+δ</sub> could originate from e.g. the fact that the van-Hove singularities are 350$`meV`$ below the Fermi surface, and hence ineffective here, or from $`V_p<`$0 for some relevant parts of the phase space in this material, or possibly from both facts. Yet another possibility which may coexist with the above is that $`V_p`$ is strongly peaked close to zero momentum owing to the band structure of this material. All these factors can result in the $`s`$-wave $`T_c`$ being higher than the $`d`$-wave $`T_c`$, and hence in the dominance of the former channel over the latter.
We obtain the following (near) optimum set of parameters for the $`d`$-wave $`T_c`$. We take $`t=250meV`$ and $`\xi =1`$ (increasing the latter leads to a reduction of $`T_c`$, see below). For form (A) - eq. (27) - we obtain $`T_c125^oK`$, for $`aX_o=0.2eV^1`$, $`\omega _c=4t`$, $`t^{}=0.11t`$, $`t^{\prime \prime }=0.5t`$, $`U=1.27eV`$ and $`n=0.88`$. For form (B) - eq. (28) - we obtain $`T_c105^oK`$, for $`aX_o=0.5eV^1`$, $`D=32t`$, $`t^{}=0.11t`$, $`t^{\prime \prime }=0.25t`$, $`U=1.47eV`$ and $`n=0.91`$. The variation of $`T_c`$ in $`{}_{}{}^{o}K`$ with $`\xi `$ is as follows for this last set of parameters : (94, 1.5), (91, 2), (90, 3), (88, 5).
Monthoux and Pines obtained similarly high transition temperatures with their approach, which uses $`𝒱(q,\omega )=g^2\chi _{MMP}(q,\omega )`$ as the effective carrier-carrier potential. Their approach has no room for the Hubbard $`U`$ and the subsequent screened carrier-carrier interaction though. Overall, our approach yields significantly enhanced $`T_c`$’s compared to the standard FLEX-type treatments .
We obtain an optimum value of $`n`$ for the following reason. $`\chi _o`$ \- and hence $`\chi _o^{eff}`$ \- is a decreasing function of $`n`$. We allow the coupling to increase up until $`b=U\chi _o^{eff}(Q,\omega =0)`$ saturates to a value close to and below unity. Then $`V`$ and $`V_p`$ (given by eqs. (6) and (32) respectively) are increasing with $`n`$. And so does the characteristic scattering rate entering $`G`$. Thus the optimum value of $`n`$ for the highest $`T_c`$ corresponds to the overall highest kernel $`V_p|G|^2`$ in the gap equation (29).
Further, we have done an extensive search of the parameter space to locate the optimum parameter set for the highest $`T_c`$. From our data it appears that the variation of $`T_c`$ as a function of the system parameters is smooth and that the optimum parameter set above is a global optimum .
Assuming that only the spin susceptibility acquires an AF enhancement according to our ansatz, yields a $`T_c`$ which is lower by 9% for $`\xi =1`$, but only lower by 2% for $`\xi =3`$, if we make use of $`\chi _{AF}^B`$, with the optimum parameters above. The variation of $`T_c`$ with $`\xi `$ is as follows here : (96,1), (92,1.5), (90,2), (89,3), (88,5).
There is an optimum value of $`a`$ for the maximum attainable $`T_c`$. This is again due to the form of the pairing potential $`V_p`$ above : the AF instability condition allows $`U\chi _o^{eff}<1`$ only. Now, for a given value of the latter product, the highest $`V_p`$ \- which in principle yields the highest $`T_c`$ as well (c.f. above) - will correspond to the highest possible $`U`$. This in turn corresponds to a smaller $`\chi _o^{eff}`$, and hence to a small but finite optimum $`a`$. In the next paragraph we discuss relevant experimental evidence.
Zheng, Kitaoka, Ishida and Asayama made a very interesting empirical observation. Namely, among the hole doped cuprates, the highest $`T_c`$ ’s correspond to a combination of both optimum total carrier concentration $`n_{x^2y^2}+2n_{p_\sigma }`$ in the planes ($`n_{x^2y^2}`$ being the concentration of holes of Cu-3d orbital character and $`n_{p_\sigma }`$ of O-2p orbital character) as well as of a reduced (probably minimum) imaginary susceptibility - as deduced from NMR experiments. In fact, Zheng et al. noticed that the highest $`T_c`$ ’s are obtained for a reduced ratio $`n_{x^2y^2}/2n_{p_\sigma }1`$, and, moreover, that such a trend is correlated with a reduced relaxation rate $`(1/T_1)_{Cu}Tlim_{\omega 0}_qA_q^2Im\chi (q,\omega )/\omega `$, in properly normalized units (here $`A_q`$ is the hyperfine coupling). This last fact points to a reduced $`Im\chi (q,\omega )`$ etc. These conclusions are in accordance with our picture, which yields both special values of the filling factor $`n`$ \- for this also c.f. e.g. \- as well as special small values of the product $`(a\chi _{AF})`$ as a prerequisite for the highest attainable $`T_c`$ ’s (also c.f. the last paragraph of section IV).
VI. Summary
To summarize, we present a single plane Fermi liquid model which for the normal state can explain the salient transport properties, the low energy dependence of $`\chi _{MMP}`$, and their relation to the existence of van-Hove singularities close to the Fermi surface. E.g. we obtain analytically a scattering rate linear in max($`T,ϵ`$), within appropriate $`T`$ and $`ϵ`$ bounds, for all momenta in the Brillouin zone. This result yields directly a linear in $`T`$ resistivity and linear in 1/$`ϵ`$ optical conductivity. The introduction of an ansatz for the susceptibility of the carriers allows for an understanding of both the appearance of the spin resonance peak in the normal state and the temperature behaviour of the Hall conductivity. Further, by using the ansatz we obtain significantly enhanced $`d_{x^2y^2}`$ wave transition temperatures $`T_c`$. Attention is paid to Nd<sub>2-x</sub>Ce<sub>x</sub>CuO<sub>4+δ</sub>, the properties of which, despite appearances, we believe to be fully consistent with those of the majority of cuprates.
In brief, let us discuss the possible connection to the physics of the underdoped cuprates. Strong experimental evidence suggests that they are in a phase separated regime, with AF domains of spins separated by stripes of holes . One can envisage that with doping increasing towards the optimal regime, the stripes melt into an effective Fermi liquid, and the physics described here is recovered. Further models on the underdoped cuprates can be found in .
The author has enjoyed discussions with Yia-Chung Chang, Gordon Baym, Joseph Betouras, Girsh Blumberg, Antonio Castro Neto, Lance Cooper, Sasha Liechtenstein, Peter Littlewood, Jörg Schmalian, Raivo Stern, Qimiao Si and Branko Stojkovic. This work was supported by the Research Board of the University of Illinois, the Office of Naval Research under N00014-90-J-1267 and NSF under DMR-89-20538.
Appendix - Disorder effects in the susceptibility of the carriers
Returning to the origin of a small $`\omega _o`$, which was discussed in section IV, another option would in principle be the disorder inherent in the cuprates. The dopants are randomly positioned in the crystal structure, thereby creating an effective disorder potential for the carriers in the planes. We calculated the effect of disorder by considering the dopants as isotropic point scatterers, with a density $`n_i`$=1% and typical scattering strength $`V_s=01`$ eV (i.e. $`8t`$). These parameters give a residual impurity scattering rate less than 2 $`meV`$, consistent with experiments on the cuprates. For the calculation of the susceptibility we used the diffuson . We only used a band structure with $`t^{\prime \prime }=0`$, and only took into account non-magnetic disorder. The non-interacting Green’s function now becomes
$$G_o^{}(k,ϵ_n)=\frac{1}{iϵ_n+\mu ϵ_k+\sigma _i(ϵ_n)},$$
(33)
with $`\sigma _i(ϵ_n)=n_iV_s/(1(V_s/N^2)_kG_o(k,ϵ_n))`$. The susceptibility is given by
$$\chi _o(q,\omega _m)=T\underset{ϵ_n}{}P(q,\omega _m;ϵ_n)\left\{\frac{\theta (ϵ_n(ϵ_n+\omega _m))}{1n_iV_s^2P(q,\omega _m;ϵ_n)}+\theta (ϵ_n(ϵ_n+\omega _m))\right\},$$
(34)
and $`P(q,\omega _m;ϵ_n)=(1/N^2)_kG(k+q,ϵ_n+\omega _m)G(k,ϵ_n)`$. However we found no evidence, for the parameters above, of $`\omega _o`$ being influenced by disorder.
Current address: 18 Giampoudi St., Iraklio, Crete 71201, Greece. E-mail: kast@iesl.forth.gr .
FIGURE CAPTIONS
Figure 1. Full Green’s function $`G^R(k_F,ϵ)`$ for (a) $`k_F=k_F^{boundary}`$ along $`(\pi ,0)(\pi ,\pi )`$ and (b) $`k_F=k_F^{diagonal}`$ along $`(0,0)(\pi ,\pi )`$, for $`t=250meV`$, $`t^{}=0.11t`$, $`t^{\prime \prime }=0.25t`$, $`U=1.5eV`$, $`n=0.91`$, at $`T=105^oK`$.
Figure 2. Self-energy $`\mathrm{\Sigma }(k_F,ϵ)`$ for $`k_F=k_F^{boundary}`$ along $`(\pi ,0)(\pi ,\pi )`$ (continuous line) and $`k_F=k_F^{diagonal}`$ along $`(0,0)(\pi ,\pi )`$ (dashed line), for the same parameters as in fig. 1. (a) depicts Im$`\mathrm{\Sigma }^R(k_F,ϵ)`$ and (b) depicts Re$`\mathrm{\Sigma }(k_F,ϵ)`$. A (quasi)linear energy dependence of Im$`\mathrm{\Sigma }^R(k_F,ϵ)`$ can be seen for energies below 0.5$`eV`$ \- also c.f. text.
Figure 3. Self-energy $`\mathrm{\Sigma }(k_F,ϵ)`$ for $`k_F=k_F^{boundary}`$ along $`(\pi ,0)(\pi ,\pi )`$ (continuous line) and $`k_F=k_F^{diagonal}`$ along $`(0,0)(\pi ,\pi )`$ (dashed line), for $`t=250meV`$, $`t^{}=0.11t`$, $`t^{\prime \prime }=0.25t`$, $`U=1.47eV`$, $`n=0.91`$, at $`T=105^oK`$. The carrier susceptibility has an antiferromagnetic enhancement according to the ansatz of eq. (26) here (see section IV), with $`\xi =1`$, $`aX_o=0.5eV^1`$ and $`D=32t`$. This is the optimum $`T_c`$ case when using form (B) - eq. (28) - of our ansatz. (a) depicts Im$`\mathrm{\Sigma }^R(k_F,ϵ)`$ and (b) depicts Re$`\mathrm{\Sigma }(k_F,ϵ)`$. A (quasi)linear energy dependence of Im$`\mathrm{\Sigma }^R(k_F,ϵ)`$ can be seen for energies below 0.5$`eV`$.
Figure 4. Evolution of the density of states. Dashed line : non-interacting system, with parameters as in fig. 3. Continuous line : same system with $`U=0.8eV`$. Notice the transfer of the central van-Hove peak towards the chemical potential, the disappearance of the two satellite peaks - see text, and the broadening of the total spectral width.
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# Interband absorption and luminescence in small quantum dots under strong magnetic fields
## I Introduction
The quasi-two-dimensional electron gas in a high magnetic field is a strongly correlated system exhibiting very complicated dynamics. At special values of the filling factor, the essential features of the ground state are captured by the Laughlin wave function , or its composite fermion generalization. The low-lying excitations can be described in the single-mode approximation of Girvin et al, or in the composite fermion picture .
Many experiments have been designed and carried out in order to test the excitation spectrum of this highly correlated system. Inelastic (Raman) light scattering experiments have tested basically the excitation gap at wavevector $`𝐤=0`$ Spin-flipped states and the magnetoroton minimum at $`k1/l_B`$ ($`l_B`$ is the magnetic length) have also been observed , although they should be activated by impurities or other mechanism to produce a trace in the Raman spectra. Evidence of the magnetoroton minimum comes also from the absorption of ballistic acoustic phonons.
On the other hand, experiments on photoluminescence (PL) related to interband electronic transitions around filling factor $`\nu =1`$ have tested the excited states with an additional electron-hole (e-h) pair. Recently, the observations have been extended to lower filling factors by increasing the magnetic field up to 60 T. The related theory is not in complete agreement with the experiment. In the infinite magnetic field limit, it was predicted that only the exciton ($`X^0`$) and the negatively charged triplet exciton ($`X_t^{}`$) are bound . The latter is expected to be dark in luminescence as a result of a hidden symmetry related to magnetic translations . In the experiments, however, very distinct singlet and triplet peaks ($`X_s^{}`$ and $`X_t^{}`$) are observed. A realistic calculation of ground state energies was presented in Ref. , where Landau level and quantum well (qwell) sub-band mixing were taken into account. The $`X_t^{}`$ peak position was reproduced, but in theory this state is dark. The problem was recently revisited by Wojs et al, who showed that in a narrow (10 nm wide) well a second bright $`X_t^{}`$ state becomes bound, thus interpreting the observed luminescence as coming from the bright state. We shall notice that both Refs. and deal with isolated three-particle systems, and thus are not able to describe the filling factor dependence of observed magnitudes for $`\nu 1/5`$.
In the present paper, we study small quantum dots (qdots) under conditions similar to the experiments reported in Refs. , i.e. quasi-two-dimensional motion, magnetic fields in the interval $`8`$ T$`B60`$ T, and temperatures well below 2 K. The laser excitation power is assumed to be low (a few mW/cm<sup>2</sup>), thus the dot works under a linear regime. The lateral confinement is modelled by a harmonic potential. Energy levels, charge densities and dipole matrix elements for absorption and luminescence are computed by exact diagonalization in the first Landau level (1LL) approximation.
Absorption or PL experiments on electron-hole qdots under very high magnetic fields are lacking. To the best of our knowledge, there is only one experiment , in which the luminescence at higher (4 K) temperature and $`B45`$ T is measured in order to estimate the e-h correlation energy.
Breaking of the magnetic translation symmetry by a lateral confinement in a qdot makes the lowest $`X^{}`$ triplet state bright. Highly nontrivial PL and absorption spectra arise even in the 1LL approximation. These spectra contain information about the energy levels and particle correlations in the system. Let us stress that a calculation of $`X^0`$ and $`X^{}`$ energy levels of in a qdot, which includes LL mixing, is available . The absorption coefficient is also reported in that paper. The differences with our work are the following. First, we consider both absorption and PL. Second, we trace the changes in the ground-state (g.s.) wave function and charge rearrangements as the magnetic field is varied. Finally, we consider larger qdots with $`X^2`$ and $`X^3`$ complexes (unbound in a qwell). It will be seen below that indications of collective effects are evident even in these relatively small systems.
The plan of the paper is as follows. The model and certain general statements are explained in section II. The next section presents results for particular systems. We start with the exciton and end up with the $`X^3`$ complex. Finally, a few concluding remarks are given.
## II The model
We consider the two-dimensional motion of $`N_e`$ electrons and $`N_h`$ holes in an external parabolic potential and a perpendicular magnetic field (along the $`z`$ axis). In particular, we will study the $`N_h=1`$ and 2 systems, which are the ones participating in interband absorption and recombination processes. The unit of length is $`\sqrt{2}`$ times the magnetic length. In the 1LL approximation, the Hamiltonian is written as
$`H(N_e,N_h)`$ $`=`$ $`\left({\displaystyle \frac{\mathrm{}\omega _c^e}{2}}+E_z^e\right)N_e+\left({\displaystyle \frac{\mathrm{}\omega _c^h}{2}}+E_z^h+E_{gap}\right)N_h`$ (1)
$`+`$ $`E_{Zeeman}^e+E_{Zeeman}^h+V_{conf}+V_{coul}.`$ (2)
The Hamiltonian (2) is intended to model a GaAs qdot with a thickness of 20 nm in the $`z`$-direction. The meaning of the different terms entering $`H`$ is evident. The specific qdot characteristics are reflected in the confinement energies along the $`z`$-direction, $`E_z`$, the in-plane confinement potential, $`V_{conf}=_iv_{conf}(r_i)`$, and the $`z`$-averaged Coulomb interactions, $`V_{coul}=_{i,j}v_{coul}(r_{ij})`$. We will use the expression
$$v_{coul}(r)=\pm 3.316\beta \sqrt{B}\frac{1}{r}[\mathrm{meV}],$$
(3)
for the pair Coulomb interactions ($`B`$ in Teslas), and
$$v_{conf}(r)=\frac{3.316}{B}\beta Kr^2[\mathrm{meV}],$$
(4)
for the one-particle confinement potential. Even these simple expressions lead to very interesting physical results. Notice that $`1/\sqrt{2}`$ times the characteristic Coulomb energy, $`e^2/(\kappa l_B)`$, equals 3.316 $`\sqrt{B}`$ in our units. The constant $`\beta =0.6`$ is used to simulate the $`z`$-averaging of the Coulomb interactions in the 20 nm - wide qdot . We fixed it by requiring the binding energy of the unconfined ($`v_{conf}`$ set to zero) $`X_t^{}`$ relative to the $`X^0`$ to be 0.6 meV at $`B=17`$ T. This is a representative value. On the other hand, the dimensionless constant $`K`$ will be fixed to 7.0 in order to obtain a “filling factor” around 1/3 for $`B30`$ T, also a common situation met in the experiments .
The only nontrivial terms entering (2) are $`V_{conf}`$ and $`V_{coul}`$. They should be diagonalized in a basis of Slater 1LL functions. The energies coming from the diagonalization processes will be denoted $`ϵ`$, and the wave functions will be used to compute physical observables. Note that, in the 1LL, the electron (hole) angular momentum is a non-positive (non-negative) number. Thus, the total angular momentum is written $`M=M_e+M_h=|M_e|+M_h`$.
In a GaAs electron system, the validity of the 1LL approximation can be stated by comparing the excitation energy to the 2LL, $`\mathrm{}\omega _c^e=1.728B`$ meV, with the Coulomb energy, $`3.316\beta \sqrt{B}`$ meV. Thus, for $`B1`$ T the 1LL approximation works. Spin excitations are lower in energy, $`\mathrm{\Delta }E_{Zeeman}0.025B`$ meV. However, at temperatures below 2 K and for $`B>8`$ T, they can not be thermally excited. It means that in both absorption and luminescence the transition starts from the lowest optically active state. When holes are created, the 1LL approximation becomes valid at higher fields. If we take for the heavy hole mass in the $`xy`$ plane the value $`\mu _h=0.11m_0`$, then $`\mathrm{}\omega _c^h1B`$ meV. The 1LL approximation works for $`B4`$ T. Below, we present results obtained in the 1LL approximation for $`8`$ T$`B60`$ T.
On the other hand, the expression (2) assumes that the particles are sitting on the first qwell sub-band. As it was stressed in Ref. , this may be a rough approximation. For a 20 nm qwell, the second electronic sub-band is around 30 meV higher, but the second hole sub-band is only 6 meV higher (a heavy hole mass $`\mu _h^z0.38m_0`$ is assumed). Our first sub-band approximation is qualitatively valid in the present situation, and will improve for narrower wells.
### A Interband absorption and luminescence (general grounds)
Interband absorption and luminescence will be studied at temperatures $`T2`$ K, i.e. typically lower than spin excitation gaps. Thus, the processes proceed from a unique initial state, which is the g.s. of the polarized $`(N_e,0)`$ system in absorption, and the lowest optically active state of the $`(N_e+1,1)`$ system in emission. In general, these processes take place in different angular momentum channels. For absorption, the incident light is supposed to be circularly polarized and propagating along the $`z`$-direction. Also circularly polarized light is supposed to be measured from the qdot luminescence.
A simple two-band model, with bands split by the Zeeman energy, will be used. The conduction-band ($`m_s=\pm 1/2`$) mass is $`\mu _e=0.067m_0`$, and the heavy hole band, $`m_j=\pm 3/2`$, shows anisotropic effective masses, $`\mu _h=\mu _h^{xy}=0.11m_0`$, $`\mu _h^z=0.38m_0`$. LL mixing in the $`m_j=3/2`$ branch will be neglected. $`m_j=3/2`$ will be called the spin-up hole branch, and $`m_j=3/2`$ – the spin-down branch. For propagation along the $`z`$ axis, the allowed transitions are $`m_j=3/2m_s=1/2`$ for right-handed circular polarization (RHCP), and $`m_j=3/2m_s=1/2`$ for left-handed circular polarization (LHCP) .
The dipole approximation is used for the interaction Hamiltonian, i.e. $`𝐃`$. In the 1LL, the interband dipole operator takes the form
$$𝐃=\frac{e𝐩_{cv}}{m_0\omega }\underset{l0}{}\left(e_{l,}^{}h_{l,}^{}+e_{l,}^{}h_{l,}^{}\right)+H.C.,$$
(5)
where $`𝐩_{cv}`$ is the GaAs interband constant. The reason for not including the light hole in (5) is twofold. First, $`E_z`$ is around 6 meV higher ($`\mu _{lh}^z0.09m_0`$), thus its absorption or luminescence lines are shifted. Second, the constant $`𝐩_{cv}^2`$ is three times smaller for light holes. Notice that the interaction Hamiltonian preserves total angular momentum.
In our $`(N_e,N_h)`$ systems with $`N_h=0,1`$, the states may be classified according to the symmetry of the electronic subsystem. For example, the $`N_e=2`$ system may be in a spatially antisymmetric (triplet) state, or in a spatially symmetric (singlet) state. We will present calculations only for spatially antisymmetric states. They are the only ones appearing in LHCP, and the ones associated to the most intense lines in RHCP . The wave functions may be written as $`\psi =\varphi _{coord}^{antisymm}\chi _{spin}^{symm}`$, or in a second quantization formalism,
$$|\psi (N_e,0)=C_{l_1l_2\mathrm{}l_{N_e}}e_{l_1,}^{}e_{l_2,}^{}\mathrm{}e_{l_{N_e},}^{}|0,$$
(6)
$`|\psi _{LHCP}(N_e+1,1)=`$ (7)
$`{\displaystyle C_{l_1l_2\mathrm{}l_{N_e+1},l_h}e_{l_1,}^{}e_{l_2,}^{}\mathrm{}e_{l_{N_e+1},}^{}h_{l_h,}^{}|0},`$ (8)
$`|\psi _{RHCP}(N_e+1,1)={\displaystyle \frac{1}{\sqrt{N_e}}}{\displaystyle C_{l_1l_2\mathrm{}l_{N_e+1},l_h}}`$ (9)
$`\times (e_{l_1,}^{}e_{l_2,}^{}\mathrm{}e_{l_{N_e+1},}^{}`$ (10)
$`+e_{l_1,}^{}e_{l_2,}^{}\mathrm{}e_{l_{N_e+1},}^{}`$ (11)
$`+\mathrm{}+e_{l_1,}^{}e_{l_2,}^{}\mathrm{}e_{l_{N_e+1},}^{})h_{l_h,}^{}|0.`$ (12)
$`\psi _{LHCP}`$ corresponds to a spin-polarized electronic subsystem, and $`\psi _{RHCP}`$ to a not completely polarized state. In the pure electron system, the sum runs over angular momentum states obeying $`0l_1<l_2<\mathrm{}<ł_{N_e}`$ and fixed $`M=l_1l_2\mathrm{}l_{N_e}`$. In the one-hole system, the total angular momentum $`M=l_1l_2\mathrm{}l_{N_e+1}+l_h`$ is fixed.
Diagonalization of $`V_{conf}+V_{coul}`$ in (2) leads to the determination of eigenenergies and wave functions. Transition energies, transition probabilities and charge densities of the relevant states are computed from these results. The transition energies are given by
$`\mathrm{}\omega `$ $`=`$ $`E_{gap}+E_z^e+E_z^h+{\displaystyle \frac{\mathrm{}\omega _c^e}{2}}+{\displaystyle \frac{\mathrm{}\omega _c^h}{2}}`$ (13)
$`+`$ $`E_{Zeeman}^e+E_{Zeeman}^h+ϵ(N_e+1,1)ϵ(N_e,0),`$ (14)
where $`ϵ`$ are the energies coming from $`V_{conf}+V_{coul}`$. We took the values $`E_{gap}=1510`$, $`E_z^e=11`$, $`E_z^h=2`$, $`\mathrm{}\omega _c^e/2=0.864B`$, $`\mathrm{}\omega _c^h/2=0.526B`$, $`E_{Zeeman}^e=0.025m_s^eB`$, $`E_{Zeeman}^h=0.016m_s^hB`$, for the quantities entering (14), where energies are given in meV and $`B`$ in Teslas. Our treatment of Zeeman energies of both electrons and holes is very simple. We used the value $`g_e=0.44`$ for the electron Landé factor and extracted the hole energy from the observed splitting of $`X^0`$ luminescence lines in RHCP and LHCP . The hole spin projection is conventionally written as $`m_s^h=\pm 1/2`$. Actually, the Zeeman energy shows a nonlinear dependence on $`B`$. Notice, however, that $`E_{gap}`$, $`\mathrm{}\omega _c`$ and $`E_{Zeeman}`$ are important in determining the absolute position of a given absorption or PL line, but not its relative position with respect to $`X^0`$ in the same polarization.
The absorption coefficient of a dot is given by
$$\alpha (\omega )=\frac{4\pi ^2\omega }{\mathrm{}cV}\underset{f}{}|f|𝐞𝐃|i|^2\delta (\omega \omega _{fi}),$$
(15)
where $`|i`$ is the g.s. of the $`(N_e,0)`$ system, $`f`$ are the states of the $`(N_e+1,1)`$ system in the same angular momentum tower and $`\mathrm{}\omega _{fi}`$ is their energy difference computed from (14). $`𝐞`$ is the light polarization vector, $`c`$ – the light velocity, and $`V`$ is the volume of absorption. We have used a phenomenological width, $`\mathrm{\Gamma }=0.8`$ meV, to replace the delta function by a Lorentzian
$$\delta (x)\frac{\mathrm{\Gamma }/\pi }{\mathrm{\Gamma }^2+x^2}.$$
(16)
In luminescence, we compute the matrix elements $`|f|𝐞𝐃|i|^2`$, assuming that $`|i`$ is the lowest state of the $`N_h=1`$ system.
## III Results
We present results in the following interval of magnetic field values, 8 T $`B60`$ T. Computations are carried out for spin polarized electronic systems, with total spin $`M_s^e=N_e/2`$, which contribute to the LHCP spectra. The energies of the incompletely polarized states with $`M_s^e=N_e/21`$, entering the RHCP spectra, are obtained by adding the corresponding Zeeman shifts.
### A Binding energies of excitonic complexes
We draw in Fig. 1 the g.s. energies, $`ϵ`$, coming from the diagonalization of $`V_{conf}+V_{coul}`$ in (2) as a function of the applied magnetic field. The polarized systems $`(N_e+1,N_h)`$= (1,1), (2,1), (3,1) and (4,1) are shown. The common notation for the excitonic systems (1,1) and (2,1) are $`X^0`$ and $`X^{}`$, so that the charged complex (4,1) may be denoted $`X^3`$. Note that the slopes of the (2,1), (3,1) and (4,1) curves are very similar. It means that the relative binding energies vary smoothly with $`B`$, and that the magnetic moments of these states take almost the same values. For example, $`X^3`$ is 14.77 meV above $`X^{}`$ at $`B=30`$ T, and 14.29 meV above $`X^{}`$ at $`B=50`$ T.
The total angular momenta in the g.s. is a constant, independent of $`B`$, in the smallest systems. It is $`M_{gs}=0`$ in the exciton, and $`M_{gs}=1`$ in the triplet $`X^{}`$ at any $`B`$. The larger systems, however, undergo abrupt rearrangements at particular $`B`$ values. The interplay between g.s. rearrangements in the $`(N_e,0)`$ and $`(N_e+1,1)`$ systems as $`B`$ is varied has direct consequences on absorption and luminescence, as will be seen below.
Note that, unlike pure electron systems, when holes are present the Hilbert space in a given $`M=|M_e|+M_h`$ sector is not finite. We enlarged the included subspace until convergence is reached. For example, in the (4,1) system at $`B=40`$ T, 2374 many-particle states (i.e. all of the states in $`15|M_e|35`$) are enough to reach convergence for the lowest energy eigenvalue in the $`M=15`$ tower.
The low-lying energy levels of $`X^3`$ at $`B=35`$ T are shown in Fig. 2 as an example. Energy distances between the lowest adjacent levels are around 0.5 meV, the same as in the three-electron system at this value of the magnetic field.
### B Interband absorption
As previously stated, temperatures are low enough for absorption to proceed from the g.s. of the $`N_e`$-electron system. It means that spin flips should not be thermally induced, i.e. $`T2`$ K for $`B>8`$ T.
We show in Fig. 3 the absorption coefficient for the $`N_e=0`$ qdot at $`B=40`$ T. The process under consideration, $`(0,0)(1,1)`$, goes through the $`M=0`$ channel. The main properties of the curve drawn in Fig. 3, i.e. dominance of the exciton g.s. and monotony, are visible also at any other value of the magnetic field. The main effect of $`B`$ is to reinforce the dominance of the first line. The threshold for absorption is determined by the exciton g.s. energy, and the maximum dipole squared behaves like $`B^{0.78}`$.
The absorption coefficient of the negatively charged dot, $`N_e=1`$, is shown in Fig. 4. The $`(1,0)(2,1)`$ process takes place in the $`M=0`$ sector. At $`B=8`$ T, a structure of isospaced bands is seen in the spectrum at higher energies. Most of these lines are suppressed already at $`B=40`$ T. The threshold for absorption and maximum strength transition are determined by the lowest $`X^{}`$ state in the $`M=0`$ tower. As a function of $`B`$, we get $`D^2B^{0.79}`$ at the maximum.
The absorption thresholds for the smallest systems, $`N_e=0`$ and 1, are smooth functions of $`B`$, signalling that the states entering the transition $`(N_e,0)(N_e+1,1)`$ do not change qualitatively as $`B`$ is raised. For larger systems, however, there is an abrupt decrease in the threshold for fields around 10 T (“filling factor” near one), and small steps at higher fields . The steps are originated by the different rates of change of $`M_{gs}`$ in the $`(N_e,0)`$ and $`(N_e+1,1)`$ systems (see Table I). Let us consider, for example, the $`(3,0)(4,1)`$ process. For $`B10`$ T, the process goes from the g.s. of (3,0) to the excited states of $`X^3`$ with $`M=3`$. For $`B>10`$ T, the g.s. of (3,0) moves to $`M=6`$, a sector which contains the g.s. of (4,1). Thus, the threshold is lowered. Every time one of the systems rearranges, there is a step like change in the absorption threshold. The actual (experimental) profile is expected to be smoothed because of temperature effects.
Of course, not only the threshold changes, but the whole spectrum is restructured. We show in Fig. 5 the absorption in the $`N_e=2`$ dot ($`X^2`$ formation) at $`B=8`$ T and 50 T. At $`B=8`$ T, the spectrum is similar to the $`X^{}`$ spectrum. The added electron is placed in an outer orbit because the inner orbitals are filled. For higher fields, there is place for the new electron in the core region, but the minimization of energy causes a global restructuration of the charge density in the dot, as will be seen below. The added pair losses its identity. Notice that for $`B>10`$ T there are two very distinct lines in the spectrum. One is the threshold (the transition to the lowest state of (3,1)), and the second is the maximum, which is 7-4 meV above the threshold.
The dipole squared at maxima as a function of $`B`$ are drawn in Fig. 6. Besides lowest state rearrangements, there are manifestations of collective effects even in these small systems. A decrease of absorption in the $`N_e=2`$ and 3 systems at “filling factors” $`\nu 1/2`$, 1/3 and 1/5 is evident from Fig. 6.
### C Magnetoluminescence
The second part of Fig. 5 shows the square of the dipole matrix elements corresponding to the luminescence of the $`N_e=2`$ dot at $`B=40`$ T. Only transitions starting from the g.s. of (3,1) are considered. Notice that the lowest state of (2,0) gives the strongest line, approximately 50 times higher than the next one. This is the common situation in our luminescence calculations for any of the systems under study. The strongest line corresponds to the transition from the g.s. of $`(N_e+1,1)`$ to the lowest state of $`(N_e,0)`$ in the same angular momentum tower. The higher states of $`(N_e,0)`$ give negligible contributions.
Luminescence in the $`N_e=0`$, and 1 dots is monotonic with $`B`$ because the initial and final states participating in it are fixed. Exciton luminescence proceeds in the $`M=0`$ channel, and $`X^{}`$ luminescence in the $`M=1`$ sector. In the latter case, the absorption and luminescence channels are different. With increasing $`B`$, the $`X^0`$ peak intensity increases, as in absorption, but the $`X^{}`$ intensity decreases. We obtained $`D^2`$ Exp$`(0.018B)`$ at the maximum.
For larger systems, the luminescence shows non monotonic behaviour because of lowest state rearrangements and collective effects, as in absorption. As a rule, the channels for absorption and PL are different in these systems. The luminescence maxima as a function of $`B`$ are drawn in Fig. 7.
### D Charge densities
Electron and hole charge densities inside the dot for the relevant states participating in absorption and luminescence are presented in this section. For electrons, we found more convenient to draw the difference $`\rho _e^{}=\rho _e(N_e+1,1)\rho _e(N_e,0)`$, which gives the density “added” to the dot.
Figure 8 shows the final-state densities in the absorption situations discussed in Fig. 5. For the $`N_e=2`$ dot at $`B=8`$ T, the added electron and hole densities are almost identical. The exciton keeps its identity inside the dot. At $`B=50`$ T, however, the added pair causes a redistribution of the charge density of the initial two-electron state.
On the other hand, as shown above, the relevant states participating in luminescence transitions are the g.s. of $`(N_e+1,1)`$ and the lowest state of the $`N_e`$-electron system in the same angular momentum sector. We show in Fig. 9 the densities of these states in the $`N_e=2`$ dot at $`B=40`$ T. These curves are typical. The exciton is annihilated from a distribution very similar to the isolated exciton g.s. (also shown in the figure for comparison).
## IV Concluding remarks
We have studied few-electron systems and excitonic complexes (with one hole) in qdots under intense magnetic fields and low temperatures. In 1- and 2-electron qdots the g.s. angular momentum is independent of the magnetic field intensity. However, larger systems undergo abrupt rearrangements at particular $`B`$ values, a fact that is reflected in the optical absorption and PL.
We computed the interband optical properties of these systems. In absorption, the initial state is the polarized ground state of $`N_e`$ electrons (for temperatures $`2`$ K), and the final states are the states of $`N_e+1`$ electrons and one hole. The main result of these computations is the non monotonic behaviour of the absorption maxima in the larger ($`N_e=2`$ and 3) systems as the field is varied (Fig. 6). This result can be understood as a consequence of ground state rearrangements and collective effects. We have presented typical charge densities in support of this picture. We found a reduction of absorption at “filling factors” 1/2, 1/3 and 1/5.
For luminescence events, we have considered the recombination from the g.s. of $`N_e+1`$ electrons and a hole. At a given magnetic field intensity, the angular momentum of this state may be different from the $`N_e`$-electron g.s. angular momentum. Thus, intrinsic absorption and luminescence may proceed through different channels. Of particular interest is that, opposite to the qwell case, the ground state of the negatively charged exciton $`X_t^{}`$ is bright in luminescence. This is a consequence of the qdot lateral confinement. Furthermore, for very high $`B`$ the $`X_t^{}`$ state recovers its dark character as compared with the other complexes here studied. On the other hand, the maximum of the recombination oscillator strength is a monotonic function of $`B`$ for qdots with 1 or 2 electrons and a hole, but it is nonmonotonic for qdots with more electrons, showing collective effects even in these small dots.
Although our calculations for finite systems with a smooth lateral confinement can not be easily extrapolated to the infinite limit, our results suggest that many-body effects should be taken into account in the computation of the $`X^{}`$ luminescence in a qwell. Whittaker and Shields , and Wojs et al have used a three-particle model for the $`X^{}`$. This model is indeed useful at very high magnetic fields. At intermediate values of $`B`$, the magnetoexciton size, which is $`2l_B50/\sqrt{B}`$ nm, becomes comparable to the inter-electronic distance, around 20 nm for a typical carrier density of 1-2$`\times 10^{11}`$ cm<sup>-2</sup>. Many-body effects should take care of the observed dependence of the PL maximum with the filling factor.
We have not attempted a more sophisticated calculation in these systems because of the absence of experimental results for qdots in very intense magnetic fields. Nevertheless, our simple approach (1LL, one qwell sub-band, parabolic lateral confinement, unrealistic Zeeman energies and $`z`$-averaged Coulomb interactions) captures the essential physics and indicates the importance of collective effects even in small qdots.
###### Acknowledgements.
A. G. acknowledges support by the Caribbean Network for Theoretical Physics. E. M-P acknowledges the Abdus Salam ICTP, where part of this work was done. The authors are grateful to C. Trallero-Giner for many useful discussions.
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# Hopf differentials and the images of harmonic maps
### §1 Results on non-surjectivity of harmonic diffeomorphisms
In \[HTTW\], it was proved that a polynomial growth harmonic diffeomorphism from $``$ into $`^2`$ is not surjective. In \[L-W 1\], the result was generalized to higher dimensions for polynomial growth harmonic maps between a more general class of manifolds. Not very many results are known if the map grows faster than polynomial. In this section, we will give results on non-surjectivity of certain harmonic diffeomorphisms from $``$ into $`^2`$ with fast growth rate. Note that the growth rate of a harmonic diffeomorphism from $``$ into $`^2`$ can be expressed in terms of the growth rate of its Hopf differential, see \[T-Wn 1\]. In particular, such a map is of polynomial growth if and only if its Hopf differential is of the form $`Pdz^2`$ with $`P`$ to be a polynomial.
###### Lemma 1.1
Let $`\mathrm{\Omega }`$ be a domain in $``$ which contains every disk $`𝔻(\sqrt{1}y,R(y))`$ with center $`\sqrt{1}y`$ and radius $`R(y)2\sqrt{2}(1+ϵ)\mathrm{log}y`$ for all $`yy_0>0`$, where $`ϵ>0`$ is a constant. Suppose $`h`$ is an orientation preserving harmonic diffeomorphism from $`\mathrm{\Omega }`$ into $`^2`$ with Hopf differential $`\mathrm{\Phi }=dz^2`$. Then the length of the image of the half line $`\mathrm{}zy_0`$, $`\mathrm{}z=0`$ under $`h`$ is bounded by a constant depending only on $`ϵ`$ and $`y_0`$.
###### Demonstration Proof
Let $`\mathrm{exp}(w)=h`$ be the norm of $`h`$ and let $`e`$ be the energy density of $`h`$ with respect to the Euclidean metric in the domain, then the pull-back metric under $`h`$ is given by
$$h^{}(ds_^2^2)=(e+2)dx^2+(e2)dy^2=2\left(\mathrm{cosh}(2w)+1\right)dx^2+2\left(\mathrm{cosh}(2w)1\right)dy^2.$$
$`1.1`$
As in \[Wf, My\] and page 63 in \[Hn\] we can prove that there is $`y_0>0`$ such that if $`yy_0`$
$$0<w(\sqrt{1}y)C_1\mathrm{exp}\left(\frac{R(y)}{2\sqrt{2}}\right)$$
$`1.2`$
where $`C_1`$ is an absolute constant. Hence the length $`\mathrm{}`$ of the image of $`\{\mathrm{}zy_0,\mathrm{}z=0\}`$ under $`h`$ satisfies:
$$\begin{array}{cc}\hfill \mathrm{}& =_{y_0}^{\mathrm{}}\left[2\left(\mathrm{cosh}(2w)1\right)\right]^{\frac{1}{2}}𝑑y\hfill \\ & C_3_{y_0}^{\mathrm{}}\mathrm{exp}\left(\frac{R(y)}{2\sqrt{2}}\right)𝑑y\hfill \\ & C_4_{y_0}^{\mathrm{}}y^{1ϵ}𝑑y\hfill \\ & =C_5\hfill \end{array}$$
where $`C_3C_5`$ are constants depending only on $`ϵ`$ and $`y_0`$, and we have used (1.1), (1.2) and the assumption that $`R(y)2\sqrt{2}(1+ϵ)\mathrm{log}y`$ if $`yy_0`$. The lemma then follows.
The following lemma basically says that if $`Q(z)=\frac{1}{2}z+o(1)`$ as $`\mathrm{}z\mathrm{}`$, then the behavior of $`\mathrm{exp}(Q(z))𝑑z`$ is similar to that of $`\mathrm{exp}(\frac{1}{2}z)𝑑z`$.
###### Lemma 1.2
Let $`2\pi A>0`$ and let $`Q(z)`$ be an analytic function on the half strip
$$𝒮=\{z|\mathrm{}z>\alpha >0\text{and}\theta A<\mathrm{}z<\theta +A\}$$
where $`\theta `$ is constant. Suppose $`Q(z)=\frac{1}{2}z+q(z)`$ such that $`|q(z)|g(\mathrm{}z)`$ where $`g(t)0`$ is a function defined on $`\mathrm{}>t\alpha `$ which satsifies $`lim_t\mathrm{}g(t)=0`$. Then for any $`\frac{1}{4}A>\delta >0`$ there exists $`𝔞>0`$ depending only on $`A`$, $`\delta `$, $`\alpha `$ and the function $`g`$ such that if $`z_0=x_0+\sqrt{1}\theta `$ with $`x_0>𝔞`$ and if
$$\zeta (z)=_{z_0}^z\mathrm{exp}\left(Q(\xi )\right)𝑑\xi ,$$
then $`\zeta `$ maps $`𝒮_\delta `$ injectively into $`\zeta `$-plane, and $`\zeta (𝒮_\delta )_{2\delta }\zeta (𝒮_{4\delta })`$. Here
$$𝒮_\delta =\{z𝒮|\mathrm{}z>x_0+\delta ,\theta A+\delta <\mathrm{}z<\theta +A\delta \},$$
$`𝒮_{4\delta }`$ is defined similarly and
$$\begin{array}{cc}\hfill _{2\delta }& =2\mathrm{exp}(\frac{1}{2}z_0)+\{\zeta ||\zeta |>2\mathrm{exp}\left(\frac{1}{2}(x_0+2\delta )\right),\hfill \\ & \frac{1}{2}(\theta A+2\delta )<\mathrm{arg}\zeta <\frac{1}{2}(\theta +A2\delta )\}.\hfill \end{array}$$
###### Demonstration Proof
Since $`lim_t\mathrm{}g(t)=0`$, for any $`ϵ>0`$, there is $`𝔞>\alpha `$ depending only on $`\alpha `$ and $`g`$ such that if $`\mathrm{}z>𝔞`$, then
$$\left|\mathrm{exp}\left(Q(z)\right)\mathrm{exp}(\frac{z}{2})\right|ϵ\mathrm{exp}(\frac{1}{2}\mathrm{}z).$$
$`1.3`$
Let $`x_0>𝔞`$ and let $`f(z)=2\left(\mathrm{exp}(\frac{1}{2}z)\mathrm{exp}(\frac{1}{2}z_0)\right)`$ with $`z_0=x_0+i\theta `$. Let $`z_1=x_1+\sqrt{1}y_1`$, $`z_2=x_2+\sqrt{1}y_2`$ in $`𝒮`$ such that $`x_1`$, $`x_2`$ are larger than $`x_0`$. Suppose $`x_1>x_2`$, then by (1.3)
$$\begin{array}{cc}& \left|\zeta (z_1)\zeta (z_2)f(z_1)+f(z_2)\right|\hfill \\ & =|z_1z_2||_0^1(\mathrm{exp}\left(Q(tz_1+(1t)z_2)\right)\mathrm{exp}(\frac{1}{2}(tz_1+(1t)z_2))dt|\hfill \\ & ϵ\left(|x_1x_2|+4\pi \right)\mathrm{exp}(\frac{1}{2}x_2)_0^1\mathrm{exp}(\frac{t}{2}(x_1x_2))𝑑t\hfill \\ & ϵ\mathrm{exp}(\frac{1}{2}x_2)\left[2\left(\mathrm{exp}(\frac{1}{2}(x_1x_2))1\right)+4\pi \mathrm{exp}\left(\frac{1}{2}(x_1x_2)\right)\right]\hfill \\ & C_1ϵ\mathrm{exp}(\frac{1}{2}x_1)\hfill \end{array}$$
where $`C_1`$ is an absolute constant. Obviously, the inequality is still true if $`x_1=x_2`$. Hence, we have
$$\left|\zeta (z_1)\zeta (z_2)f(z_1)+f(z_2)\right|C_1ϵ\mathrm{exp}(\frac{1}{2}\mathrm{max}\{\mathrm{}z_1,\mathrm{}z_2\})$$
$`1.4`$
provided $`\mathrm{}z_1,\mathrm{}z_2>x_0`$.
On the other hand, for any $`0<\delta _1<A`$ and $`z_1z_2𝒮`$, with $`\mathrm{}z_1\mathrm{}z_2`$ and $`|\mathrm{}(z_1z_2)|2A2\delta _1`$, we have
$$\begin{array}{cc}\hfill |f(z_1)f(z_2)|& =2\mathrm{exp}(\frac{1}{2}\mathrm{}z_1)|1\mathrm{exp}(\frac{1}{2}(z_2z_1))|\hfill \\ & \tau \mathrm{exp}(\frac{1}{2}\mathrm{}z_1)\hfill \end{array}$$
where $`\tau >0`$ depends only on $`A\delta _1`$ and the lower bound of $`|z_1z_2|`$ where we have used the fact that $`|\mathrm{}(z_1z_2)|2A2\delta _14\pi 2\delta _1`$. Hence for any $`z_1z_2𝒮`$,
$$|f(z_1)f(z_2)|\tau \mathrm{exp}(\frac{1}{2}\mathrm{max}\{\mathrm{}z_1,\mathrm{}z_2\})$$
$`1.5`$
where $`\tau >0`$ depending only on the lower bound of $`|z_1z_2|`$.
Let $`0<\delta _1<A`$, for any $`a`$ in $`𝒮_{\delta _1}\{z|\mathrm{}z<\beta \}`$ and $`z`$ the boundary of $`𝒮_{\frac{1}{2}\delta _1}\{z|\mathrm{}z<\beta +\frac{1}{2}\delta _1\}`$ where $`\beta `$ is a large number, we have
$$\begin{array}{cc}\hfill |f(z)f(a)|& \tau \mathrm{exp}(\frac{1}{2}\mathrm{}z)\hfill \\ & \frac{\tau }{C_1ϵ}|\zeta (z)f(z)|\hfill \end{array}$$
by (1.4) with $`z_1=z`$ and $`z_2=z_0`$, and (1.5), where $`\tau >0`$ is a constant depending only on $`\delta _1`$. Here we take $`z_1=z`$ and $`z_2=z_0`$ in (1.4). Choose $`ϵ`$ small enough depending only on $`A`$ and $`\delta _1`$ such that $`\frac{\tau }{C_1ϵ}>1`$, we have
$$|f(z)f(a)|>|\zeta (z)f(z)|.$$
Apply the Rouché Theorem to the functions $`\zeta f(a)`$, $`ff(a)`$ on $`𝒮_{\frac{1}{2}\delta _1}\{z|\mathrm{}z<\beta +\frac{1}{2}\delta _1\}`$ and then let $`\beta \mathrm{}`$ we conclude that for any $`a𝒮_{\delta _1}`$ there is one and only one $`zS_{\frac{1}{2}\delta _1}`$ such that $`\zeta (z)=f(a)`$.
On the other hand for such an $`a`$, we have
$$\begin{array}{cc}\hfill |\zeta (z)\zeta (a)|& |f(z)f(a)||\zeta (z)\zeta (a)f(z)+f(a)|\hfill \\ & \frac{1}{2}|f(z)f(a)|\hfill \\ & >|\zeta (z)f(z)|\hfill \end{array}$$
provided $`ϵ`$ is chosen to be small enough (depending only on $`A`$ and $`\delta _1`$). Hence there is also exactly one $`zS_{\frac{1}{2}\delta _1}`$ such that $`f(z)=\zeta (a)`$. From these the lemma follows by considering the image of $`f`$.
In the next theorem we will study the surjectivity of those harmonic diffeomorphisms from $``$ into $`^2`$ whose Hopf differentials are of finite order.
###### Theorem 1.1
Let $`h`$ be an orientation preserving harmonic diffeomorphism from $``$ into $`^2`$ with Hopf differential $`\mathrm{\Phi }=P\mathrm{exp}(Q)dz^2`$ such that
Then $`h`$ is not surjective. In particular, if $`P`$ is a polynomial then $`h`$ is not surjective.
###### Demonstration Proof
By the Hadamard factorization theorem, $`P(z)=z^me^{a(z)}A(z)`$, where $`m`$ is the order of $`z=0`$, $`a(z)`$ is a polynomial of degree less than $`\rho <n`$, and $`A(z)`$ is a canonical product of order less than or equal to $`\rho `$ formed by the zeros of $`P`$. So we can absorb $`a(z)`$ to the lower order terms of $`Q(z)`$ and assume $`P(z)`$ has the form $`z^mA(z)`$.
Let $`\zeta _1=z^n`$, which will map $`\delta <\mathrm{arg}z<\delta `$ bijectively onto $`n\delta <\mathrm{arg}\zeta _1<n\delta `$. In the region
$$_1=\{|\zeta _1|>R_0^n\}\{n\delta <\mathrm{arg}\zeta _1<n\delta \}.$$
$$\mathrm{\Phi }=n^2\zeta _1^{\frac{2(n1)}{n}}\stackrel{~}{P}(\zeta _1)\mathrm{exp}(\zeta _1)d\zeta _1^2$$
where $`\stackrel{~}{P}(\zeta _1)=P(z(\zeta _1))\mathrm{exp}(_{j=1}^na_j\zeta _1^{1j/n})`$. By (ii), without loss of generality we may assume that for $`|\zeta _1|>R_0^n`$, $`|\stackrel{~}{P}(\zeta _1)|\mathrm{exp}(|\zeta _1|^ϵ)`$ for some $`ϵ>0`$ which is small enough such that $`ϵ<1`$. By (iii) and lemma 2.6.18 in \[B\], we have for any $`\eta >0`$, $`\mathrm{log}|A(z)|>|z|^{\rho +\eta }`$ on $`\{z||z|>R_0\text{and}\delta <\mathrm{arg}z<\delta \}`$ for a possibly larger $`R_0`$ and a smaller $`\delta `$. Therefore, we have $`\left|\mathrm{log}|\stackrel{~}{P}(\zeta _1)|\right|=O(|\zeta _1|^ϵ)`$ as $`\zeta _1\mathrm{}`$ for some $`ϵ<1`$. From this, we have $`\left|\mathrm{log}|\stackrel{~}{P}(\zeta _1)|\right|=o(1)`$ on $`\stackrel{~}{}_1=_1\{n\delta +\delta _1<\mathrm{arg}\zeta _1<n\delta \delta _1\}`$ for any $`\delta _1>0`$; and hence $`|\frac{d}{d\zeta _1}\mathrm{log}\stackrel{~}{P}|=o(1)`$ as $`\zeta _1\mathrm{}`$ and $`\zeta _1\stackrel{~}{}_1`$. We conclude that in $`\stackrel{~}{}_1`$,
$$\mathrm{\Phi }=\mathrm{exp}(Q_1(\zeta _1))d\zeta _1^2=\varphi d\zeta _1^2$$
where
$$Q_1(\zeta _1)=\zeta _1+Q_2(\zeta _1)$$
with $`n\delta +\delta _1<\mathrm{arg}\zeta _1<n\delta \delta _1`$ and $`|\zeta _1|>R_0^n`$. Here $`Q_2(\zeta _1)=o(|\zeta _1|)`$ and $`\frac{d}{d\zeta _1}Q_2(\zeta _1)=o(1)`$ as $`\zeta _1\mathrm{}`$.
Let $`\zeta _2=\zeta _1+Q_2(\zeta _1)`$. This will map $`\{\pi <\mathrm{arg}\zeta _1<\pi \}\{|\zeta _1|>R_1\}`$ injectively onto its image for some $`R_1>R_0^n`$. Moreover, there exists $`R_2>0`$ such that $`_2=\{\pi +2\delta _1<\mathrm{arg}\zeta _2<\pi 2\delta _1\}\{|\zeta _2|>R_2\}`$ is in the image of the map $`\zeta _1\zeta _2`$. In $`_2`$, $`\mathrm{\Phi }`$ can be written in the form
$$\mathrm{\Phi }=\left[1+\frac{dQ_2}{d\zeta _1}\right]^2\mathrm{exp}(\zeta _2)d\zeta _2^2=\mathrm{exp}(Q_3(\zeta _2))d\zeta _2^2$$
where $`Q_3(\zeta _2)=\zeta _2+o(1)`$ as $`|\zeta _2|\mathrm{}`$. Let $`\zeta (\zeta _2)=^{\zeta _2}\mathrm{exp}(Q_3(\xi ))𝑑\xi `$. By Lemma 1.2, we conclude that $`\zeta (\zeta _2)`$ will map a subdomain of $`_2`$ bijectively onto the region
$$=\{\zeta ||\zeta |>R\text{and}\frac{1}{2}\pi \frac{1}{2}\delta _2<\mathrm{arg}\zeta <\frac{1}{2}\pi +\frac{1}{2}\delta _2\}$$
for some $`R>0`$ and $`\delta _2>0`$. On $``$, $`\mathrm{\Phi }=d\zeta ^2`$. The map $`\zeta \zeta _2\zeta _1z`$ is injective when restricted on $``$. Hence $`h(z(\zeta ))`$ is an orientation harmonic diffeomorphism from $``$ into $`^2`$. By Lemma 1.1, we conclude that the length of the image of the half line $`\mathrm{}\zeta >a_0`$, $`\mathrm{}\zeta =0`$ under $`h`$ is finite. Here $`a_0`$ is a large constant. By the definition of $`\zeta `$, $`\mathrm{}\zeta \mathrm{}`$ with $`\mathrm{}\zeta =0`$ implies that $`z\mathrm{}`$. Hence $`h`$ cannot be surjective.
Please see the appendix for figures showing the behaviour of the horizontal trajectories for some typical examples of holomorphic quadratic differentials discussed here (figures 1-6). If we refine the method of proof in Theorem 1.1, we can generalize the result to some cases that the Hopf differentials grow very fast (see figure 7 in appendix). First we have the following:
###### Lemma 1.3
Let $`R_0>0`$ and $`\delta >0`$ be constants and let $`h`$ be an orientaion preserving harmonic diffeomorphism from $`\mathrm{\Omega }_\delta =\{|z|>R_0,|\mathrm{arg}z|<\pi \delta \}`$ into $`^2`$ with Hopf differential of the form
$$\mathrm{\Phi }=\mathrm{exp}\left[g_1+\mathrm{exp}\left[g_2+\mathrm{}+\mathrm{exp}[g_k+Q]\mathrm{}\right]\right]dz^2$$
where $`Q(z)=z^n+_{j=1}^na_jz^{nj}`$ is a polynomial in $`z`$ and for each $`j=1,\mathrm{},k`$, $`|g_j(z)|=O(\mathrm{log}|z|)`$ as $`|z|\mathrm{}`$. Then there exists a path in $`\mathrm{\Omega }_\delta `$ diverging to infinity such that its image under $`h`$ has finite length.
###### Demonstration Proof
We will prove by induction on $`k`$. For $`k=1`$, $`\mathrm{\Phi }=\mathrm{exp}[g_1+Q]dz^2`$ and we can apply the same proof as in Theorem 1.1 to conclude the existence of such path.
For $`k2`$, we consider the map $`\zeta =f(z)=Q(z)+g_k(z)`$ on a convex subdomain $``$ in $`\mathrm{\Omega }_\delta `$ defined by
$$=\{z\mathrm{\Omega }_\delta |\mathrm{}zR_1,\frac{\pi }{2(n1)}+ϵ<\mathrm{arg}z<\frac{\pi }{2(n1)}ϵ\}$$
where $`R_1>R_0`$ and $`ϵ>0`$ will be chosen later. It is clear that for $`z=re^{i\theta }`$ with $`r`$ sufficiently large,
$$\mathrm{}f^{}(z)=nr^{n1}\mathrm{cos}[(n1)\theta ]+o(r^{n1})>0.$$
Therefore, if we choose $`R_1`$ sufficiently large, $`\zeta =f(z)`$ maps $``$ one-one onto its image $`f()`$ as $``$ is convex (proposition 1.10 in \[P\]). Since
$$f(re^{\pm i(\frac{\pi }{2(n1)}ϵ)})=r^ne^{\pm i(n\frac{\pi }{2(n1)}nϵ)}+o(r^n),$$
it is clear that $`f()`$ contains a subset of the form
$$\{|\mathrm{arg}\zeta |<n\pi /2(n1)ϵ_1,|\zeta |>R_2\}.$$
If we choose $`ϵ<\frac{\pi }{2n(n1)}`$, then we can choose accordingly an $`ϵ_1`$ such that $`n\frac{\pi }{2(n1)}ϵ_1>\pi /2`$. Hence, under this choice of $`R_1`$ and $`ϵ`$, $`f()`$ contains a half-plane $`\{\mathrm{}\zeta >R_2\}`$ for some $`R_2>0`$.
On $`f()`$, in particular on the half-plane $`\{\mathrm{}\zeta >R_2\}`$, $`\mathrm{\Phi }`$ can be written as
$$\mathrm{\Phi }=\mathrm{exp}\left[\stackrel{~}{g}_1(\zeta )+\mathrm{exp}\left[\stackrel{~}{g}_2(\zeta )+\mathrm{}+\mathrm{exp}[\stackrel{~}{g}_{k1}(\zeta )+\mathrm{exp}\zeta ]\mathrm{}\right]\right]d\zeta ^2,$$
where $`\stackrel{~}{g}_1(\zeta )=g_1(f^1(\zeta ))2\left(\frac{d}{dz}\mathrm{log}f\right)\left(f^1(\mathrm{log}\xi )\right)`$ and $`\stackrel{~}{g}_j(\zeta )=g_j(f^1(\zeta ))`$ for $`j=2,\mathrm{},k1`$.
Now, for any small $`\delta _1>0`$, lets consider the half strip
$$𝒮=\{\mathrm{}\zeta >R_2,|\mathrm{}\zeta |<\pi \delta _1\}.$$
The exponential map $`\xi =\mathrm{exp}\zeta `$ maps $`𝒮`$ one-one onto the domain
$$\mathrm{\Omega }_{\delta _1}=\{|\xi |>R_3,|\mathrm{arg}\xi |<\pi \delta _1\},$$
where $`R_3=e^{R_2}`$. And on $`\mathrm{\Omega }_{\delta _1}`$, $`\mathrm{\Phi }`$ takes the form
$$\mathrm{\Phi }=\mathrm{exp}\left[\stackrel{~}{\stackrel{~}{g}}_1(\xi )+\mathrm{exp}\left[\stackrel{~}{\stackrel{~}{g}}_2(\xi )+\mathrm{}+\mathrm{exp}[\stackrel{~}{\stackrel{~}{g}}_{k1}(\xi )+\xi ]\mathrm{}\right]\right]d\xi ^2,$$
where
$$\stackrel{~}{\stackrel{~}{g}}_1(\xi )=g_1(f^1(\mathrm{log}\xi ))2\left(\frac{d}{dz}\mathrm{log}f\right)\left(f^1(\mathrm{log}\xi )\right)2\mathrm{log}\xi $$
and $`\stackrel{~}{\stackrel{~}{g}}_j(\xi )=g_j(f^1(\mathrm{log}\xi ))`$ for $`j=2,\mathrm{},k1`$.
Since $`|g_j(z)|=O(\mathrm{log}|z|)`$ and $`\zeta =f(z)z^n`$ as $`|z|\mathrm{}`$, we see that
$$|g_j(f^1(\mathrm{log}\xi ))|=O(\mathrm{log}\mathrm{log}|\xi |)\text{for }j=1,\mathrm{},k1\text{.}$$
We also conclude from $`\frac{df}{dz}(z)z^{n1}\zeta ^{(n1)/n}`$ as $`|\zeta |\mathrm{}`$ that
$$|\left(\frac{d}{dz}\mathrm{log}f\right)\left(f^1(\mathrm{log}\xi )\right))|=O(\mathrm{log}\mathrm{log}|\xi |).$$
All together we have, as $`|\xi |\mathrm{}`$,
$$|\stackrel{~}{\stackrel{~}{g}}_j(\xi )|=\stackrel{~}{O}(\mathrm{log}|\xi |).$$
Therefore, induction hypothesis implies that there exists a divergent path $`\xi =\gamma (t)`$ in $`\mathrm{\Omega }_{\delta _1}`$ such that its image under the harmonic maps $`hf^1\mathrm{log}`$ has finite length. That is, there exists a path $`f^1(\mathrm{log}\gamma (t))`$ in $`\mathrm{\Omega }_\delta `$ such that its image under $`h`$ has finite length.
###### Theorem 1.2
Let $`h`$ be a harmonic diffeomorphism from $``$ into $`^2`$ with Hopf differential of the form
$$\mathrm{\Phi }=P_1\mathrm{exp}\left[P_2\mathrm{exp}\left[\mathrm{}\mathrm{exp}\left[P_k\mathrm{exp}(Q)\right]\mathrm{}\right]\right]dz^2,$$
where $`Q(z)=a_nz^n+\mathrm{}`$ and $`P_j`$, $`j=1,\mathrm{},k`$ are polynomials and $`k1`$. Then $`h`$ is not surjective.
###### Demonstration Proof
By making a change of parameter of the form $`zr_0\mathrm{exp}(\sqrt{1}\theta _0)z`$ for some constants $`r_0`$ and $`\theta _0`$, we may assume that $`Q=z^n+_{j=1}^na_jz^{nj}`$. As $`P_j`$, $`j=1,\mathrm{},k`$, are polynomials, there exists $`R_0>0`$ such that there is no zeros of any $`P_j`$ in the set $`\{|z|>R_0\}`$. Then for any $`\delta >0`$, one can define $`g_j=\mathrm{log}P_j`$ on the set $`\mathrm{\Omega }_\delta =\{|z|>R_0,|\mathrm{arg}z|<\pi \delta \}`$ and the Hopf differential $`\mathrm{\Phi }`$ can be written in the form required in Lemma 1.3. Therefore, there exists a path diverging to infinity in $`\mathrm{\Omega }_\delta `$ such that its image under $`h`$ has finite length. Hence, $`h`$ is not surjective.
Let us finish the dicussion of this section by given a different type of condition for the nonsurjectivity. Recall that a complex number $`a`$ in the extended complex plane is said to be an asymptotic value of an entire function $`f(z)`$ if there is a path $`z(t)`$, $`0t<1`$ such that $`lim_{t1}z(t)=\mathrm{}`$ and $`lim_{t1}f(z(t))=a`$. If $`a`$ is a finite number, then it is called a finite asymptotic value.
###### Theorem 1.3
Let $`f`$ be an entire function. Suppose there exist $`\delta >0`$ and $`R>0`$ such that
Suppose $`h`$ is an orientation harmonic diffeomorphism from $``$ into $`^2`$ with Hopf differential $`(f^{})^2dz^2`$, then $`h`$ is not surjective.
###### Demonstration Proof
Let $`\mathrm{\Omega }`$ be a component of $`f^1()`$. By (ii), $`f`$ is a local diffeomorphism on $`\mathrm{\Omega }`$. By (i), we can conclude that every path in $``$ begins at $`\zeta _0`$ can be lifted to a path in $`\mathrm{\Omega }`$ which begins at a point $`z_0`$ with $`f(z_0)=\zeta _0`$. Since $``$ is simply connected, $`f`$ maps $`\mathrm{\Omega }`$ bijectively to $``$. Hence $`hf^1(\zeta )`$ is a harmonic diffeomorphism from $``$ into $`^2`$ with Hopf differential $`d\zeta ^2`$. Moreover, $`f^1(\zeta )\mathrm{}`$ if $`\zeta `$ and $`\zeta \mathrm{}`$. The result follows from Lemma 1.1.
### §2 Maximal $`\mathrm{\Phi }`$-radius and quasi-conformal harmonic maps
Let us recall the definition of maximal $`\mathrm{\Phi }`$-radius of a holomorphic quadratic differential $`\mathrm{\Phi }`$ on a domain in $``$. Let $`\mathrm{\Omega }`$ be a domain in $``$ and let $`\mathrm{\Phi }=\varphi dz^2`$ be a holomorphic quadratic differential on $`\mathrm{\Omega }`$. Let $`z_0\mathrm{\Omega }`$ such that $`\varphi (z_0)0`$. Choose a branch of $`\sqrt{\varphi }`$ near $`z_0`$, and let
$$w=f(z)=_{z_0}^z\sqrt{\varphi (\zeta )}𝑑\zeta .$$
Let $`B(R)=\{w||w|<R\}`$ be the maximal disk in the $`w`$-plane such that $`f^1`$ is a conformal diffeomorphism from $`B(R)`$ into $`\mathrm{\Omega }`$. Then $`R_{z_0,\mathrm{\Omega }}=R`$ is called the maximal $`\mathrm{\Phi }`$-radius of $`\mathrm{\Phi }`$ at $`z_0`$ with respect to $`\mathrm{\Omega }`$ and $`V_{z_0,\mathrm{\Omega }}=f^1(B(R))`$ is called the maximal $`\mathrm{\Phi }`$ disk around $`z_0`$ with respect to $`\mathrm{\Omega }`$. We will drop the subscript $`\mathrm{\Omega }`$ if this will not cause any confusion. Moreover, by convention if $`\varphi (z_0)=0`$ we define $`R_{z_0}=0`$. In \[Hn\], it was proved that if $`h`$ is an orientation preserving harmonic diffeomorphism from a domain $`\mathrm{\Omega }`$ in $``$ into $`^2`$ with Hopf differential $`\mathrm{\Phi }`$, and if $`z_n`$ is a sequence in $`\mathrm{\Omega }`$ with $`R_{z_n}\mathrm{}`$, then the modulus the complex dilatation of $`h`$ at $`z_n`$ will tend to 1. Conversely, one would like to know whether $`h`$ would be quasi-conformal on a set with bounded maximal $`\mathrm{\Phi }`$-radius. In this section, we will prove that this is the case under certain assumptions. The result will be useful to study images of harmonic diffeomorphisms from $``$ into $`^2`$.
Let $`\mathrm{\Omega }`$ be a hyperbolic domain in $``$, i.e. its universal cover is conformal to the unit disk. Let $`\rho ^2|dz|^2`$ be the hyperbolic metric on $`\mathrm{\Omega }`$, i.e. the complete metric with constant Gaussian curvature $`1`$. Then it is known that \[Ah\] for any $`z\mathrm{\Omega }`$
$$\rho (z)\frac{2}{d(z,\mathrm{\Omega })},$$
where $`d(z,\mathrm{\Omega })`$ is the Euclidean distance from $`z`$ to $`\mathrm{\Omega }`$. If in addition, we have
$$\rho (z)\frac{C}{d(z,\mathrm{\Omega })}$$
for some positive constant $`C`$ for all $`z\mathrm{\Omega }`$, then we say that $`\mathrm{\Omega }`$ is strongly hyperbolic. Please note that our definition is slightly different from that in \[A-M-M\]. It is shown in Theorem 5 of \[A-M-M\] that if $`\mathrm{\Omega }`$ is bounded hyperbolic and the diameters of the boundary components are uniformly bounded from below by a positive constant then $`\mathrm{\Omega }`$ is strongly hyperbolic. Moreover, being strongly hyperbolic is conformally invariant:
###### Lemma 2.1
Let $`\mathrm{\Omega }_1`$ and $`\mathrm{\Omega }_2`$ be conformally equivalent domains. Suppose $`\mathrm{\Omega }_1`$ is strongly hyperbolic, so is $`\mathrm{\Omega }_2`$.
###### Demonstration Proof
Obviously $`\mathrm{\Omega }_2`$ is hyperbolic. Let $`w=f(z)`$ be a conformal diffeomorphism from $`\mathrm{\Omega }_1`$ onto $`\mathrm{\Omega }_2`$. Let $`\rho _2^2|dw|^2`$ be the hyperbolic metric on $`\mathrm{\Omega }_2`$ and let $`\rho _1^2|dz|^2`$ be the hyperbolic metric on $`\mathrm{\Omega }_1`$. Let $`d_1(z)=\text{dist}(z,\mathrm{\Omega }_1)`$ and $`d_2(w)=\text{dist}(w,\mathrm{\Omega }_2)`$, where both distances are Euclidean distances. Then by well-known fact \[V, p. 147\], we have
$$\frac{1}{4}d_1(z)|f^{}(z)|d_2(f(z)).$$
Hence
$$\begin{array}{cc}\hfill \rho _2(f(z))& =|f^{}(z)|^1\rho _1(z)\hfill \\ & \frac{C}{d_1(z)|f^{}(z)|}\hfill \\ & \frac{C}{4}d_2(f(z))\hfill \end{array}$$
for some positive contant $`C`$, where we have used the fact that $`\mathrm{\Omega }_1`$ is strongly hyperbolic. Hence $`\mathrm{\Omega }_2`$ is strongly hyperbolic.
Let $`\mathrm{\Omega }`$ be a strongly hyperbolic domain and let $`\mathrm{\Phi }=\varphi dz^2`$ be a holomorphic quadratic differential on $`\mathrm{\Omega }`$ with hyperbolic metric $`\rho ^2|dz|^2`$. For $`z\mathrm{\Omega }`$, let $`\mathrm{\Phi }(z)=\rho ^2(z)|\varphi |(z)`$ be the norm of $`\mathrm{\Phi }`$ at $`z`$ and let $`|\mathrm{\Phi }|=sup_{z\mathrm{\Omega }}\mathrm{\Phi }(z)`$. The following is proved in \[A-M-M\] (the equation (2) and the lemma 1.2)
###### Theorem
(Anić-Marković-Mateljević) With the above notations with $`\mathrm{\Omega }`$ being the unit disk $`𝔻`$, there exists an absolute constant $`C>0`$ such that for any holomorphic quadratic differential $`\mathrm{\Phi }`$ on $`𝔻`$ we have
$$\mathrm{\Phi }(z)C^1R_z^2$$
$`2.1`$
for all $`z𝔻`$, and
$$|\mathrm{\Phi }|CR_{\mathrm{}}^2$$
$`2.2`$
where $`R_{\mathrm{}}=sup_{z𝔻}R_z`$.
They actually proved that (2.1) is true for any hyperbolic domain in $``$. We will obtain a pointwise estimate for strongly hyperbolic domain which implies (2.2). The estimate will be useful in applications.
###### Proposition 2.1
Let $`\mathrm{\Omega }`$ be a hyperbolic domain and let $`\mathrm{\Phi }=\varphi dz^2`$ be a holomorphic quadratic differential defined on $`\mathrm{\Omega }`$. Then there exists an absolute positive constant $`C`$ such that for $`z\mathrm{\Omega }`$
$$R_zC\mathrm{\Phi }^{\frac{1}{2}}(z).$$
$`2.3`$
If in addition $`\mathrm{\Omega }`$ is strongly hyperbolic then there is a positive constant $`C^{}`$ depending only on $`\mathrm{\Omega }`$ such that for $`z\mathrm{\Omega }`$ with $`\mathrm{\Phi }(z)0`$
$$R_zC^{}\frac{\mathrm{\Phi }^2(z)}{|\mathrm{\Phi }|^{\frac{3}{2}}}$$
$`2.4`$
We need the following lemmas.
###### Lemma 2.2
Denote $`B(r)`$ to be the set of complex numbers with modulus less than $`r`$. Let $`f:B(r)`$ be an analytic function, such that $`f(0)=0`$ and $`f^{}(0)0`$. Suppose $`|f(z)|M`$ for all $`z`$, then (i) $`f`$ is one to one on $`B(r_1)`$, where $`r_1=\frac{r^2|f^{}(0)|}{8M}`$; (ii) $`f(B(r_1))B(\frac{r^2|f^{}(0)|^2}{32M})`$.
###### Demonstration Proof
Let us first assume that $`r=1`$, and $`f^{}(0)=1`$. Then
$$f(z)=z+\underset{n=2}{\overset{\mathrm{}}{}}a_nz^n.$$
By Cauchy theorem, we have $`|a_n|M`$, and $`1M`$. Suppose $`z_1z_2`$ are in $`B(\frac{1}{8M})`$, and $`r=\mathrm{max}\{|z_1|,|z_2|\}`$, then
$$\begin{array}{cc}\hfill |f(z_1)f(z_2)|& =\left|(z_1z_2)+\underset{n=2}{\overset{\mathrm{}}{}}a_n(z_1^nz_2^n)\right|\hfill \\ & |z_1z_2|\left|1M\underset{n=2}{\overset{\mathrm{}}{}}nr^{n1}\right|\hfill \\ & =|z_1z_2|\left|1Mr\frac{2r}{(1r)^2}\right|\hfill \\ & |z_1z_2|(18Mr)\hfill \\ & >0\hfill \end{array}$$
where we have used the facts that $`M1`$, and $`r<\frac{1}{8M}<\frac{1}{2}`$. Hence $`f`$ is one to one on $`B(\frac{1}{8M})`$. Using the fact that the Koebe’s constant is $`\frac{1}{4}`$ \[V, p. 149\], we have $`f\left(B(\frac{1}{8M})\right)`$ contains $`B(\frac{1}{32M})`$. In general, if $`f`$ is defined on $`B(r)`$ with $`f(0)=0`$ and with $`b=f^{}(0)0`$. Define $`\stackrel{~}{f}(\zeta )=\frac{f(r\zeta )}{rb}`$ for $`\zeta 𝔻`$. Then $`\stackrel{~}{f}(0)=0`$, and $`\stackrel{~}{f}^{}(0)=1`$. Let $`M_1=\frac{M}{r|b|}`$, then $`M_1|\stackrel{~}{f}(\zeta )|`$ for all $`\zeta `$. Hence $`\stackrel{~}{f}`$ is one to one on $`B(\frac{1}{8M_1})`$ and $`\stackrel{~}{f}\left(B(\frac{1}{8M_1})\right)`$ contains $`B(\frac{1}{32M_1})`$. Hence $`f`$ is one to one on $`B(\frac{r}{8M_1})=B(\frac{r^2|f^{}(0)|}{8M})=B(r_1)`$, and $`f\left(B(r_1)\right)`$ contains $`B(\frac{r|b|}{32M_1})=B(\frac{r^2|f^{}(0)|^2}{32M})`$.
The following lemma is proved in \[A-M-M, see Lemma 1.2\].
###### Lemma 2.3
Let $`f:B(r)`$ be analytic, and $`M|f(z)|`$ for all $`z`$. Suppose $`f(0)0`$, then $`f(z)0`$ for all $`zB(\frac{r|f(0)|}{2M}).`$
###### Demonstration Proof of Proposition 2.1
(2.3) was proved in the Lemma 2.3 of \[A-M-M\]. In order to prove (2.4), let $`z_0\mathrm{\Omega }`$, with $`\varphi (z_0)0`$. $`\varphi `$ is analytic on $`B(z_0,r_0)`$ where $`r_0=\frac{1}{2}d(z_0,\mathrm{\Omega })`$, where $`d`$ is the Euclidean distance. By Lemma 2.3, $`\varphi `$ is never zero on $`B(z_0,r)`$, where
$$r=\frac{r_0|\varphi (z_0)|}{2M_0}$$
and $`M_0=sup_{B(z_0,r_0)}|\varphi |`$. Hence we can take a branch of square root of $`\varphi `$ in $`B(z_0,r)`$. Let $`f(z)=_{z_0}^z\sqrt{\varphi }(\zeta )𝑑\zeta `$, for $`zB(z_0,r)`$, then $`f`$ is analytic, $`f(z_0)=0`$ and $`|f^{}(z_0)|=|\varphi (z_0)|^{\frac{1}{2}}0`$. By Lemma 2.2, $`f`$ is one to one on $`B(r_1)`$ where $`r_1=\frac{r^2|f^{}(z_0)|}{8M_1}`$, where $`M_1=sup_{B(z_0,r)}|f|`$. Moreover, $`f(B(r_1))`$ contains the disk $`B(R)=B(\frac{r^2|f^{}(z_0)|^2}{M_1})`$. Now $`M_1rM_0^{\frac{1}{2}}`$implies that
$$R_{z_0}R=\frac{r^2|\varphi (z_0)|}{32rM_0^{\frac{1}{2}}}=\frac{r_0|\varphi (z_0)|^2}{64M_0^{\frac{3}{2}}}.$$
$`2.5`$
This will imply the proposition because $`\mathrm{\Omega }`$ is strongly hyperbolic.
From the proof of the proposition, we have the following corollary which will be used in §3 and §4.
###### Corollary 2.1
Suppose $`\varphi `$ is analytic on $`B_{z_0}(R)`$ such that $`\alpha |\varphi |(z_0)|\varphi |(z)`$ for some constant $`\alpha >0`$ for all $`zB_{z_0}(R)`$. Let $`\mathrm{\Phi }=\varphi dz^2`$. Then the maximal $`\mathrm{\Phi }`$-radius of $`z_0`$ with respect to $`B_{z_0}(R)`$ is bounded below by $`\frac{R|\varphi |^{\frac{1}{2}}(z_0)}{64\alpha ^{\frac{3}{2}}}`$.
###### Demonstration Proof
This is a direct consequence of (2.5).
###### Lemma 2.4
Let $`\mathrm{\Omega }`$ be a simply connected domain in $``$ and $`\mathrm{\Phi }=\varphi dz^2`$ be a holomorphic quadratic differential on $`\mathrm{\Omega }`$. Let $`z_0\mathrm{\Omega }`$ such that $`\varphi (z_0)0`$ and let $`R`$ be the maximal $`\mathrm{\Phi }`$-radius of $`z_0`$ with maximal $`\mathrm{\Phi }`$-disk $`V`$. Suppose $`R<\mathrm{}`$ and suppose
$$w=\psi (z)=_{z_0}^z\sqrt{\varphi }𝑑\zeta $$
$`zV`$. Then for any $`0<\delta <R`$, there exists a point $`zV`$ with $`|\psi (z)|=\delta `$ such that the $`\mathrm{\Phi }`$-radius of $`z`$ is exactly $`R\delta `$.
###### Demonstration Proof
By the definitions of $`V`$ and $`R`$, $`\psi ^1:𝔻_RV`$ is a bijective conformal diffeomorphism, where
$$𝔻_R=\{w||w|<R\}.$$
It is easy to see that if $`zV`$ with $`|\psi (z)|=\delta `$ then the $`\mathrm{\Phi }`$-radius of $`z`$ is at least $`R\delta `$. Suppose the lemma is not true, then $`R_z>R\delta `$ for any $`z\psi ^1\left(\{w||w|=\delta \}\right)`$. Because $`R_z`$ is continuous, there is $`ϵ>0`$ such that $`R_zR\delta +ϵ`$ for all $`z`$ with $`|\psi (z)|=\delta `$. Hence $`\psi ^1`$ can be extended to an analytic function from $`𝔻_{R+ϵ}`$ to $`\mathrm{\Omega }`$ such that it is a local diffeomorphism. In particular, $`\varphi `$ is not zero in $`\psi ^1(𝔻_{R+ϵ})`$. By the definition of $`R`$, there exist two sequences $`w_n`$ and $`\stackrel{~}{w}_n`$ such that for each $`n`$ both $`w_n`$ and $`\stackrel{~}{w}_n`$ are in $`𝔻_{R+ϵ/n}`$, $`w_n\stackrel{~}{w}_n`$ but $`\psi ^1(w_n)=\psi ^1(\stackrel{~}{w}_n)`$. Without loss of generality, we may assume that $`w_na`$ and $`\stackrel{~}{w}_nb`$, and
$$\underset{n\mathrm{}}{lim}\psi ^1(w_n)=\underset{n\mathrm{}}{lim}\psi ^1(\stackrel{~}{w}_n)=c.$$
Since $`\psi ^1`$ is a local diffeomorphism, $`ab`$. Note that $`a`$ and $`b`$ are in $`\overline{𝔻}_R`$. Let $`\gamma `$ be the straight line joining $`a`$ and $`b`$ and let $`\mathrm{\Gamma }=\psi ^1(\gamma )`$. Then $`\mathrm{\Gamma }`$ is a smooth simple closed curve in $`\mathrm{\Omega }`$ because $`\psi ^1`$ is one to one on $`𝔻_R`$. Let $`\theta `$ be the interior angle at $`c`$. Apply the Gauss-Bonnet Theorem for the metric $`(|\varphi |+\eta )|dz|^2`$ on $`\mathrm{\Omega }`$ with $`\eta >0`$, we have
$$\frac{1}{2}_{\mathrm{\Omega }_1}\mathrm{\Delta }\mathrm{log}(|\varphi |+\eta )+_\mathrm{\Gamma }\kappa _\eta =\pi +\theta $$
where $`\mathrm{\Omega }_1`$ is the interior of $`\mathrm{\Gamma }`$, $`\kappa _\eta `$ is the geodesic curvature of $`\mathrm{\Gamma }`$ with respect to the metric $`(|\varphi |+\eta )|dz|^2`$ and $`\mathrm{\Delta }`$ is the Euclidean Laplacian. Here we have used the fact that $`\mathrm{\Omega }`$ is simply connected. Let $`c_1,\mathrm{},c_{\mathrm{}}`$ be the zeros of $`\varphi `$ inside $`\mathrm{\Omega }_1`$ with multiplicities $`k_1,\mathrm{},k_{\mathrm{}}`$ with $`k_j0`$. Let $`r>0`$ be small enough so that $`1j\mathrm{}`$ the disks $`D_j`$ of radius $`r`$ and centers at $`c_j`$ are disjoint and are inside of $`\mathrm{\Omega }_1`$. Then we have
$$\begin{array}{cc}\hfill \pi +\theta & =\frac{1}{2}\underset{j=1}{\overset{\mathrm{}}{}}_{D_j}\mathrm{\Delta }\mathrm{log}(|\varphi |+\eta )\frac{1}{2}_{\mathrm{\Omega }_1_{j=1}^{\mathrm{}}D_j}\mathrm{\Delta }\mathrm{log}(|\varphi |+\eta )+_\mathrm{\Gamma }\kappa _\eta \hfill \\ & =\frac{1}{2}\underset{j=1}{\overset{\mathrm{}}{}}_{D_j}\frac{}{r}\mathrm{log}(|\varphi |+\eta )\frac{1}{2}_{\mathrm{\Omega }_1_{j=1}^{\mathrm{}}D_j}\mathrm{\Delta }\mathrm{log}(|\varphi |+\eta )+_\mathrm{\Gamma }\kappa _\eta \hfill \\ & \frac{1}{2}\underset{j=1}{\overset{\mathrm{}}{}}_{D_j}\frac{}{r}\mathrm{log}|\varphi |+_\mathrm{\Gamma }\kappa \hfill \end{array}$$
as $`\eta 0`$, where $`\kappa `$ is the geodesic curvature with respect to the metric $`|\varphi ||dz|^2`$ and we have used the fact that $`\mathrm{log}|\varphi |`$ is harmonic. Since $`|\varphi ||dz|^2`$ is the pull-back metric under $`\psi ^1`$ of the flat metric in the $`𝔻_R`$, we have $`\kappa =0`$ on $`\mathrm{\Gamma }`$. Let $`r0`$, we conclude that
$$\pi +\theta =\underset{j=1}{\overset{\mathrm{}}{}}k_j\pi .$$
Since $`\theta 0`$ and $`k_j0`$ for all $`j`$, this is impossible.
###### Theorem 2.1
Let $`\mathrm{\Omega }`$ be a strongly hyperbolic domain in $``$ with hyperbolic metric $`e^{2v}|dz|^2`$ and $`\mathrm{\Omega }_1\mathrm{\Omega }`$. Let $`h`$ be an orientation preserving harmonic diffeomorphism from $``$ into $`^2`$ and let $`w=\mathrm{log}|h|`$, where $`|h|`$ is the norm of $`h`$ with respect to the Euclidean metric on $`\mathrm{\Omega }`$ and hyperbolic metric on $`^2`$. Let $`\mathrm{\Phi }=\varphi dz^2`$ be the Hopf differential of $`h`$ and let $`R_z`$ be the maximal $`\mathrm{\Phi }`$-radius of $`z`$ with respect to $``$. Suppose $`sup_{z\mathrm{\Omega }}R_z=R<\mathrm{}`$ and $`wvC`$ on $`\mathrm{\Omega }_1`$ for some constant $`C`$ and $`inf_{z\mathrm{\Omega }_1}R_z>0`$. Then $`h`$ is quasi-conformal on $`\mathrm{\Omega }_1`$.
###### Demonstration Proof
Suppose that $`h`$ is not quasi-conformal on $`\mathrm{\Omega }_1`$. Then there exists $`z_n\mathrm{\Omega }_1`$ such that $`\varphi (z_n)e^{2w(z_n)}1`$ as $`n\mathrm{}`$. Since $`R_{z_n}R`$, we may assume that $`lim_n\mathrm{}R_{z_n}=R_0`$. Suppose $`R_0>0`$. Let $`V_{z_n}`$ be the maximal $`\mathrm{\Phi }`$-disk with image $`𝔻_{R_{z_n}}`$ and let $`\zeta =_{z_n}^z\sqrt{\varphi }𝑑z=\psi _n(z)`$. Let $`\stackrel{~}{w}_n(\zeta )=(w\frac{1}{2}\mathrm{log}|\varphi |)(\psi _n^1(\zeta ))>0`$ which is considered as a function on $`𝔻_{R_{z_n}}`$. Then $`\stackrel{~}{w}_n>0`$ and
$$\mathrm{\Delta }_\zeta \stackrel{~}{w}_n=e^{2\stackrel{~}{w}_n}e^{2\stackrel{~}{w}_n}.$$
Then $`\stackrel{~}{w}_n`$ are locally uniformly bounded by proposition 1.5 in \[T-W 1\]. Passing to a subsequence if necessary, $`\stackrel{~}{w}_n`$ converges uniformly on compact subsets of $`𝔻_{R_0}`$. Since $`\stackrel{~}{w}_n(0)0`$, by mean value inequality, $`\stackrel{~}{w}_n0`$ uniformly on compact sets of $`𝔻_{R_0}`$. By Lemma 2.4, for each $`n`$ there exists $`\zeta _n`$ with $`|\zeta _n|=\frac{1}{2}R_{z_n}`$ such that the $`\mathrm{\Phi }`$-radius of $`z_n^{}=\psi _n^1(\zeta _n)`$ is $`\frac{1}{2}R_{z_n}`$. Moreover, we still have $`\varphi (z_n^{})e^{2w(z_n^{})}1`$. Continue in this way and by a diagonal process, if $`h`$ is not quasi-conformal on $`\mathrm{\Omega }_1`$, then we can find $`z_n`$ such that
$$\underset{n\mathrm{}}{lim}|\varphi (z_n)|e^{2w(z_n)}=1$$
$`2.6`$
and
$$\underset{n\mathrm{}}{lim}R_{z_n}=0.$$
$`2.7`$
Since $`inf_{z\mathrm{\Omega }_1}R_z>0`$, we may assume that $`z_n\mathrm{\Omega }_1`$ for all $`n`$. By Proposition 2.1, we have
$$\begin{array}{cc}\hfill R_{z_n}& C_3\frac{\mathrm{\Phi }^2(z_n)}{|\mathrm{\Phi }|^{\frac{3}{2}}}\hfill \\ & \frac{C_4}{R^3}\left[|\varphi (z_n)|e^{2v(z_n)}\right]^2\hfill \\ & \frac{C_5}{R^3}\left[|\varphi (z_n)|e^{2w(z_n)}\right]^2,\hfill \end{array}$$
$`2.8`$
where the norm of $`\mathrm{\Phi }`$ is taken with respect to the metric $`e^{2v}|dz|^2`$ on $`\mathrm{\Omega }`$. Here we have used (2.3), (2.4), the fact that the $`\mathrm{\Phi }`$-radius with respect to $`\mathrm{\Omega }_1`$ or $`\mathrm{\Omega }`$ is no greater than the $`\mathrm{\Phi }`$-radius with respect to $``$, the assumption that $`R_z`$ are uniformly bounded by $`R`$ on $`\mathrm{\Omega }`$ and that $`wvC`$ on $`\mathrm{\Omega }_1`$. Let $`n\mathrm{}`$ in (2.8), we have a contradiction because of (2.6) and (2.7). This completes the proof of the theorem.
###### Remark 2.1
In the theorem, we may replace $``$ by the unit disk. Moreover, suppose $`\mathrm{\Omega }`$ is a subset of $``$ (respectively $`^2`$) and $`h`$ is an orientation preserving harmonic diffeomorphism from $``$ (respectively $`^2`$) into $`^2`$ such that $`|h|^2|dz|^2`$ is complete in $``$ (respectively $`^2`$), where the norm is taken with respect to the Euclidean metric in the domain. Then the assumption that $`wvC`$ on $`\mathrm{\Omega }_1`$ in the theorem can be replaced by $`wvC`$ on $`\mathrm{\Omega }_1`$ by the comparison principle in \[Wn\].
By (2.3), which was proved in \[A-M-M\], and the above remark, we obtain a new proof of the following result in \[Wn\] as a corollary of Theorem 2.1.
###### Corollary 2.2
Let $`h`$ be an orientation preserving harmonic diffeomorphism on $`^2`$ with Hopf differential $`\mathrm{\Phi }`$ such that $`|h|^2|dz|^2`$ is complete. Suppose $`|\mathrm{\Phi }|<\mathrm{}`$, then $`h`$ is quasi-conformal. Here the norm of $`h`$ is taken with respect to the Euclidean metric in the domain and the norm of $`\mathrm{\Phi }`$ is taken with respect to the Poincaré metric while $`|\mathrm{\Phi }|`$ is taken with respect to the Poincaré metric.
### §3 Image of harmonic diffeomorphism with Hopf differential $`P\mathrm{exp}(Q)dz^2`$
Let $`h`$ be an orientation preserving harmonic diffeomorphism from $``$ into $`^2`$ with Hopf differential $`\mathrm{\Phi }=\varphi dz^2=P\mathrm{exp}(Q)dz^2`$ where $`P`$ and $`Q`$ are polynomials. By the result of §1, we know that $`h`$ is not surjective. Assume that $`|h|^2|dz|^2`$ is complete on $``$. In \[HTTW\], it was proved that if $`Q`$ is a constant, that is, if $`\varphi `$ is a polynomial of degree $`m`$, then the closure of the image of $`h`$ in $`\overline{^2}`$ is the convex hull of an ideal polygon with $`m+2`$ vertices in $`^2`$. The result is generalized from $``$ to surfaces with finite total curvature and in higher dimensions in \[L-W 1, 2\]. The assumption that $`\varphi `$ is a polynomial is equivalent to the fact that $`h`$ is of polynomial growth. Let $`A`$ be the intersection of the closure of the image of $`h`$ with the geometric boundary of $`^2`$. If $`Q`$ is not constant, then $`A`$ will no longer be a finite set. In this section, we will prove that in this case, $`A`$ is a countable set with exactly $`n`$ distinct accumulation points where $`n`$ is the degree of $`Q`$. In fact, we will prove that the result is true for a larger class of harmonic diffeomorphisms.
First we need a lemma. For $`\alpha 0`$, let
$$_\alpha =\{z|\mathrm{}z>\alpha ,\mathrm{}<\mathrm{}z<\mathrm{}\}.$$
Let $`h`$ be an orientation preserving harmonic diffeomorphism from $`_\alpha `$ into $`^2`$ with Hopf differential $`\mathrm{\Phi }=\mathrm{exp}(Q)dz^2`$ with $`Q(z)=z+q(z)`$ such that $`|q(z)|g(\mathrm{}z)`$ for some nonnegative function $`g`$ with $`lim_t\mathrm{}g(t)=0`$.
###### Lemma 3.1
With the above notations and assumptions, we have the following:
###### Demonstration Proof
For simplicity, assume $`\alpha =0`$. To prove the existence of those $`p_k^2`$ in (i), we apply Lemma 1.2 with $`\theta =0`$ and $`A=2\pi `$ to obtain that for any $`\pi >\delta >0`$, there exists $`x_0>0`$ such that if $`z_0=x_0`$ and $`\zeta (z)=_{z_0}^z\mathrm{exp}(\frac{1}{2}Q(\xi ))𝑑\xi +\mathrm{exp}(\frac{1}{2}x_0)`$, then $`\zeta `$ is injective on $`𝒮_{\frac{1}{4}\delta }`$, and $`\zeta (𝒮_{\frac{1}{4}\delta })_{\frac{1}{2}\delta }\zeta (𝒮_\delta )`$ where $`𝒮_\delta `$ and $`_\delta `$ etc. are defined as in Lemma 1.2 (with $`A=2\pi `$ and $`\theta =0`$). Then $`h(z(\zeta ))`$ is an orientation preserving harmonic diffeomorphism from $`\zeta (𝒮_{\frac{1}{4}\delta })`$ into $`^2`$ with Hopf differential $`\mathrm{\Phi }=d\zeta ^2`$. Note that the maximal $`\mathrm{\Phi }`$-radius of any point $`\zeta =u+\sqrt{1}v`$ in $`_{\frac{1}{2}\delta }`$ is at least $`uC_1`$ for some constant $`C_1`$ depending only on $`\delta `$. As in \[HTTW, p.109\], we can prove that the image of any horizontal half line $`\zeta (t)=t+\sqrt{1}v_0`$ with $`t`$ being larger than some constant in $`_{\frac{1}{2}\delta }`$ under $`h`$ is asymptotically a geodesic near infinity and tends to a point in $`^2`$ as $`t\mathrm{}`$. By the proof of Lemma 1.1, we can conclude that the image of any vertical line $`u=`$constant in $`_{\frac{1}{2}\delta }`$ under $`h`$ has uniformly bounded length. Hence if $`\zeta _n_{\frac{1}{2}\delta }`$, $`\zeta _n\mathrm{}`$ then $`h(z(\zeta _n))p_0`$ for some $`p_0^2`$. Since $`\zeta (𝒮_\delta )_{\frac{1}{2}\delta }`$ and $`z_n\mathrm{}`$ implies that $`\zeta (z_n)\mathrm{}`$ for $`z_n𝒮_\delta `$, we have
$$\underset{n\mathrm{}}{lim}h(z_n)=p_0.$$
Similarly, one can prove that for any integer $`k`$ there exists $`p_k^2`$ such that for any $`\delta >0`$ and $`z_n`$ with $`(2k1)\pi +\delta <\mathrm{}z_n<(2k+1)\pi \delta `$ such that $`z_n\mathrm{}`$, then
$$\underset{n\mathrm{}}{lim}h(z_n)=p_k.$$
To prove the remaining of (i) and (ii), we use Lemma 1.2 again to conclude that for all $`\delta >0`$ small enough, there exist $`a_j>0`$, $`b_j>0`$, $`j=1,2`$ such that for any integer $`k`$, there is a analytic function $`\zeta =\zeta ^{(k)}(z)`$ which maps
$$𝒮_1=\{z|\mathrm{}z>a_1,(2k1)\pi +\frac{1}{2}\delta <\mathrm{}z<(2k+3)\pi \frac{1}{2}\delta \}$$
and
$$𝒮_2=\{z|\mathrm{}z>a_2,2k\pi \delta <\mathrm{}z<2(k+1)\pi +\delta \}$$
injectively into $`\zeta `$-plane. Moreover, $`\zeta (𝒮_1)_1`$, $`\zeta (𝒮_2)_2`$, for $`j=1,2`$. Here
$$_1=\{\zeta ||\zeta |>b_1,\frac{1}{2}(2k1)\pi +\delta <\mathrm{arg}\zeta <\frac{1}{2}(2k+3)\pi \delta \},$$
$$_2=\left\{\zeta \right||\zeta |>b_2,\frac{1}{2}\left(2k\pi \frac{3}{2}\delta \right)<\mathrm{arg}\zeta <\frac{1}{2}\left(2(k+1)\pi +\frac{3}{2}\delta \right)\}$$
for $`j=1,2`$. Moreover, the Hopf differential $`\mathrm{\Phi }`$ of $`h`$ in the $`\zeta `$ coordinates is of the form $`d\zeta ^2`$. We will write $`h(\zeta )`$ instead of $`h(z(\zeta ))`$ if no confusion will arise. Let us consider the case when $`k`$ is even. The case that $`k`$ is odd is similar. By the previous result, we know that if $`\zeta _n_1`$ with $`\mathrm{}\zeta _n\mathrm{}`$ along a half line $`\mathrm{}\zeta `$=constant, then $`h(\zeta _n)p_k`$, and if $`\mathrm{}\zeta _n\mathrm{}`$, then $`h(\zeta _n)p_{k+1}`$.
In order to prove that $`p_kp_{k+1}`$ and that $`p_k`$ is monotone, we notice that the length of the curve $`h(z(\zeta ))`$ is infinite where $`\zeta =u+\sqrt{1}v_1`$ with $`v_1`$ to be a constant and $`\mathrm{}<u<\mathrm{}`$. Moreover, by \[Wf, M\] or \[HTTW, p.109\], the geodesic curvature of this curve is bounded by $`ϵ`$ provided $`v_1`$ is large. From this, it is easy to see that $`p_kp_{k+1}`$. Since $`h`$ is an orientation preserving diffeomorphism, we conclude that $`p_kp_j`$ if $`kj`$, and $`p_k`$ is monotone on $`𝕊^1`$. In particular, $`p_+=lim_k\mathrm{}p_k`$ and $`p_{}=lim_k\mathrm{}p_k`$ exist.
To prove (ii), we observe that for any $`C>0`$ there is $`v_0>0`$ independent of $`k`$ such that the $`\mathrm{\Phi }`$-radius of $`\zeta _1`$ is larger than $`C`$ for all $`\zeta `$ with $`\mathrm{}\zeta >v_0`$. By the argument in \[HTTW, p.102\], we conclude that for any $`ϵ>0`$, there is $`v_0>0`$ independent of $`k`$ such that if $`\mathrm{}\zeta >v_0`$, then $`d_^2(h(\zeta )),\gamma _k)ϵ`$, where $`\gamma _k`$ is the geodesic joining $`p_k`$ and $`p_{k+1}`$. From the proof of Lemma 1.2, we see that given $`v_0`$, there exists $`a>0`$ independent of $`k`$ such that if $`z𝒮_2`$ and $`\mathrm{}z>a`$, then $`\mathrm{}\zeta (z)>v_0`$. From this we can conclude that (ii) is true.
In order to prove (iii), let $`\delta >0`$ as above but small and let $`b=a_2`$ which is in the definition of $`𝒮_2`$. Suppose $`z_n_\alpha `$ with $`\mathrm{}z_n>b`$. Let $`k_n`$ be such that $`2k_n\pi \mathrm{}z_n<2(k_n+1)\pi `$. Then $`lim_n\mathrm{}k_n=\mathrm{}`$. For each $`n`$, let $`\zeta =\zeta ^{(k_n)}`$ as above then $`\zeta _n=\zeta (z_n)`$ can be defined and $`\zeta _n_2`$. By Lemma 1.1, for all $`\zeta _2`$ with $`\mathrm{}\zeta >0`$ and $`\mathrm{}\zeta =\mathrm{}\zeta _n`$, $`d_^2(h(\zeta _n),h(\zeta ))C_2`$ for some constant $`C_2`$ independent of $`n`$. From (ii), we conclude that $`d_^2(h(\zeta _n),\gamma _{k_n})C_3`$ for some constant $`C_3`$ independent of $`n`$. From this, the result follows.
###### Theorem 3.1
Let $`h`$ be an orientation preserving harmonic diffeomorphism from $``$ into $`^2`$ with Hopf differential $`\mathrm{\Phi }=\varphi dz^2=P\mathrm{exp}(Q)dz^2`$ such that $`|h|^2|dz|^2`$ is complete on $``$ and such that
Then the closure of the image of $`h`$ is the convex hull of a countable set $`A`$ of $`^2`$ with exactly $`n`$ accumulation points.
###### Demonstration Proof
We claim that for any $`ϵ>0`$ with $`nϵ<\frac{\pi }{2}`$, there exists a constant $`C_1>0`$ such that the maximal $`\mathrm{\Phi }`$-radius $`R_z`$ of $`z`$ satisfies
$$R_zC_1$$
$`3.1`$
for all $`zW_k`$, $`0kn1`$, where $`W_k`$ is the wedge
$$W_k=\left\{z\right|\left|\mathrm{arg}z\frac{(2k+1)\pi }{n}\right|<\frac{\pi }{2n}ϵ\}.$$
for $`0kn1`$. To prove the claim, note that there exists $`\tau >0`$ such that for $`zW_k`$, $`\mathrm{}(z^n)\tau |z|^n`$. By the assumptions (i) and (ii), for any $`zW_k`$, let $`\gamma `$ be the half ray $`\gamma (t)=t\mathrm{exp}(\sqrt{1}\mathrm{arg}z)`$ for $`t|z|`$, then
$$\begin{array}{cc}\hfill _{|z|}^{\mathrm{}}|\varphi |^{\frac{1}{2}}(\gamma (t))𝑑t& _{|z|}^{\mathrm{}}\mathrm{exp}\left(\frac{1}{2}\tau t^n+C_2(1+t^{n1}+t^{\stackrel{~}{\rho }})\right)𝑑t,\text{ for any }\rho <\stackrel{~}{\rho }<n,\hfill \\ & C_3\hfill \end{array}$$
where $`C_2`$ and $`C_3`$ are constants independent of $`z`$. Hence the maximal $`\mathrm{\Phi }`$-radius of $`zW_k`$ is uniformly bounded. This proves the claim.
Next, for each $`0kn1`$, and for $`\delta >4ϵ>0`$, let
$$V_{k,ϵ}=\left\{z\right|\left|\mathrm{arg}z\frac{2k\pi }{n}\right|<\frac{\pi }{2n}+ϵ\}.$$
Define $`V_{k,4ϵ}`$ similarly. By assumption (iii), we can take a branch of $`\mathrm{log}P`$ in $`\{zV_{k,4ϵ}||z|>R_0\}`$. As in the proof of Theorem 1.1, there exist positive constants $`R_2>R_1>R_0`$, $`T_2>T_1`$, $`ϵ_1>0`$ and a conformal map $`\zeta _k(z)`$ which is of polynomial growth as a function of $`z`$ and which will map $`𝒮_1^{(k)}=\{zV_{k,4ϵ}||z|>R_1\}`$ injectively onto its image. For simplicity, we write $`\zeta =\zeta _k`$. Moreover, if
$$𝒮_2^{(k)}=\{zV_{k,ϵ}||z|>R_2\}$$
$$_1=\{\zeta ||\mathrm{arg}\zeta |<\frac{\pi }{2}+2ϵ_1\text{and}|\zeta |>T_1\}$$
and
$$_2=\{\zeta ||\mathrm{arg}\zeta |<\frac{\pi }{2}+ϵ_1\text{and}|\zeta |>T_2\}$$
then $`\zeta (𝒮_1^{(k)})_1\zeta (𝒮_2^{(k)})_2`$. Moreover, in $`_1`$ the Hopf differential of $`h`$ is of the form $`\mathrm{\Phi }=\mathrm{exp}(\zeta +Q_1(\zeta ))d\zeta ^2`$ where $`Q_1(\zeta )0`$ as $`\zeta \mathrm{}`$. Choose $`a>b>T_2`$. As in the proof of (3.1), we have
$$R_zC_2$$
$`3.2`$
for some constant $`C_2`$ for all $`z𝒮_2^{(k)}\zeta ^1(\{\mathrm{}\zeta a\})`$. Moreover, on $`\mathrm{}\zeta =b`$, $`|\mathrm{exp}(\zeta +Q_1(\zeta ))|C_3`$ for some positive constant $`C_3`$. Hence if $`\stackrel{~}{w}=\mathrm{log}|_\zeta h|`$ and if $`e^{2\stackrel{~}{v}}|d\zeta |^2`$ is the hyperbolic metric on $`\zeta (𝒮^{(k)})\{\mathrm{}\zeta >a\}`$, then $`\stackrel{~}{w}\stackrel{~}{v}C_4`$ for some constant $`C_4`$ because $`e^{2\stackrel{~}{w}(\zeta )}|\mathrm{exp}(\zeta +Q_1(\zeta ))|<1`$ and $`\stackrel{~}{v}C`$ on $`\mathrm{}z=b`$ for some positive constant $`C`$. Let $`\mathrm{\Gamma }_k=\zeta ^1(\{\mathrm{}z=a\})`$ and $`\gamma _k=\zeta ^1(\{\mathrm{}\zeta =b\})`$. Note that for fixed $`c>T_2`$, $`\mathrm{arg}(\zeta ^1(c+\sqrt{1}t))(2k\pm \frac{1}{2})\pi /n`$ as $`t\pm \mathrm{}`$. Let $`\mathrm{\Omega }`$ be the component containing the origin of $`_{k=0}^{n1}\mathrm{\Gamma }_k`$, and let $`\mathrm{\Omega }_1`$ be the component containing the origin of $`_{k=0}^{n1}\gamma _k`$. By (3.1) and (3.2) if we choose $`ϵ>0`$ in (3.2) and then choose $`ϵ>0`$ in (3.1) small enough then we have $`R_zC_1+C_2`$ for all $`z\mathrm{\Omega }`$, and if $`e^{2v}|dz|^2`$ is the hyperbolic metric on $`\mathrm{\Omega }`$ then $`w=\mathrm{log}|_zh|vC`$ for some constant $`C`$ for all $`z\mathrm{\Omega }_1`$. Here we have used the fact that the hyperbolic metric on $`\mathrm{\Omega }`$ is dominated by the hyperbolic metric on its subdomain.
Next we want to show that $`inf_{z\mathrm{\Omega }_1}R_z>0`$. In fact, if $`z\mathrm{\Omega }_1`$, then there is $`k`$ such that $`\mathrm{}\zeta _k(z)b`$. Apply Corollary 2.1 on the disk with center $`\zeta _k`$ and radius 1, we can conclude that on $`\mathrm{}\zeta _kb`$ the maximal $`\mathrm{\Phi }`$-radius is bounded below by a positive constant independent of $`\zeta _k`$, because $`\mathrm{\Phi }=\mathrm{exp}(\zeta _k+o(1))d\zeta _k^2`$.
Since $`\mathrm{\Omega }`$ is strongly hyperbolic and $`|h|^2|dz|^2`$ is complete in $``$, $`h`$ is quasi-conformal on $`\mathrm{\Omega }_1`$ by Theorem 2.1 and Remark 2.1.
On the other hand, by Lemma 3.1, if we choose $`a`$ and $`b`$ large enough, then for each $`k`$, there exist $`p_j^{(k)}^2`$, $`j`$, which are monotone in $`𝕊^1`$ such that the intersection of the closure of the image under $`h`$ of the set $`\{\zeta _2,\mathrm{}\zeta b\}`$ with $`^2`$ is equal to
$$𝒜=\underset{k=0}{\overset{n1}{}}\{p_j^{(k)}|j\}\underset{k=0}{\overset{n1}{}}\{p_+^{(k)},p_{}^{(k)}\},$$
where $`p_\pm ^{(k)}=lim_{j\pm \mathrm{}}p_j^{(k)}`$. Moreover, if $`\zeta _n_2`$ with $`\mathrm{}\zeta _nb`$ and $`\mathrm{}z_n+\mathrm{}`$ (respectively $`\mathrm{}z_n\mathrm{}`$) then $`h(\zeta _n)p_+^{(k)}`$ (respectively $`h(\zeta _n)p_{}^{(k)}`$).
Since $`h`$ is at most linear growth in $`_2`$ with respect to $`\zeta `$, $`h`$ is of polynomial growth on $`V_{k,ϵ}`$, provided $`ϵ>0`$ is small enough. It is easy to see that $`h`$ is at most of linear growth on $`W_k`$. By the definition of $`\mathrm{\Omega }_1`$, we see that $`h`$ is of polynomial growth on $`\mathrm{\Omega }_1`$. Namely, there exist positive constants $`\mathrm{}`$ and $`C`$ such that
$$d_^2(h(z),o)C\left(d_{}(z,0)+1\right)^{\mathrm{}}$$
$`3.3`$
for all $`z\mathrm{\Omega }_1`$, where $`o`$ is a fixed point in $`^2`$ and $`0`$ is the origin of $``$. We claim that the image of $`h`$ is the convex hull of $`𝒜`$ together with at most finitely many points $`q_j^2`$. By theorem 4.8 in \[C-T\] and theorem 5 in \[Wn\], it is sufficient to show that $`\overline{h()}^2`$ is $`\{p_j^{(k)}|j\}\{p_+^{(k)},p_{}^{(k)}\}`$ together with at most finitely many points $`q_j`$. Suppose $`q_1,\mathrm{},q_m`$ are distinct points in
$$\left(\overline{h()}^2\right)𝒜.$$
There exist disjoint neighborhoods $`U_1,\mathrm{},U_m`$ of $`q_1,\mathrm{},q_m`$ respectively in $`\overline{}^2`$. We may choose $`U_j`$, $`1jm`$ small enough so that $`h^1(U_j)\mathrm{\Omega }_1`$. For if this is not true, then there exists $`q_j`$ and a sequence of neighborhoods $`U_{j,n}`$ such that $`_{n=1}^{\mathrm{}}U_{j,n}=\{q_j\}`$ and such that $`h^1(U_{j,n})`$ is not contained in $`\mathrm{\Omega }_1`$ for each $`n`$. By choosing a subsequence, we may assume that there is $`z_nU_{j,n}`$ such that $`\mathrm{}\zeta _k(z_n)b`$ under the map $`\zeta _k`$ described above. Since $`h(z_n)q_j`$ by construction, we conclude that $`q_j`$ must be $`p_l^{(k)}`$ or $`p_\pm ^k`$ for some $`k`$ and $`l`$. This is a contradiction. Hence we may choose $`U_j`$ such that $`h^1(U_j)`$ is contained in $`\mathrm{\Omega }_1`$. Moreover, we may assume that $`U_j`$ is bounded by a geodesic line in $`^2`$. Let $`f_j(z)=d_^2(u(z),^2U_j)`$, then $`f_j`$ harmonic because $`d_^2(,^2U_j)`$ is convex by \[B-O\]. Note that $`f_j`$ is smooth in $`h^1(U_j)`$, $`f_j(z)=0`$ for $`zh^1(U_j)`$ and there exists a constant $`C_3`$
$$f_j(z)C_3\left(d_{}(z)+1\right)^{\mathrm{}}$$
for all $`z`$ and for all $`1jm`$ by (3.3). Since $`h^1(U_j)`$, $`1jm`$, are disjoint and nonempty, $`m`$ is bounded from above by a constant depending only on $`\mathrm{}`$ by Theorem 3.4 in \[L-W 1\]. This proves the claim.
Observe that each $`q_j`$ must lie between $`p_+^{(k)}`$ and $`p_{}^{(k+1)}`$ for some $`k`$. Here we use the convention that $`p_+^{(n)}=p_{}^{(0)}`$. Since if $`\gamma _k(t)=\zeta ^1(b+\sqrt{1}t)`$, then $`lim_{t\pm \mathrm{}}h(\gamma _k(t))=p_{k,\pm }`$, we conclude that $`h(\mathrm{\Omega }_1)`$ is bounded by $`h(\gamma _k)`$ and the geodesics joining consecutive points of $`p_+^k`$, $`q_j`$ and $`p_{}^{(k+1)}`$, with $`q_j`$ between $`p_+^k`$ and $`p_{}^{(k+1)}`$, and they are oriented positively. Since $`h`$ is quasi-conformal on $`\mathrm{\Omega }_1`$, for each $`k`$ if $`\zeta _n_1`$ with $`\mathrm{}\zeta _nb`$ and $`\mathrm{}z_n+\mathrm{}`$ (respectively $`\mathrm{}z_n\mathrm{}`$) then $`h(\zeta _n)p_+^{(k)}`$ (respectively $`h(\zeta _n)p_{}^{(k)}`$). Again, using the fact that $`h`$ is quasi-conformal on $`W_k`$, we conclude that for $`zW_k`$, and if $`z\mathrm{}`$ then $`h(z)`$ will converge to a point $`q_k`$ in $`\overline{^2}`$. But $`q_k`$ must be equal to $`p_+^{(k)}`$ and $`p_{}^{(k+1)}`$ at the same time. Hence the closure of the image of $`h`$ is the convex hull of the set $`A`$ consisting of $`p_j^{(k)}`$, $`q_k`$ which is countable and has exactly $`n`$ accumulation points.
It is clear that the theorem is true for any polynomial $`Q`$ without requiring the leading coefficient to be $`1`$ as long as the zeros of the entire function $`P`$ are distributed in the corresponding sections. For instance, we conclude immediately from the theorem the following.
###### Corollary 3.1
Let $`h`$ be an orientation preserving harmonic diffeomorphism from $``$ into $`^2`$ with Hopf differential $`\mathrm{\Phi }=P\mathrm{exp}(Q)dz^2`$ where $`P`$ and $`Q`$ are polynomials with $`\mathrm{deg}Q=n`$. Suppose $`|h|^2|dz|^2`$ is complete in $``$. Then the image of $`h`$ is the convex hull of a countable set $`A`$ of $`^2`$ with exactly $`n`$ accumulation points.
Figures 1, 5, and 6 show the horizontal trajectories structures of holomorphic quadratic differentials which are included in the corollary 3.1. Figure 1 also shows the image of the harmonic map corresponding to $`e^zdz^2`$ which is the basis of all the discussion in this paper.
### §4 Images of harmonic diffeomorphisms with Hopf differential $`f(e^z)dz^2`$
In this section, we will study the images of certain harmonic diffeomorphisms from $``$ into $`^2`$ with Hopf differentials of the form $`f(e^z)dz^2`$. As before for $`\alpha 0`$, let $`_\alpha =\{z=x+\sqrt{1}y|x>\alpha \}`$.
###### Lemma 4.1
Let $`h:_0^2`$ be an orientation preserving harmonic diffeomorphic injection with Hopf differential $`\mathrm{\Phi }=dz^2`$. Suppose that for some $`x_0>0`$, $`lim_{y+\mathrm{}}h(x_0+\sqrt{1}y)=p_1`$ and $`lim_y\mathrm{}h(x_0+\sqrt{1}y)=p_2`$ for some $`p_1`$, $`p_2`$ in $`^2`$. Then $`p_1=p_2=p`$, and for all $`x_0>0`$
$$\underset{\genfrac{}{}{0pt}{}{|z|\mathrm{}}{\mathrm{}zx_0}}{lim}h(z)=p.$$
###### Demonstration Proof
By proposition 1.5 in \[ T-W 1\], see also \[Wn\], the energy density of $`h`$ in the half-plane $`0<\mathrm{}z<\mathrm{}`$ is bounded. Since for some $`x_0>0`$, $`lim_{y+\mathrm{}}h(x_0+\sqrt{1}y)=p_1`$ and $`lim_y\mathrm{}h(x_0+\sqrt{1}y)=p_2`$, where $`p_1,p_2^2`$, we conclude that for all $`x_1>x_0>0`$,
$$\underset{\genfrac{}{}{0pt}{}{\mathrm{}z+\mathrm{}}{x_0<\mathrm{}z<x_1}}{lim}h(z)=p_1$$
and
$$\underset{\genfrac{}{}{0pt}{}{\mathrm{}z\mathrm{}}{x_0<\mathrm{}z<x_1}}{lim}h(z)=p_2.$$
Identify $`\overline{^2}=^2^2`$ with the unit disk. We claim that for any $`x_0>0`$, the closure of $`h(_{x_0})`$ in $`\overline{^2}`$ is $`\overline{\mathrm{\Omega }}`$ where $`\mathrm{\Omega }`$ is the domain bounded by the curve $`h(x_0+\sqrt{1}y)`$, $`\mathrm{}<y<\mathrm{}`$ and one of the arc on $`𝕊^1`$ with end points $`p_1`$ and $`p_2`$. Obviously, $`h(_{x_0})`$ is contained in such an $`\mathrm{\Omega }`$ because $`h`$ is injective. Suppose that the claim is not true, then there is $`q`$ on the boundary of $`h(_{x_0})`$ such that $`q^2`$ and there is a geodesic arc $`\gamma `$ in $`h(_{x_0})^2`$ from a point $`q_1`$ in $`h(_{x_0})`$ to $`q`$ with $`\gamma (\mathrm{})=q`$, where $`\mathrm{}`$ is the length of $`\gamma `$ and is finite. Without loss of generality, we may asumme that $`\gamma ([0,\mathrm{}))h(_{x_0})`$. Let $`\beta =h^1(\gamma )`$. Then $`\beta `$ is a path in $`_{x_0}`$ such that $`\beta (t)\mathrm{}`$ as $`t\mathrm{}`$ because $`q`$ is in the boundary of $`h(_{x_0})`$. Moreover, $`\mathrm{}\beta (t)+\mathrm{}`$. Otherwise, we would have $`\beta (t)p_1`$ or $`p_2`$. However, the pull-back metric under $`h`$ is given by $`(e+2)dx^2+(e2)dy^2`$, where $`e`$ is the energy density of $`h`$, and $`e>2`$. We then have
$$\begin{array}{cc}\hfill \mathrm{}& =_0^{\mathrm{}}\left[(e+2)\left(\frac{dx}{dt}\right)^2+(e2)\left(\frac{dy}{dt}\right)^2\right]^{\frac{1}{2}}𝑑t\hfill \\ & \sqrt{2}\left(\underset{t\mathrm{}}{lim}x(\mathrm{})x(0)\right)\hfill \\ & =\mathrm{}\hfill \end{array}$$
which is a contradiction. Hence $`h(_{x_0})=\mathrm{\Omega }`$. Suppose $`p_1p_2`$, then $`\overline{h(_{x_0})}`$ contains a nontrivial arc on $`𝕊^1`$. However, for $`x_0>0`$, $`h`$ is of at most linear growth. By theorem 3.4 in \[L-W 1\], we conclude that $`\overline{h(_{x_0})}^2`$ consists of only finitely many points. Hence we must have $`p_1=p_2=p`$. Since $`h`$ is a diffeomorphism, we must have
$$\underset{\genfrac{}{}{0pt}{}{|z|\mathrm{}}{\mathrm{}zx_0}}{lim}h(z)=p.$$
###### Lemma 4.2
Let $`0<\beta \pi `$ and let $`h:e^{\sqrt{1}\beta }_0^2`$ be an orientation preserving harmonic diffeomorphic injection with Hopf differential $`\mathrm{\Phi }=dz^2`$. Suppose that for some $`x_0>0`$,
$$\underset{y\mathrm{}}{lim}h\left(e^{\sqrt{1}\beta }(x_0+\sqrt{1}y)\right)=p_1\text{and}\underset{y\mathrm{}}{lim}h\left(e^{\sqrt{1}\beta }(x_0+\sqrt{1}y)\right)=p_2$$
for some $`p_1`$, $`p_2`$ in $`^2`$. Then $`p_1p_2`$ and for all $`x_0>0`$, $`\overline{h(e^{\sqrt{1}\beta }_{x_0})}^2=\{p_1,p_2\}`$.
###### Demonstration Proof
As in the proof of Lemma 4.1, we conclude that for any $`x_0>0`$,
$$\underset{y\mathrm{}}{lim}h\left(e^{\sqrt{1}\beta }(x_0+\sqrt{1}y)\right)=p_1,$$
and
$$\underset{y\mathrm{}}{lim}h\left(e^{\sqrt{1}\beta }(x_0+\sqrt{1}y)\right)=p_2.$$
Let $`x_0>0`$. Since $`0<\beta \pi `$, by Lemma 1.1, suppose $`z_ne^{\sqrt{1}\beta }_{x_0}`$, if $`\mathrm{}z_n\mathrm{}`$, then $`lim_n\mathrm{}h(z_n)=p_1`$; and if $`\mathrm{}z_n\mathrm{}`$, then $`lim_n\mathrm{}h(z_n)=p_2`$. If $`z_n\mathrm{}`$ and for all $`n`$, $`x_0\mathrm{}z_n<x_1`$ for some $`x_1`$, then $`h(z_n)`$ are uniformly bounded. Hence $`\overline{h(e^{\sqrt{1}\beta }_{x_0})}_{\mathrm{}}^2=\{p_1,p_2\}`$.
To prove that $`p_1p_2`$. Note that $`\sqrt{1}ne^{\sqrt{1}\beta }`$ for any positive integer $`n`$. Moreover, it is easy to see that $`z_n=e^{\sqrt{1}\frac{\beta }{3}}\sqrt{1}n=n\mathrm{sin}\frac{\beta }{3}+\sqrt{1}n\mathrm{cos}\frac{\beta }{3}`$ and $`\stackrel{~}{z}_n=n+\sqrt{1}n\mathrm{cos}\frac{\beta }{3}`$ are in $`e^{\sqrt{1}\beta }C_{x_0}`$ if $`n`$ is large. Let $`L_n`$ be the horizontal line joining $`z_n`$ and $`\stackrel{~}{z}_n`$. By the arguments in section 3 of \[HTTW\], we conclude that $`h(L_n)`$ is of uniformly bounded distance from the geodesic passing through $`h(z_n)`$ and $`h(\stackrel{~}{z}_n)`$. Since $`h(z_n)p_1`$ and $`h(\stackrel{~}{z}_n)p_2`$, if $`p_1=p_2=p`$ then $`h(L_n)p`$ as $`n\mathrm{}`$. On the other hand, $`h(\sqrt{1}n\mathrm{cos}\frac{\beta }{3})`$ are uniformly bounded. This is a contradiction. Therefore, $`p_1p_2`$.
###### Theorem 4.1
Let $`m`$, $`n`$ be nonnegative intergers and let $`P(t)`$ be a nonconstant rational function of the form
$$P(t)=\underset{k=m}{\overset{n}{}}a_kt^k,$$
with $`a_m0a_n`$. Suppose $`h`$ is an orientation preserving harmonic diffeomorphism from $``$ into $`^2`$ with Hopf differential given by
$$\mathrm{\Phi }=P(e^z)dz^2$$
such that $`|h|^2|dz|^2`$ is a complete metric. Then $`𝒜=\overline{h()}^2`$ is countable which has exactly one accumulation point if $`m`$ or $`n=0`$, and $`a_00`$; and has two accumulation points otherwise. Moreover $`\overline{h()}`$ is the convex hull of $`𝒜`$.
###### Remark
Figures 1 and 2 in the appendix show horizontal trajectories structures for the case that $`m`$ or $`n=0`$, and $`a_00`$. In fact, in both figures, $`m=0`$, and $`a_0=0`$ and $`1`$ respectively. The other case are showed by the figures 3 and 4. In figure 3, $`m=0`$ but $`a_0=1`$. The image of the corresponding harmonic map has 2 accumulations both are limits from 1 side. In figure 4, both $`m`$ and $`n`$ are not zero and the image of the corresponding harmonic map has 2 accumulations both are limits from 2 sides.
###### Demonstration Proof
Suppose that $`m>0`$ and $`n>0`$. By the proof of Lemma 3.1, we can conclude that there exist $`p_k`$, $`k`$ such that $`\overline{h(\{z|\mathrm{}z0\}}^2`$ is equal to $`\overline{\{p_k\}_k}`$ and the $`p_k`$ are monotone on $`𝕊^1`$. Moreover, if $`p_kp_\pm `$ as $`k\pm \mathrm{}`$, then $`lim_{\mathrm{}z+\mathrm{},\mathrm{}z0}h(z)=p_+`$ and $`lim_{\mathrm{}z\mathrm{},\mathrm{}z0}h(z)=p_{}`$.
Similarly, there exist $`q_k`$, $`k`$ such that $`\overline{h(\{z|\mathrm{}z0\}}^2`$ is equal to $`\overline{\{q_k\}_k}`$ and the $`q_k`$ are monotone on $`𝕊^1`$, and if $`q_kq_\pm `$ as $`k\pm \mathrm{}`$, then $`lim_{\mathrm{}z+\mathrm{},\mathrm{}z0}h(z)=q_+`$ and $`lim_{\mathrm{}z\mathrm{},\mathrm{}z0}h(z)=q_{}`$. Hence $`q_+=p_+`$ and $`q_{}=p_{}`$. Since $`h`$ is a diffeomorphism, $`p_+p_{}`$ and $`𝒜=\{p_k,q_k\}_k\{p_+,p_{}\}`$, which has two accumulation points.
Next, let us consider the case that $`m`$ or $`n=0`$. Without loss of generality, we may assume that $`m=0`$. As before, there exist $`p_k`$, $`k`$ such that $`\overline{h(\{z|\mathrm{}z0\}}^2`$ is equal to $`\overline{\{p_k\}_k}`$ and the $`p_k`$ are monotone on $`𝕊^1`$. Let $`p_+`$ and $`p_{}`$ defined as above.
Suppose $`a_0=0`$. Then we can conclude as in the proof of Theorem 3.1 that $`p_+=p_{}=p`$ and $`𝒜=\{p_k\}_k\{p\}`$ which has only one accumulation points.
Suppose $`m=0`$ and $`a_00`$, let $`a_0=\rho ^2e^{2\sqrt{1}\beta }`$ with $`0\beta <\pi `$, $`\rho >0`$. There exists $`\delta >0`$ such that if $`|t|<\delta `$, we can take a branch of the square root of $`P(t)`$ and
$$\sqrt{P(t)}=\rho e^{\sqrt{1}\beta }+tg(t),$$
where $`g(t)`$ is analytic and
$$|g(t)|C_1$$
$`4.1`$
for some constant $`C_1`$ for $`|t|<\delta `$. Let $`\stackrel{~}{g}(t)`$ be such that $`\stackrel{~}{g}^{}=g`$ on $`|t|<\delta `$ and $`\stackrel{~}{g}(0)=0`$. Let $`x_0<0`$ be small enough so that $`|e^z|<\delta `$ on $`\mathrm{}zx_0`$. Define
$$\begin{array}{cc}\hfill \zeta (z)& =_{x_0}^z\sqrt{P(e^\xi )}𝑑\xi \hfill \\ & =\rho e^{\sqrt{1}\beta }(zx_0)+_{x_0}^ze^\xi g(e^\xi )𝑑\xi \hfill \\ & =\rho e^{\sqrt{1}\beta }z+\stackrel{~}{g}(e^z)+\zeta _0\hfill \end{array}$$
for all $`z`$ with $`\mathrm{}zx_0`$, where $`\zeta _0`$ is a constant. Here the integration is along the straight line from $`x_0`$ to $`z`$. Then $`\zeta `$ is analytic. By (4.1), if we choose $`x_0`$ small enough, then $`\zeta `$ is injective. Since $`|\stackrel{~}{g}(e^z)|C_2|e^z|`$ for some constant $`C_2`$, if we choose $`x_0`$ small enough, then the analytic map $`z\zeta _1=(\zeta \zeta _0)`$ will map $`\{\mathrm{}zx_0\}`$ injectively onto its image $``$. Moreover
$$e^{\sqrt{1}\beta }\{\zeta _1|\mathrm{}z_1\rho x_0+1\}e^{\sqrt{1}\beta }\{\zeta _1|\mathrm{}z_1\rho x_01\}.$$
The Hopf differential of $`h`$ in the $`\zeta _1`$ plane is given by $`d\zeta _1^2`$. As before, we have
$$\underset{y\pm \mathrm{}}{lim}h(x_0+\sqrt{1}y)=p_\pm .$$
Hence if $`\beta =0`$, we have $`p_+=p_{}=p`$ by Lemma 4.1 and $`𝒜`$ is countable with only one accumulation point. If $`\beta >0`$, we have $`p_+p_{}`$ by Lemma 4.2 and $`𝒜`$ is countable with exactly two accumulation points. The last statement of the theorem follows from theorem 4.8 in \[C-T\] and theorem 5 in \[Wn\].
As an application, we use Theorem 4.1 to study harmonic diffeomorphic injection from a flat cylinder to a hyperbolic cylinder. Let $`N`$ be a hyperbolic cylinder and Let $`^{}=\{0\}`$. Let $`\mathrm{\Phi }(^{},N)`$ be the set of all Hopf differentials of orientation preserving harmonic diffeomorphic injections $`h`$ from $`^{}`$ to $`N`$ such that $`|h|^2|dz|^2`$ is complete on $`^{}`$. Let $`𝒫(N)`$ be the set of holomorphic quadratic differentials on $`^{}`$ defined by
$$𝒫(N)=\left\{\frac{P(z)}{z^2}dz^2\right|P(z)=\underset{k=m}{\overset{n}{}}a_kz^k\text{for some }0m,n\text{, and }Pa_0\}.$$
$$𝒫_1(N)=\left\{\frac{P(z)}{z^2}dz^2𝒫(N)\right|P(z)=\underset{k=m}{\overset{n}{}}a_kz^k\text{with }m\text{ or }n=0\text{, and }a_00\}$$
and $`𝒫_2(N)=𝒫(N)𝒫_1(N)`$.
###### Corollary 4.1
With the above notations we have $`\mathrm{\Phi }(^{},N)𝒫(N)`$ is either a subset of $`𝒫_1(N)`$ or a subset of $`𝒫_2(N)`$. Moreover, if $`\mathrm{\Phi }(^{},N)𝒫(N)\mathrm{}`$, then it is a subset of $`𝒫_1(N)`$ if and only if $`N`$ has a cusp.
###### Demonstration Proof
Let $`z^2P(z)dz^2\mathrm{\Phi }(C^{},N)`$ be the Hopf differential of an orientation preserving harmonic diffeomorphic injections $`h`$ from $`^{}`$ into $`N`$. Lifting $`h`$ to the universal coverings, we have an orientation preserving harmonic diffeomorphic injection, denoted by $`h`$ again, from $``$ into $`^2`$, with Hopf differential given by
$$P(e^z)dz^2$$
and an element $`\rho `$ of the Möbius group which generates $`\pi _1(N)`$ such that
$$h(z+2\pi i)=\rho (h(z)).$$
Note that $`|h|^2|dz|^2`$ is complete on $``$. Let $`𝒜=\overline{h()}_{\mathrm{}}^2`$. Since $`h`$ is equivariant, $`𝒜`$ is invariant under $`\rho `$. This implies that the set of fixed points of $`\rho `$ is exactly the set of accumulation points of $`A`$. The corollary then follows easily from Theorem 4.1.
###### Remark 4.1
It was proved in \[Wn, W-A, T-W 1\] that given a holomorphic quadratic differential $`\mathrm{\Phi }`$ on $``$ or on $`^2`$ there exists an orientation preserving harmonic diffeomorphic injection from $``$ or $`^2`$ to $`^2`$ whose Hopf differential is the given $`\mathrm{\Phi }`$. Corollary 4.1 shows that the prescribed Hopf differential problem is not alway solvable from $`^{}`$ into $`N`$ where $`N`$ is a hyperbolic cylinder.
Our next result is to consider the image of a harmonic map with Hopf differential with infinite order.
###### Theorem 4.2
Let $`h`$ be an orientation preserving harmonic diffeomorphic injection from $``$ into $`^2`$ such that $`|h|^2|dz|^2`$ is complete. Suppose that the Hopf differential of $`h`$ is given by
$$\mathrm{\Phi }=\mathrm{exp}^{(k)}(z)dz^2,$$
for some positive integer $`k`$, where $`\mathrm{exp}^{(k)}(z)`$ is defined inductively by $`\mathrm{exp}^{(0)}(z)=1`$ and $`\mathrm{exp}^{(j)}(z)=\mathrm{exp}(\mathrm{exp}^{(j1)}(z))`$. Let $`𝒜=\overline{h()}^2`$. Then $`𝒜=_{j=0}^k𝒜_j`$ such that
Please see figure 7 in the appendix for the horizontal trajectories structure of $`e^{e^z}dz^2`$ and the corresponding image of the harmonic map.
###### Demonstration Proof
We may assume that $`k2`$ because $`k=1`$ is a special case of Theorem 3.1. First of all, we want to find out the domains such that $`\mathrm{\Phi }`$ can be written in the form of Lemma 3.1. Given any $`\alpha `$, $`1lk1`$ and $`(n_1,n_2,\mathrm{},n_l)^l`$, we define the open subsets $`𝒮_{(n_1,\mathrm{},n_l)}`$ inductively by
$$𝒮_{(n_1)}=\{z|\mathrm{}z>\alpha ,|\mathrm{}z2n_1\pi |<\pi \}$$
and
$$𝒮_{(n_1,\mathrm{},n_l)}=\{z𝒮_{(n_1,\mathrm{},n_{l1})}|\mathrm{}\zeta _{l1}>\mathrm{exp}^{(l1)}(\alpha ),|\mathrm{}\zeta _{l1}2n_l\pi |<\pi \},$$
where $`\zeta _{l1}=\mathrm{exp}^{(l1)}(z)`$. Then $`\zeta _l=e^{\zeta _{l1}}=\mathrm{exp}^{(l)}(z)`$ maps $`𝒮_{(n_1,\mathrm{},n_l)}`$ one-one onto the open set
$$\mathrm{\Omega }_l=\left(\{\zeta _l|\zeta _l0\}\{\zeta _l||\zeta _l|\mathrm{exp}^{(l)}(\alpha )\}\right),$$
and in terms of $`\zeta _l`$
$$\mathrm{\Phi }=\frac{\mathrm{exp}^{(kl)}(\zeta _l)}{_{j=0}^{l1}(\mathrm{log}^{(j)}\zeta _l)^2}d\zeta _l^2.$$
$`4.2`$
In particular, for $`l=k1`$,
$$\begin{array}{cc}\hfill \mathrm{\Phi }& =\frac{\mathrm{exp}(\zeta _{k1})}{_{j=0}^{k2}(\mathrm{log}^{(j)}\zeta _{k1})^2}d\zeta _{k1}^2\hfill \\ & =\mathrm{exp}\left(\zeta _{k1}2\underset{j=1}{\overset{k1}{}}\mathrm{log}^{(j)}\zeta _{k1}\right)d\zeta _{k1}^2\hfill \end{array}$$
$`4.3`$
on
$$\mathrm{\Omega }_{k1}=\left(\{\zeta _{k1}|\zeta _{k1}0\}\{\zeta _{k1}||\zeta _{k1}|\mathrm{exp}^{(k1)}(\alpha )\}\right).$$
A further transformation
$$\eta =\zeta _{k1}2\underset{j=1}{\overset{k1}{}}\mathrm{log}^{(j)}\zeta _{k1}$$
will put the Hopf differential into the form of Lemma 3.1 and we can conclude on the boundary behaviour of the harmonic map $`h`$. However, to ensure that there are no other ideal boundary point, we need to show that $`h`$ is quasiconformal in certain domain.
In order to do so, given any $`\beta >>1`$, we define $`E_\beta =\{\zeta _{k1}\mathrm{\Omega }_{k1}|\mathrm{}\eta (\zeta _{k1})>\beta \}`$ and claim that for any $`\alpha `$, there are simply-connected domains $`V_0\stackrel{~}{V}_0\{z|\mathrm{}z>\alpha 1\}`$ such that
$$\stackrel{~}{V}_0=\{\mathrm{}z<\alpha 1\}\left(\underset{(n_1\mathrm{},n_{k1})}{}\stackrel{~}{T}_{(n_1\mathrm{},n_{k1})}\right)$$
and
$$V_0=\{\mathrm{}z<\alpha \}\left(\underset{(n_1\mathrm{},n_{k1})}{}T_{(n_1\mathrm{},n_{k1})}\right),$$
where $`\stackrel{~}{T}_{(n_1\mathrm{},n_{k1})}`$, respectively $`T_{(n_1\mathrm{},n_{k1})}`$, is the component of the preimage of $`E_{\beta +1}`$, respectively $`E_\beta `$, under the map $`\mathrm{exp}^{(k1)}(z)`$ corresponding to the branch of $`\mathrm{log}`$ given by $`(n_1,\mathrm{},n_{k1})`$. Moreover, there are constants $`C_0`$, $`M_0`$, $`\delta `$ with $`C_0`$ and $`\delta >0`$ such that
$$\underset{\stackrel{~}{V}_0}{sup}R_zC_0,\underset{V_0}{inf}R_z\delta ,\text{and}\underset{V_0}{inf}(wv)(z)M_0,$$
$`4.4`$
where $`R_z`$ is the maximal $`\mathrm{\Phi }`$-radius at $`z`$, $`w=\mathrm{log}|h|`$ and $`e^{2v}|dz|^2`$ is the Poincaré metric on $`\stackrel{~}{V}_0`$. If the claim is true, then the last inequality of (4.4) implies that $`wvM_0`$ for all $`zV_0`$, and hence, by Theorem 2.1, one can conclude that $`h`$ is quasiconformal on $`V_0`$.
To prove the claim, we note that for $`zT_{(n_1,\mathrm{},n_{k1})}`$ then the image of $`z`$ under $`\mathrm{exp}^{(k1)}`$ is in $`E_\beta `$. By Corollary 2.1 as in the proof of Theorem 3.1, we conclude that $`R_z\delta >0`$ for some $`\delta >0`$ independent of $`(n_1,\mathrm{},n_{k1})`$. On the other hand, since $`e^{(k)}(z)e^{(k1)}(0)`$ if $`\mathrm{}z\mathrm{}`$, we also have $`R_z\delta `$ by Corollary 2.1 if $`\mathrm{}z\alpha `$ by choosing a possible smaller $`\delta `$. The second inequality of (4.4) is proved.
Let $`zV_0`$, then either $`\mathrm{}z=\alpha `$ or the image of $`z`$ under $`\mathrm{exp}^{(k1)}`$ is on the boundary of $`E_\beta `$. In the first case, $`e^{2w(z)}|e^{(k1)}(z)|C`$ for some constant $`C>0`$ independent of $`z`$. Hence it is easy to see that $`w(z)v(z)M_0`$ for some constant $`M_0`$ because $`\stackrel{~}{V}_0`$ is strongly hyperbolic. In the second case, then we can proceed as in the proof of Theorem 3.1 and obtain the third inequality in (4.4).
To prove the first inequality in (4.4), we let
$$V_{k1}=\mathrm{\Omega }_{k1}E_\beta \stackrel{~}{V}_{k1}=\mathrm{\Omega }_{k1}E_{\beta +1}.$$
Then it is easy to see that $`V_{k1}\stackrel{~}{V}_{k1}`$ are simply-connected domains in $`\mathrm{\Omega }_{k1}`$ and there is $`C_{k1}`$ such that
$$\underset{\stackrel{~}{V}_{k1}}{sup}R_zC_{k1}.$$
$`4.5`$
In fact, for all $`z\stackrel{~}{V}_{k1}`$, there is a divergent path $`\gamma `$ in $`\stackrel{~}{V}_{k1}`$ such that
$$L_\mathrm{\Phi }(\gamma )<C_{k1}.$$
Now, for $`l=k2`$, we consider subsets in $`\mathrm{\Omega }_{k2}`$ containing the preimage of $`V_{k1}`$ and $`\stackrel{~}{V}_{k1}`$ under the exponential map $`\zeta _{k1}=\mathrm{exp}(\zeta _{k2})`$. It is clear from the property of the exponential map that
$$\begin{array}{cc}\hfill V_{k2}& =\mathrm{exp}^1(V_{k1})\left[\left(\overline{V_{k2}^+}\overline{V_{k2}^{}}\right)\mathrm{\Omega }_{k2}\right]\hfill \\ \hfill \stackrel{~}{V}_{k2}& =\mathrm{exp}^1(\stackrel{~}{V}_{k1})\left[\left(\overline{V_{k2}^+}\overline{V_{k2}^{}}\right)\mathrm{\Omega }_{k2}\right]\hfill \end{array}$$
are simply-connected domains in $`\mathrm{\Omega }_{k2}`$ such that
$$V_{k2}^\pm V_{k2}\stackrel{~}{V}_{k2},$$
where, for $`l=1,\mathrm{},k1`$,
$$\begin{array}{cc}\hfill V_{l1}^+& =\{z𝒮_{(n_1,\mathrm{},n_{l1})}|\mathrm{}\zeta _{l1}<\mathrm{exp}^{(l1)}(\alpha ),\mathrm{}\zeta _{l1}>0\}\hfill \\ \hfill V_{l1}^{}& =\{z𝒮_{(n_1,\mathrm{},n_{l1})}|\mathrm{}\zeta _{l1}<\mathrm{exp}^{(l1)}(\alpha ),\mathrm{}\zeta _{l1}<0\}.\hfill \end{array}$$
We note that, for $`l=1,\mathrm{},k1`$,
$$𝒮_{(n_1,\mathrm{},n_{l1})}=V_{l1}^+V_{l1}^{}\left[_{n_l}\left(\overline{𝒮_{(n_1,\mathrm{},n_{l1},n_l)}}𝒮_{(n_1,\mathrm{},n_{l1})}\right)\right].$$
We want to show that there exists $`C_{k2}^{}>0`$ such that for all $`z\stackrel{~}{V}_{k2}`$, there is a divergent path $`\gamma `$ in $`\stackrel{~}{V}_{k2}`$ with $`L_\mathrm{\Phi }(\gamma )<C_{k2}^{}`$. This will immediately implies that
$$\underset{\stackrel{~}{V}_{k2}}{sup}R_zC_{k2}$$
for some $`C_{k2}>0`$. To prove this, we note that for all $`r_0>\mathrm{exp}^{(k2)}(\alpha )`$,
$$\begin{array}{cc}\hfill _{r=r_0,\mathrm{\hspace{0.17em}0}<\theta <\pi }|\mathrm{\Phi }||d\zeta _{k2}|& C_{r=r_0,\mathrm{\hspace{0.17em}0}<\theta <\pi }\frac{rd\theta }{|\zeta _{k2}||\mathrm{log}\zeta _{k2}|\mathrm{}|\mathrm{log}^{(k2)}\zeta _{k2}|}\hfill \\ & \frac{C}{\mathrm{log}r_0\mathrm{}\mathrm{log}^{(k2)}r_0}\hfill \\ & 0\text{ as }r_0+\mathrm{}.\hfill \end{array}$$
Then for all point $`\zeta _{k2}V_{k2}^\pm `$, it can be connected to a point on the vertical line $`\{\mathrm{}\zeta _{k2}=\mathrm{exp}^{(k2)}(\alpha )\}`$ by a circular arc with uniformly bounded $`\mathrm{\Phi }`$-length. By using $`\zeta _{k1}=\mathrm{exp}(\zeta _{k2})`$ to map a point on $`\mathrm{\Omega }_{k1}`$, we can find a divergent path in $`\stackrel{~}{V}_{k1}`$ with $`\mathrm{\Phi }`$-length bounded by $`C_{k1}`$. Lifting this path to $`\stackrel{~}{V}_{k2}`$ and together with the circular arc, we find a path starting from any point in $`V_{k2}^\pm `$ a divergent path in $`\stackrel{~}{V}_{k2}`$. The same is obviously true for other $`z\stackrel{~}{V}_{k2}`$ since they belong to $`\mathrm{exp}^1(\stackrel{~}{V}_{k1})`$. This proves our assertion that
$$\underset{\stackrel{~}{V}_{k2}}{sup}R_zC_{k2}$$
for some $`C_{k2}>0`$.
Continue in this way, for all $`lk1`$ we can define $`V_l`$, $`\stackrel{~}{V}_l`$ and $`V_l^\pm `$, such that
$$V_0=\mathrm{exp}^1(V_1)\text{and }\stackrel{~}{V}_0=\mathrm{exp}^1(\stackrel{~}{V}_1)\{z|\alpha 1<\mathrm{}z<\alpha \}.$$
Moreover, we can prove inductively that
$$\underset{\stackrel{~}{V}_l}{sup}R_zC_l$$
for some constant $`C_l`$. Finally, it is easy to see that the $`R_zC`$ for some contant $`C`$ if $`\alpha 1<\mathrm{}z<\alpha `$ because such a point can be joined by a line with bounded $`|\varphi |`$-length to a point in $`\mathrm{exp}^1(\stackrel{~}{V}_1)`$. This completes the proof of the first inequality in (4.4).
Now we can study the structure of the boundary points of $`h()`$. Firstly, for any $`(n_1,\mathrm{},n_{k1})^{k1}`$, using Lemma 3.1, we can argue as before to conclude that there exists monotone sequence $`p_{(n_1,\mathrm{},n_{k1});j_0}`$, $`j_0`$ and $`p_{(n_1,\mathrm{},n_{k1});+}`$, $`p_{(n_1,\mathrm{},n_{k1});}`$ in $`^2`$ such that, for $`\beta `$ sufficiently large,
$$\begin{array}{cc}& \overline{h(𝒮_{(n_1,\mathrm{},n_{k1})}\{z|\mathrm{}\eta (z)\beta \})}^2\hfill \\ & =\{p_{(n_1,\mathrm{},n_{k1});j_0}\}_{j_0}\{p_{(n_1,\mathrm{},n_{k1});+},p_{(n_1,\mathrm{},n_{k1});}\}\hfill \end{array}$$
$`4.6`$
$$\underset{j_0\pm \mathrm{}}{lim}p_{(n_1,\mathrm{},n_{k1});j_0}=p_{(n_1,\mathrm{},n_{k1});\pm }.$$
Moreover, if $`z_n𝒮_{(n_1,\mathrm{},n_{k1})}\{z|\mathrm{}\eta (z)\beta \}`$ and $`z_n\mathrm{}`$ then
$$h(z_n)\{\begin{array}{cc}p_{(n_1,\mathrm{},n_{k1});+},\hfill & \text{if }\mathrm{}\eta (z_n)\mathrm{}\hfill \\ p_{(n_1,\mathrm{},n_{k1});},\hfill & \text{if }\mathrm{}\eta (z_n)\mathrm{}\hfill \end{array}$$
In fact, we conclude by (4.3) that the energy density of $`h`$ is bounded on the set $`z𝒮_{(n_1,\mathrm{},n_{k1})}`$ such that $`a<\mathrm{}\eta (z)<b`$ and $`|\mathrm{}\eta (z)|R`$ for any $`a`$, $`b`$ and $`R`$ provided $`R`$ is large enough. Hence we still have $`h(z_n)p_{(n_1,\mathrm{},n_{k1});\pm }`$ if $`\mathrm{}\eta (z_n)\pm \mathrm{}`$, $`z_n𝒮_{(n_1,\mathrm{},n_{k1})}`$ and $`a<\mathrm{}\eta (z_n)<b`$.
Secondly, for any $`n_1,\mathrm{},n_{k2}^{k2}`$ and for any $`j_1`$, the map $`\zeta _{k1}=\mathrm{exp}(\zeta _{k2})`$ will map
$$\{z𝒮_{(n_1,\mathrm{},n_{k2})}|\mathrm{}\zeta _{k2}(z)>\mathrm{exp}^{(k1)}(\alpha ),|\mathrm{}\zeta _{k2}(z)(2j_1+1)\pi |<\pi \}$$
one-one onto
$$\stackrel{~}{\mathrm{\Omega }}_{k1}=\left(\{\zeta _{k1}|\zeta _{k1}0\}\{\zeta _{k1}||\zeta _{k1}|\mathrm{exp}^{(k1)}(\alpha )\}\right).$$
The corresponding curves given by $`\mathrm{}\eta =\beta `$ in $`𝒮_{(n_1,\mathrm{},n_{k2},j_1)}`$ and $`𝒮_{(n_1,\mathrm{},n_{k2},j_11)}`$ give us two branches of curve $`\gamma _+`$ and $`\gamma _{}`$ satisfying $`\mathrm{}\eta =\beta `$ on $`\stackrel{~}{\mathrm{\Omega }}_{k1}\{\mathrm{}\zeta _{k1}>0\}`$ and $`\stackrel{~}{\mathrm{\Omega }}_{k1}\{\mathrm{}\zeta _{k1}<0\}`$ respectively. Joining the two branches of curve by a compact curve $`\gamma `$ in $`\stackrel{~}{\mathrm{\Omega }}_{k1}`$, for instance a circular arc with sufficiently large radius centered at the origin, gives a subset $`U`$ with $`U=\gamma _+\gamma _{}\gamma `$ on which $`h`$ is quasi-conformal. By (4.6), $`h`$ will maps $`U`$ to a curve in $`^2`$ such that if $`\mathrm{}\eta \mathrm{}`$ the image under $`h`$ will tends to the point $`p_{(n_1,\mathrm{},n_{k2},j_1+1);}`$, and if $`\mathrm{}\eta \mathrm{}`$ the image under $`h`$ will tends to the point $`p_{(n_1,\mathrm{},n_{k2},j_1);+}`$. As in the proof of Theorem 3.1, we see that $`p_{(n_1,\mathrm{},n_{k2},j_1+1);}=p_{(n_1,\mathrm{},n_{k2},j_1);+}`$ which will be denoted by $`p_{(n_1,\mathrm{},n_{k2});j_1}`$. It is then not hard to see that
$$\overline{h(S_{(n_1,\mathrm{},n_{k1})})}^2=\{p_{(n_1,\mathrm{},n_{k1});j_0}\}_{j_0}\{p_{(n_1,\mathrm{},n_{k1});+},p_{(n_1,\mathrm{},n_{k1});}\},$$
and
$$𝒜_0=\{p_{(n_1,\mathrm{},n_{k2},n_{k1});j_0}\}_{(n_1,\mathrm{},n_{k1})^{(k1)},j_0},$$
is countable and discrete. Now for each $`(n_1,\mathrm{},n_{k2})^{(k1)}`$ the set
$$\{p_{(n_1,\mathrm{},n_{k2});j_1}\}_{j_1}$$
is monotone in $`j_1`$ and we denote $`p_{(n_1,\mathrm{},n_{k2});\pm }=lim_{j_1\pm \mathrm{}}p_{(n_1,\mathrm{},n_{k2});j_1}`$. Since the Hopf differential on $`𝒮_{(n_1,\mathrm{},n_{k1})}`$ is of the same form (4.3), by the proof of Lemma 3.1, for each $`(n_1,\mathrm{},n_{k1})`$, there is a point $`z_{(n_1,\mathrm{},n_{k1})}𝒮_{(n_1,\mathrm{},n_{k1})}`$ and two consecutive points in $`\{p_{(n_1,\mathrm{},n_{k1};j_0)}\}_{j_0}`$ such that
$$\mathrm{}\eta (z_{(n_1,\mathrm{},n_{k1})})=\beta $$
$`4.7`$
$$\left|\mathrm{}\eta (z_{(n_1,\mathrm{},n_{k1})})\right|\pi ,$$
$`4.8`$
and that the distance from $`h\left(z_{(n_1,\mathrm{},n_{k1})}\right)`$ to the geodesic joining these two consecutive points is bounded by $`C_1`$ for some constant $`C_1>0`$ which is independent of $`(n_1,\mathrm{},n_{k1})`$. From (4.7) and (4.8) we have
$$\underset{j_1\pm \mathrm{}}{lim}h\left(z_{(n_1,\mathrm{},n_{k2},j_1)}\right)=p_{(n_1,\mathrm{},n_{(k2)});\pm }.$$
$`4.9`$
Using (4.2) and (4.9), we can argue as before to conclude that
$$p_{(n_1,\mathrm{},n_{k3},j_2+1);}=p_{(n_1,\mathrm{},n_{k3},j_2);+}$$
which will be denoted by $`p_{(n_1,\mathrm{},n_{k3});j_2}`$, and
$$\overline{h(S_{(n_1,\mathrm{},n_{k2})})}^2=𝒜_0\{p_{(n_1,\mathrm{},n_{k3},n_{k2});+},p_{(n_1,\mathrm{},n_{k3},n_{k2});}\}.$$
Let
$$𝒜_1=\{p_{(n_1,\mathrm{},n_{k2});j_1}\}_{(n_1,\mathrm{},n_{k2})^{k2},j_1}.$$
Then $`𝒜_21`$ is countable and each point in $`𝒜_1`$ is an isolated accumulation point of $`𝒜_0`$. The accumulation points of $`𝒜_1`$ are $`p_{(n_1,\mathrm{},n_{k3});j_2}`$, $`(n_1,\mathrm{},n_{k3})^{(k2)}`$ and $`j_2`$. Continue in this way, we can find $`𝒜_j^2`$, $`0jk`$ such that each $`𝒜_j`$ is countable and discrete for $`0jk1`$ and $`𝒜_j`$ consists of all isolated accumulation points of $`𝒜_{j1}`$ for $`1jk`$. Moreover,
$$\overline{h()}^2=_{j=0}^k𝒜_j.$$
Finally, We want to prove that $`𝒜_k`$ consists of only one point. From the proof, we can see that $`𝒜_k`$ consists of at most two points $`p`$ and $`q`$ satisfying
$$\underset{y\mathrm{}}{lim}h(\sqrt{1}y)=p$$
and
$$\underset{y\mathrm{}}{lim}h(\sqrt{1}y)=q.$$
Since
$$\mathrm{exp}^{(k1)}(t)=\mathrm{exp}^{(k1)}(0)+tg(t)$$
on $`|t|1`$, where $`g(t)`$ is analytic. One can proceed as in the proof of Theorem 4.1 to show that $`p=q`$ and
$$\overline{h(\{z|\mathrm{}z0\}}^2=\{p\}.$$
Hence $`𝒜_k`$ is a singleton and this completes the proof of the theorem.
###### Remark 4.2
The Theorem 4.1 is also true for the Hopf differential
$$\mathrm{exp}^{(k1)}(e^zdz^2)=\mathrm{exp}^{(k)}(z)\underset{j=1}{\overset{k1}{}}\left[\mathrm{exp}^{(j)}(z)\right]^2dz^2.$$
In fact, the proof is much easier and can be done by induction since the form of the Hopf differential is not change under the map $`\zeta =e^z`$.
###### Remark 4.3
The Theorem 4.1 is not necessary true in general. In fact, it becomes very complicated for the general form as in Theorem 1.2. Even for $`\mathrm{\Phi }=P(z)\mathrm{exp}^{(k)}(z)dz^2`$, the Theorem 4.1 need modification. For instance, if $`P(z)=\sqrt{1}`$, then the same argument as in the proof of Theorem 4.1 and using Lemma 4.2 instead of Lemma 4.1 on the region $`\{z<\alpha \}`$, we see that the set $`𝒜_k`$ consists of two points whether than one. So the best to hope for is that $`𝒜_k`$ has at most two points for the general form in Theorem 1.2.
### §5 Harmonic diffeomorphisms on hyperbolic plane
The result in §2, in particular Proposition 2.1, can be applied to study a conjecture of Schoen, which says that any quasi-symmetric homeomorphism on $`𝕊^1`$ can be extended to a unique quasi-conformal harmonic diffeomorphism on $`^2`$. The existence part of the conjecture is still open, but there are many partial results, see \[Ak, L-T 1--3, T-W 2, S-T-W, H-W, Y\]. Schoen’s conjecture can be reformulated as follows. Let $`\text{BQD}(^2)`$ be the space of holomorphic quadratic differentials $`\mathrm{\Phi }`$ on $`^2`$ such that
$$|\mathrm{\Phi }|=\underset{z^2}{sup}\mathrm{\Phi }(z)<\mathrm{}$$
where $`\mathrm{\Phi }(z)`$ is the norm of $`\mathrm{\Phi }`$ at $`z`$ with respect to the Poincaré metric. In \[Wn\], the third author proved that for any $`\mathrm{\Phi }\text{BQD}(^2)`$, there is a unique quasi-conformal harmonic diffeomorphism $`u`$ on $`^2`$ with $`\mathrm{\Phi }`$ as Hopf differential. This defines a map $`𝔅`$ from $`\text{BQD}(^2)`$ to the universal Teichmüller space $`𝒯`$ by sending $`\mathrm{\Phi }`$ to the the class of quasi-symmetric homeomorphism containing the boundary value of $`u`$. The existence part of the conjecture of Schoen is equivalent to the surjectivity of the map $`𝔅`$. Let $``$ be a subset of $`\text{BQD}(^2)`$, we say that $`𝔅`$ is proper on $``$ if for any $`\mathrm{\Phi }_n`$ with $`|\mathrm{\Phi }_n|\mathrm{}`$ we have $`d_𝒯(𝔅(\mathrm{\Phi }_n),0))\mathrm{}`$, where $`d_𝒯`$ is the Teichmüller metric on $`𝒯`$. It is not hard to see that $`𝔅`$ is surjective if $`𝔅`$ is proper on $`\text{BQD}(^2)`$. It is also not hard to see that if $`𝔅`$ is proper on the set of $`\mathrm{\Phi }\text{BQD}(^2)`$ with $`_^2\mathrm{\Phi }𝑑v_^2<\mathrm{}`$ or even on the set $`\mathrm{\Phi }=\varphi dz^2`$ with $`\varphi `$ to be a polynomial, then $`𝔅`$ is proper on $`\text{BQD}(^2)`$. Here, we identify $`^2`$ with $`𝔻`$ with the Poincaré metric. For the sake of completeness, we give a proof of this fact below. Denote
$$=\{\mathrm{\Phi }\text{BQD}(^2)|_^2||\mathrm{\Phi }||dv_^2<\mathrm{}\}.$$
Note that $`\mathrm{\Phi }=\varphi dz^2`$, then $`_^2\mathrm{\Phi }𝑑v_^2=_𝔻|\varphi |𝑑x𝑑y`$.
###### Proposition 5.1
Let
$$𝒢=\{\mathrm{\Phi }|\mathrm{\Phi }=\varphi dz^2,\varphi \text{ is a polynomial}\}.$$
Then
In particular, if $`𝔅`$ is proper on $`𝒢`$, then $`𝔅`$ is surjective.
###### Demonstration Proof
(i) First we prove that if $`𝔅`$ is proper on $`𝒢`$, then $`𝔅`$ is proper on $``$. Let $`\mathrm{\Phi }_n`$ such that $`|\mathrm{\Phi }_n|\mathrm{}`$. Suppose that there is a constant $`C_1`$ such that $`d_𝒯(𝔅(\mathrm{\Phi }_n),0)C_1`$ for all $`n`$. Since $`𝔅`$ is continuous, there exist $`\delta _n>0`$ such that if $`|\mathrm{\Phi }_n\mathrm{\Psi }|\delta _n`$, then $`d_𝒯(𝔅(\mathrm{\Psi }),0)C_1+1`$. Hence it is sufficient to prove that $`𝒢`$ is dense in $``$. Let $`\mathrm{\Phi }=\varphi dz^2`$, then
$$_𝔻|\varphi |𝑑x𝑑y=_^2\mathrm{\Phi }𝑑v_^2<\mathrm{}.$$
Apply the mean value inequality on the disk $`𝔻_{z,r}`$ with center at $`z`$ and radius $`r=\frac{1}{2}(1|z|)`$, we can conclude that $`\mathrm{\Phi }(z)0`$ uniformly as $`|z|1`$. For $`0<R<1`$, let $`\mathrm{\Phi }_R(z)=\mathrm{\Phi }(Rz)`$. For any $`ϵ>0`$, we can find $`1>\delta >0`$ such that if $`1\delta |z|<1`$ then $`\mathrm{\Phi }(z)\frac{1}{2}ϵ`$. Then for $`1\frac{1}{2}\delta |z|<1`$ and for $`R`$ large enough, so that $`R|z|1\delta `$
$$\mathrm{\Phi }_R(z)=\frac{(1|z|^2)^2}{(1|Rz|^2)^2}\mathrm{\Phi }(Rz)\frac{1}{2}ϵ.$$
On the other hand, for $`|z|1\frac{1}{2}\delta `$, $`\varphi _R(z)\varphi (z)`$ uniformly, as $`R1`$. Hence we can find $`R`$ large enough, so that
$$\mathrm{\Phi }_R(z)\mathrm{\Phi }(z)ϵ$$
for all $`z𝔻`$. Hence $`|\mathrm{\Phi }_R\mathrm{\Phi }|_{BQD}ϵ.`$ But $`\mathrm{\Phi }_R`$ is analytic on $`|z|<\frac{1}{R}`$ which is large than 1. So it can be approximated uniformly on $`𝔻`$ by polynomials. This completes the proof of (i)
(ii) We will prove that if $`𝔅`$ is proper on $``$, then $`𝔅`$ is proper on $`\text{BQD}(^2)`$. Let $`\mathrm{\Phi }`$ and let $`𝔅(\mathrm{\Phi })=[f]`$ where $`f`$ is a quasi-symmetric homeomorphism of $`𝕊^1`$ fixing $`1,i,i`$. Let $`d_𝒯([f],0)=C_1`$. Then there exist smooth quasi-symmetric functions $`g_k`$ fixing $`1,i,i`$ such that $`g_kf`$ in $`C^\alpha `$ norm for some $`1>\alpha >0`$ and such that $`d_𝒯([g_k],0)C_2`$ which depends only on $`C_1`$. Moreover, $`C_1\mathrm{}`$ if and only if $`C_2\mathrm{}`$. These follow from theorem 2 and remark (1) in \[D-E\]. By theorem 6.4 in \[L-T 3\], see also \[T-W 2\], for each $`k`$ there exists a unique $`\mathrm{\Psi }_k\text{BQD}(^2)`$ such that $`𝔅(\mathrm{\Psi }_k)=[g_k]`$ with $`\mathrm{\Psi }_k`$. By the assumption, we have $`|\mathrm{\Psi }_k|C_3`$ for all $`k`$, where $`C_3`$ depends only on $`C_1`$. Note that $`\mathrm{\Psi }_k`$ is the Hopf differential of quasi-conformal harmonic diffeomorphism on $`^2`$ with boundary value $`g_k`$. Hence $`\mathrm{\Psi }_k(z)\mathrm{\Phi }(z)`$ for all $`z𝔻`$ and so
$$\underset{k\mathrm{}}{lim\; sup}|\mathrm{\Psi }_k||\mathrm{\Phi }|.$$
From this, it is easy to see that $`𝔅`$ is proper on $`\text{BQD}(^2)`$.
(iii) We will prove that if $`𝔅`$ is proper on $`\text{BQD}(^2)`$, then $`𝔅`$ is surjective. Let $`[f]`$ be a class of quasi-symmetric homeomorphism on $`𝕊^1`$ such that $`[f]`$ is in the closure of $`𝔅(\text{BQD}(^2))`$. Then there exists $`f_n`$ quasi-symmetric homeomorphisms on $`𝕊^1`$ fixing $`1,i,i`$ such that $`f_nf`$ uniformly, and $`[f_n]=𝔅(\mathrm{\Phi }_n)`$. Since $`[f_n]`$ are uniformly bounded on $`𝒯`$, $`\mathrm{\Phi }_n`$ are uniformly bounded in $`\text{BQD}(^2)`$. By theorem 13 in \[Wn\], the quasi-conformal harmonic diffeomorphisms $`u_n`$ with Hopf differentials $`\mathrm{\Phi }_n`$ has complex dilatation $`\mu _n`$ satisfying $`|\mu _n|\mu <1`$ for some constant $`\mu `$ independent of $`n`$. Passing to a subsequence if necessary, $`u_n`$ converges uniformly on $`\overline{𝔻}`$ to a quasi-conformal harmonic diffeomorphism on $`^2`$ with boundary value $`f`$. Hence $`[f]`$ is in $`𝔅(\text{BQD}(^2))`$. Combine with the theorem 4.1 in \[T-W 2\], we conclude that $`𝔅`$ is surjective.
###### Proposition 5.2
Let $`\mathrm{\Phi }_n\text{BQD}(^2)`$ satisfying $`_^2\mathrm{\Phi }_n<\mathrm{}`$ and $`|\mathrm{\Phi }_n|\mathrm{}`$. Suppose for all $`k>0`$,
$$\underset{n\mathrm{}}{lim}\frac{_{U_n}\mathrm{\Phi }_n𝑑v_^2}{_^2\mathrm{\Phi }_n𝑑v_^2}=0$$
where
$$U_n=\{z^2|\mathrm{\Phi }_n(z)k|\mathrm{\Phi }_n|^{\frac{3}{4}}\}.$$
Then $`d_𝒯(𝔅(\mathrm{\Phi }_n),0)\mathrm{}.`$
###### Demonstration Proof
Identify $`^2`$ with the unit disk $`𝔻`$ equipped with the Poincaré metric. Let $`\mathrm{\Phi }_n=\varphi _ndz^2`$. Let $`\mu _{n,0}`$ be the supremum of the modulus of the complex dilation of $`𝔅(\mathrm{\Phi }_n)`$, then
$$d_𝒯(𝔅(\mathrm{\Phi }_n),0)=\frac{1}{2}\mathrm{log}\frac{1+\mu _{n,0}}{1\mu _{n,0}}.$$
Since
$$_𝔻|\varphi _n|𝑑x𝑑y=_^2\mathrm{\Phi }_n𝑑v_n<\mathrm{},$$
we can apply the main inequality in \[R-S\] and conclude that
$$\frac{1}{1+\mu _{n,0}}_^2\mathrm{\Phi }_n𝑑v_^2_^2\frac{1}{1+|\mu _n|}\mathrm{\Phi }_n𝑑v_^2$$
$`5.1`$
where $`\mu _n`$ is the complex dilatation of the quasi-conformal harmonic diffeomorphism with Hopf differential $`\mathrm{\Phi }_n`$. Let $`1>\delta >0`$ and $`k>1`$ be fixed numbers. Define
$$D_n=\{z^2|\mathrm{\Phi }_n(z)k\mathrm{\Phi }_n_{BQD}^{\frac{3}{4}}\},$$
$$U_n=\{z^2|\mathrm{\Phi }_n(z)<k||\mathrm{\Phi }_n||_{BQD}^{\frac{3}{4}}\}.$$
We have
$$_^2\frac{1}{1+|\mu _n|}\mathrm{\Phi }_n𝑑v_^2=\left(_{D_n}+_{U_n}\right)\frac{1}{1+|\mu _n|}\mathrm{\Phi }_ndv_^2.$$
$`5.2`$
By Proposition 2.1, for any $`zD_n`$, the maximal $`\mathrm{\Phi }_n`$-radius $`R_{n,z}`$ satisfies $`R_{n,z}C_1k`$ for some absolute constant $`C_1>0`$. By the result on page 63 of \[Hn\], we have $`|\mu _n|\eta (k)`$ for some constant $`\eta (k)`$ such that $`\eta (k)1`$ as $`k\mathrm{}`$. Hence we have
$$_{D_n}\frac{1}{1+|\mu _n|}\mathrm{\Phi }_n𝑑v_^2\frac{1}{1+\eta (k)}_{D_n}\mathrm{\Phi }_n𝑑v_^2.$$
$`5.3`$
For any $`0<\delta <1`$, there exists $`n_0`$ such that if $`nn_0`$ then
$$_{U_n}\mathrm{\Phi }_n𝑑v_^2\delta _^2\mathrm{\Phi }_n𝑑v_^2.$$
Combine this with (5.1)–(5.3), we have for $`nn_0`$,
$$\frac{1}{1+\mu _{n,0}}_^2\mathrm{\Phi }_n𝑑v_^2\frac{1}{1+\eta (k)}_{D_n}\mathrm{\Phi }_n𝑑v_^2+\delta _^2\mathrm{\Phi }_n𝑑v_^2.$$
Hence
$$\frac{1}{1+\mu _{n,0}}\frac{1}{1+\eta (k)}+\delta .$$
Let $`k\mathrm{}`$, and then let $`\delta 0`$, we have
$$\underset{n\mathrm{}}{lim\; sup}\frac{1}{1+\mu _{n,0}}\frac{1}{2}.$$
Since $`0\mu _{n,0}<1`$ for all $`n`$, we have $`lim_n\mathrm{}\mu _{n,0}=1`$. From this, it is easy to see that $`d_𝒯(𝔅(\mathrm{\Phi }_n),0)\mathrm{}.`$
###### Corollary 5.1
Let $``$ as in Proposition 5.1. Let $`\mathrm{\Phi }_n=\varphi _ndz^2`$. Suppose $`|\mathrm{\Phi }_n|=1`$. Let $`A_n=_^2\mathrm{\Phi }_n𝑑v_^2=_𝔻|\varphi _n|𝑑x𝑑y`$. Suppose $`lim_n\mathrm{}\varphi _n/A_n=\psi `$ such that $`_𝔻|\psi |𝑑x𝑑y=1`$, then $`𝔅(t_n\mathrm{\Phi }_n)\mathrm{}`$ for any $`t_n\mathrm{}`$. In particular $`𝔅`$ is proper on any finite dimensional subspace of $``$.
###### Demonstration Proof
For any $`1>\delta >0`$, we can find $`1>r_0>0`$, such that
$$_{𝔻_{r_0}}|\psi |𝑑x𝑑y=1\delta ,$$
where $`𝔻_{r_0}=\{|z|<r_0\}`$. We have
$$_{𝔻_{r_0}}\frac{|\varphi _n|}{A_n}12\delta ,$$
$`5.4`$
provided $`n`$ is large enough. For any $`k>0`$, let $`ϵ_n=kt_n^{\frac{1}{4}}`$, where $`t_n\mathrm{}`$. Let
$$U_n=\{z^2|t_n\mathrm{\Phi }_n(z)k|t_n\mathrm{\Phi }_n|^{\frac{3}{4}}\}=\{z^2|\mathrm{\Phi }_n(z)ϵ_n\}.$$
Note that if $`\mathrm{\Phi }_n=\varphi _ndz^2`$, then for $`zU_n𝔻_{r_0}`$
$$|\varphi _n|(z)C_1(1r_0)^2\mathrm{\Phi }_n(z)ϵ_n(1r_0)^2$$
for some absolute constant $`C_1`$. Hence by (5.4),
$$\begin{array}{cc}\hfill _{U_n}\mathrm{\Phi }_n𝑑v_^2& _{U_n𝔻_{r_0}}|\varphi _n|𝑑x𝑑y+_{𝔻𝔻_{r_0}}|\varphi _n|𝑑x𝑑y\hfill \\ & \frac{\pi C_1ϵ_n}{(1r_0)^2}+2\delta A_n.\hfill \end{array}$$
$`5.5`$
Since $`|\mathrm{\Phi }_n|=1`$, by applying the mean value inequality to $`|\varphi _n|(z)`$ at a point $`z_n`$ with $`\mathrm{\Phi }_n(z_n)\frac{1}{2}`$, we conclude that $`A_nC_2`$ for some abolute constant $`C_2>0`$. By the definition of $`ϵ_n`$, we have $`ϵ_n0`$ as $`n\mathrm{}`$. Hence (5.5) implies that we have
$$\begin{array}{cc}\hfill \underset{n\mathrm{}}{lim\; sup}\frac{_{U_n}t_n\mathrm{\Phi }_n𝑑v_^2}{_^2||t_n\mathrm{\Phi }_n|dv_^2}& =\underset{n\mathrm{}}{lim\; sup}\frac{_{U_n}\mathrm{\Phi }_n𝑑v_^2}{A_n}\hfill \\ & 2\delta .\hfill \end{array}$$
Since $`\delta `$ is arbitrary, $`t_n\mathrm{\Phi }_n`$ satisfies the conditions in the Proposition 5.2. Hence the first part of corollary is proved.
To prove the second part of the corollary, let $``$ be a finite dimensional subspace of $``$ with basis $`\mathrm{\Psi }_1,\mathrm{},\mathrm{\Psi }_k`$. If $`\mathrm{\Phi }_n`$ is such that $`|\mathrm{\Phi }_n|\mathrm{}`$, then $`\mathrm{\Phi }_n=t_n\stackrel{~}{\mathrm{\Phi }}_n`$ such that $`|\stackrel{~}{\mathrm{\Phi }}_n|=1`$ and $`\stackrel{~}{\mathrm{\Phi }}_n=_{j=1}^ka_{n,j}\mathrm{\Psi }_j`$ with $`_{j=1}^ka_{n,j}^2`$ being uniformly bounded. From this, it is easy to see that the second part of the corollary follows.
Let $`\mathrm{\Phi }=\varphi dz^2`$. Define
$$|\mathrm{\Phi }|_{\mathrm{}}=\underset{a𝔻}{inf}\underset{z𝔻}{sup}|\varphi (z)|\frac{|1\overline{a}z|^4}{(|1|a|^2)^2}=\underset{a𝔻}{inf}\underset{\zeta 𝔻}{sup}|\stackrel{~}{\varphi }_a(\zeta )|,$$
where $`\stackrel{~}{\varphi }_ad\zeta ^2=\mathrm{\Phi }`$ and $`\zeta =\frac{za}{1\overline{a}z}`$.
$$|\mathrm{\Phi }|_{L^1}=_𝔻|\varphi |𝑑x𝑑y=_^2\mathrm{\Phi }𝑑v_{H^2}.$$
Note that
$$_𝔻|\varphi |𝑑x𝑑y=_𝔻|\stackrel{~}{\varphi }_a|𝑑x𝑑y$$
if $`\varphi `$ and $`\stackrel{~}{\varphi }_a`$ are related as above.
###### Corollary 5.2
Let $`_1=\{\mathrm{\Phi }BQD||\mathrm{\Phi }|=1,|\mathrm{\Phi }|_{\mathrm{}}<\mathrm{}\}`$. Let $`\mathrm{\Phi }_n_1`$, and $`t_n\mathrm{}`$ be a sequence such that $`|\mathrm{\Phi }_n|_{\mathrm{}}|\mathrm{\Phi }_n|_{L^1}^2=o(t_n^{\frac{1}{4}})`$, then $`d_𝒯(𝔅(t_n\mathrm{\Phi }_n),0)\mathrm{}`$. In particular, if $`|\mathrm{\Phi }_n|_{\mathrm{}}|\mathrm{\Phi }_n|_{L^1}^2C`$ for some constant $`C`$ independent of $`n`$, then $`d_𝒯(𝔅(t_n\mathrm{\Phi }_n),0)\mathrm{}`$, for any $`t_n\mathrm{}`$.
###### Demonstration Proof
For each $`n`$, by the definition, by a linear fractional transformation of $`𝔻`$ if necessary, we may assume that $`\mathrm{\Phi }_n=\varphi _ndz^2`$, with $`\left(sup_{z𝔻}|\varphi _n|(z)\right)|\varphi _n|_{L^1}^2=o(t_n^{\frac{1}{4}})`$. Let $`M_n=sup_{z𝔻}|\varphi _n|(z)`$ and $`I_n=|\varphi _n|_{L^1}`$. We claim that
$$I_n^2C_1M_n$$
$`5.6`$
for all $`n`$ and for some absolute constant $`C_1`$. Fix $`n`$, take $`r_0`$ such that
$$_{𝔻(r_0)}|\varphi _n|𝑑x𝑑y=\frac{1}{2}I_n,$$
where $`𝔻(r_0)=\{z||z|<r_0\}`$. Since $`|\mathrm{\Phi }_n|=1`$, we have
$$\frac{1}{2}I_n=_{𝔻(r_0)}|\varphi _n|𝑑x𝑑yC_2_0^{r_0}(1r)^2\frac{C_2}{1r_0}$$
$`5.7`$
where $`C_2`$ is an absolute constant. On the other hand
$$\frac{1}{2}I_n=_{𝔻𝔻(r_0)}|\varphi _n|𝑑x𝑑yM_n2\pi (1r_0).$$
Hence
$$\begin{array}{cc}\hfill \frac{1}{2}I_n& \frac{C_2}{1r_0}\hfill \\ & \frac{2\pi M_nC_2}{\frac{1}{2}I_n}.\hfill \end{array}$$
From this (5.6) follows.
Now for any $`k>0`$, let
$$\begin{array}{cc}\hfill U_n& =\{z^2|t_n\mathrm{\Phi }_n(z)k|t_n\mathrm{\Phi }_n|^{3/4}\}\hfill \\ & =\{z^2|\mathrm{\Phi }_n(z)kϵ_n|\mathrm{\Phi }_n|^{3/4}\}\hfill \end{array}$$
where $`ϵ_n=t_n^{\frac{1}{4}}`$. Let $`\delta _n=\left(\frac{kϵ_n}{M_nI_n^2}\right)^{\frac{1}{2}}.`$ By (5.6) and the fact that $`ϵ_n0`$, we have $`\delta _n0`$ as $`n\mathrm{}`$. As in the proof of Corollary 5.1, $`I_nC_3`$ for some absolute constant $`C_3`$. Hence $`\delta _nI_n^10`$ as $`n\mathrm{}`$. If $`n`$ is large enough so that $`\delta _nI_n^1<1`$, then we take $`r_n`$ such that $`1r_n=\delta I_n^1`$. Denote $`𝔻(r_n)=\{z||z|<r_n\}`$. Then
$$\begin{array}{cc}\hfill _{(𝔻𝔻_n)U_n}|\varphi _n|𝑑x𝑑y& 2\pi M_n(1r_n^2)\hfill \\ & 4\pi (M_nI_n^2)\delta _nI_n.\hfill \end{array}$$
$`5.8`$
On the other hand, as in the proof of (5.7), we have
$$\begin{array}{cc}\hfill _{𝔻(r_n)U_n}|\varphi _n|𝑑x𝑑y& C_4kϵ_n(1r_n)^1\hfill \\ & =C_4kϵ_n\delta _n^1I_n\hfill \end{array}$$
for some constant $`C_4`$ independent of $`n`$. Hence
$$\begin{array}{cc}\hfill _{U_n}|\varphi _n|𝑑x𝑑y& C_5\left\{M_nI_n^2\delta _n+kϵ_n\delta _n^1\right\}I_n\hfill \\ & =2C_5\left\{kϵ_nM_nI_n^2\right\}^{\frac{1}{2}}I_n\hfill \end{array}$$
for some constant $`C_5`$ independent of $`n`$. Hence if $`n`$ is large enough,
$$\frac{_{U_n}|\varphi _n|𝑑x𝑑y}{_𝔻|\varphi _n|𝑑x𝑑y}2C_5\left\{kϵ_nM_nI_n^2\right\}^{\frac{1}{2}}.$$
By the assumption, the right side of the above inequality tends to zero as $`n0`$ and the corollary follows.
Example 1. Let $`\mathrm{\Phi }_n=\varphi _ndz^2=c_nn^2z^ndz^2`$, where $`c_n`$ is chosen so that $`|\mathrm{\Phi }_n|=1`$. Direct computations show that $`C^1c_nC`$ for some positive constant $`C`$ independent of $`n`$. Then $`sup_{z𝔻}|\varphi _n|(z)=c_nn^2`$, and $`_𝔻|\varphi _n|𝑑x𝑑y=\frac{2\pi n^2c_n}{n+2}`$. Hence the $`|\mathrm{\Phi }_n|_{\mathrm{}}|\mathrm{\Phi }_n|_{L_1}^2C^{}`$ for some constant $`C^{}`$ independent of $`n`$. Hence by Corollary 5.2, for any subsequence $`n_k`$ and for any $`t_k\mathrm{}`$ we have $`d_𝒯(𝔅(t_k\mathrm{\Phi }_{n_k}),0)\mathrm{}`$.
Example 2. Consider $`(z1)^ndz^2`$. Then direct computation shows
$$\underset{z𝔻}{sup}(1|z|)^2|z1|^n=c_nn^22^n$$
where $`C^1c_nC`$ for some constant $`C>0`$ independent of $`n`$. Let $`\mathrm{\Phi }_n=c_n^{}n^22^n(z1)^ndz^2=\varphi _n(z)dz^2`$, where $`c_n^{}`$ is chosen so that $`|\mathrm{\Phi }_n|=1`$. Note that $`\stackrel{~}{C}^1c_n^{}\stackrel{~}{C}`$ for some $`\stackrel{~}{C}>0`$ independent of $`n`$. Let $`a=\frac{n+84\sqrt{n+4}}{n}`$, note that $`1<a<0`$.
$$\begin{array}{cc}\hfill \underset{z𝔻}{sup}|\varphi _n(z)|\frac{|1az|^4}{(1a^2)^2}& =c_n^{}n^22^n\underset{0\theta 2\pi }{sup}|e^{i\theta }1|^n|1ae^{i\theta }|^4(1a^2)^2\hfill \\ & =c_n^{}n^2\underset{0\theta 2\pi }{sup}\mathrm{sin}^n\frac{\theta }{2}\left[(1a)^2+4a\mathrm{sin}^2\frac{\theta }{2}\right]^2(1a^2)^2\hfill \\ & =c_n^{}n^2\underset{0t1}{sup}t^n\left[(1a)^2+4at^2\right]^2(1a^2)^2.\hfill \end{array}$$
$`5.9`$
Let $`f(t)=t^n\left[(1a)^2+4at^2\right]^2`$, then $`f(0)=0`$. Suppose $`f(t)`$, $`0t1`$ attains its maximum at $`t_0(0,1)`$, then $`f^{}(t_0)=0`$, and
$$nt_0^{n1}\left[(1a)^2+4at_0^2\right]^2+t_0^n2\left[(1a)^2+4at_0^2\right]8at_0=0.$$
Hence $`t_0^2=\frac{n(1a^2)}{4a(n+4)}=1`$ by the choice of $`a`$, which is impossible. Hence, for $`0t1`$, $`f(t)`$ attains its maximum at $`t=1`$. By (5.9), we have
$$\begin{array}{cc}\hfill |\mathrm{\Phi }_n|_{\mathrm{}}& c_n^{}n^2\left[(1a)^2+4a\right]^2(1a^2)^2\hfill \\ & =c_n^{}n^2(1+a)^2(1a)^2.\hfill \end{array}$$
Since $`a<0`$, $`1a>1`$. Also
$$1+a=1\frac{n+84\sqrt{n+4}}{n}=\frac{8+4\sqrt{n+4}}{n}.$$
Hence
$$|\mathrm{\Phi }_n|_{\mathrm{}}C_1n$$
$`5.10`$
for some constant $`C_1`$ independent of $`n`$. On the other hand
$$\begin{array}{cc}\hfill _𝔻|z1|^n𝑑x𝑑y& =_{2\pi }^0_0^1|r^2+12r\mathrm{cos}\theta |^{\frac{n}{2}}r𝑑r𝑑\theta \hfill \\ & =_\pi ^\pi _0^1|r^2+1+2r\mathrm{cos}\theta |^{\frac{n}{2}}r𝑑r𝑑\theta \hfill \\ & _0^{\frac{1}{\sqrt{n}}}_0^1(r+1)^n|1\frac{r\theta ^2}{(1+r)^2}|^{\frac{n}{2}}r𝑑r𝑑\theta \hfill \\ & \frac{C_2}{\sqrt{n}}_0^1(r+1)^n\left(1\frac{n}{2}\frac{r}{n(1+r)^2}\right)r𝑑r\hfill \\ & C_32^nn^{\frac{3}{2}}\hfill \\ & C_4n^{\frac{1}{2}}\hfill \end{array}$$
for some constants $`C_2C_4`$ independent of $`n`$. By Corollary 5.2, we also have $`𝔅(t_k\mathrm{\Phi }_{n_k})\mathrm{}`$ for all subsequence $`n_k`$ and for all $`t_k\mathrm{}`$.
In the last section of \[Wn\], it was proved that $`𝔅`$ is continuous, and in section 4 of \[T-W 2\], it was proved that the image of $`𝔅`$ is open and $`𝔅`$ is a diffeomorphism from $`\text{BQD}(^2)`$ into $`𝒯`$. From the proof of the proposition 14 in \[Wn\], $`𝔅`$ is in fact uniformly continuous on bounded subsets of $`\text{BQD}(^2)`$. On the other hand, we have the following:
###### Proposition 5.3
Let $`R>0`$ and let
$$B(R)=\{\mathrm{\Phi }\text{BQD}(^2)||\mathrm{\Phi }|R\}.$$
Let $`𝒯(R)=𝔅\left(B(R)\right)`$. Then $`𝔅^1`$ is uniformly continuous on $`𝒯(R)`$.
###### Demonstration Proof
For any complex measurable function $`\mu `$ on $`𝔻`$ such that $`\mu _{\mathrm{}}<1`$, denote $`F^\mu `$ to be the unique quasi-conformal map on $`𝔻`$ with boundary value $`f^\mu `$ which fixes $`1,i,i`$. Suppose $`f^\mu `$ can be extended to a quasi-conformal harmonic diffeomorphism, then the harmonic map will be denoted by $`\stackrel{~}{F}^\mu `$ and its complex dilation is denoted by $`\stackrel{~}{\mu }`$. By theorem 13 in \[Wn\], there exists $`0<k<1`$ such that if $`[f]𝒯(R)`$ then $`\mu _{\mathrm{}}k`$, where $`\mu `$ is the complex dilatation of an extremal quasi-confomal map with boundary value in $`[f]`$.
$$𝒯^{}(R)=\{\mu |\mu \text{ is measurable, }\mu _{\mathrm{}}k\text{ and }[f^\mu ]𝒯(R)\}.$$
Note that if $`\mu 𝒯^{}(R)`$, then $`f^\mu `$ can be extended to a quasi-conformal harmonic diffeomorphism with Hopf differential in $`B(R)`$. We claim that for any $`ϵ>0`$, there is $`\delta >0`$ such that if $`\mu `$, $`\nu `$ in $`𝒯^{}(R)`$ and $`\mu \nu _{\mathrm{}}\delta `$ then $`\stackrel{~}{\mu }\stackrel{~}{\nu }_{\mathrm{}}ϵ`$. If the claim is true, then by the definition of $`d_𝒯`$ and by the method as in the proof of proposition 14 in \[Wn\], one can conclude that $`𝔅^1`$ is uniformly continuous on $`𝒯(R)`$.
First we prove the following, given $`ϵ>0`$, there is $`\delta >0`$ such that if $`\mu `$ and $`\nu `$ are in $`𝒯^{}(R)`$, then $`|\stackrel{~}{\mu }(0)\stackrel{~}{\nu }(0)|ϵ`$. Suppose not, then there is $`ϵ>0`$ and two sequences $`\mu _n`$, $`\nu _n`$ in $`𝒯^{}(R)`$ such that $`\mu _n\nu _n_{\mathrm{}}0`$, but $`|\stackrel{~}{\mu }_n(0)\stackrel{~}{\nu }_n(0)|ϵ`$. Since $`\stackrel{~}{\mu }_n_{\mathrm{}}k_1`$ and $`\stackrel{~}{\nu }_n_{\mathrm{}}k_1`$ for some $`0<k_1<1`$ by \[Wn\], passing to subsequences if necessary, $`\stackrel{~}{F}^{\mu _n}`$ and $`\stackrel{~}{F}^{\nu _n}`$ converge uniformly on $`\overline{𝔻}`$ to normalized quasi-conformal harmonic diffeomorphisms $`H_1`$ and $`H_2`$ respectively. Since $`\mu _n\nu _n_{\mathrm{}}0`$ and $`\mu _n`$, $`\nu _n`$ are in $`𝒯^{}(R)`$, $`H_1`$ and $`H_2`$ must have the same boundary value and so $`H_1=H_2`$ by \[L-T 3\]. It then follows that $`|\stackrel{~}{\mu }_n(0)\stackrel{~}{\nu }_n(0)|0`$, which is a contradiction.
Now, for any $`ϵ>0`$, let $`\delta >0`$ be as above. Let $`\mu `$, $`\nu `$ be in $`𝒯^{}(R)`$ such that $`\mu \nu _{\mathrm{}}\delta `$. Let $`a𝔻`$ and $`\varphi (z)=(za)/(1\overline{a}z)`$. Define $`\mu _1`$ and $`\nu _1`$ by
$$\mu _1(\varphi (z))=\mu (z)\left(\frac{\varphi ^{}(z)}{|\varphi ^{}(z)|}\right)^2$$
and
$$\nu _1(\varphi (z))=\nu (z)\left(\frac{\varphi ^{}(z)}{|\varphi ^{}(z)|}\right)^2.$$
Then $`f^{\mu _1}=h^\mu f^\mu \varphi ^1`$, and $`f^{\nu _1}=h^\nu f^\nu \varphi ^1`$ where $`h^\mu `$ and $`h^\nu `$ are the linear fractional transformations which map $`𝔻`$ onto itself and are chosen so that $`f^{\mu _1}`$ and $`f^{\nu _1}`$ fix $`1,i,i`$ respectively. Obviously $`f^{\mu _1}`$ and $`f^{\nu _1}`$ have quasi-conformal representatives $`\stackrel{~}{F}^{\mu _1}`$ and $`\stackrel{~}{F}^{\nu _1}`$. In fact,
$$\stackrel{~}{F}^{\mu _1}=h^\mu \stackrel{~}{F}^\mu \varphi ^1$$
and
$$\stackrel{~}{F}^{\nu _1}=h^\mu \stackrel{~}{F}^\nu \varphi ^1.$$
Moreover, the Hopf differentials of $`\stackrel{~}{F}^{\mu _1}`$ and $`\stackrel{~}{F}^{\nu _1}`$ are in $`B(R)`$. Hence $`\mu _1`$, $`\nu _1`$ are in $`𝒯^{}(R)`$ becasue $`\stackrel{~}{\mu }_1_{\mathrm{}}=\mu _{\mathrm{}}k`$ and $`\stackrel{~}{\nu }_1_{\mathrm{}}=\nu _{\mathrm{}}k`$ .
We also have
$$\stackrel{~}{\mu }_1(\varphi (z))=\stackrel{~}{\mu }(z)\left(\frac{\varphi ^{}(z)}{|\varphi ^{}(z)|}\right)^2$$
and
$$\stackrel{~}{\nu }_1(\varphi (z))=\stackrel{~}{\nu }(z)\left(\frac{\varphi ^{}(z)}{|\varphi ^{}(z)|}\right)^2.$$
Note that
$$\mu _1\nu _1_{\mathrm{}}=\mu \nu _{\mathrm{}}\delta .$$
Therefore
$$|\mu (a)\nu (a)|=|\mu _1(0)\nu _1(0)|ϵ.$$
Since $`a`$ is any point in $`𝔻`$, the claim follows.
### Appendix: Trajectories and Image Accumulation
In this appendix, seven figures of horizontal trajectories defined by a quadratic differential are shown. These pictures of trajectories are produced by programming in Mathematica. Some trajectories may be broken due to slow convergence of the algorithm. In fact it should be smoothly defined for all time. Nevertheless, the qualitative behavior of the trajectory patterns is shown clearly. In some figures, the correponding harmonic map produces an image which has a good accumulation structure on the boundary. This structure is also shown on the unit disk.
## Finitely Many Accumulations
#### Figure 1. $`\mathrm{\Phi }=e^zdz^2`$ (See Theorem 3.1 and Corollary 3.1)
This example is the basis of all others. The trajectories have a $`2\pi i`$ periodicity. The image of the corresponding harmonic map has an accumulation point at $`1`$.
#### Figure 2. $`\mathrm{\Phi }=(e^z+1)dz^2`$ (See Theorem 4.1)
#### Figure 3. $`\mathrm{\Phi }=(e^z1)dz^2`$ (See Theorem 4.1)
From this and the previous one, a lower order term may significantly change the behavior of the harmonic map. This one has two accumulation points at $`\pm 1`$ while the previous one has only one.
#### Figure 4. $`\mathrm{\Phi }=\mathrm{sinh}^2zdz^2`$ (See Theorem 4.1)
This is another example that the fundamental region is different while there are also two accumulation points.
#### Figure 5. $`\mathrm{\Phi }=(z^21)e^{z^2}dz^2`$ (See Theorem 3.1 and Corollary 3.1)
The image of this example should have 2 accumulation points.
#### Figure 6. $`\mathrm{\Phi }=e^{z^3+z^2}dz^2`$ (See Theorem 3.1 and Corollary 3.1)
The image of this example should have 3 accumulation points.
## Accumulation of accumulating points
#### Figure 7. $`\mathrm{\Phi }=e^{e^z}dz^2`$ (See Theorem 4.2)
Finally, this is an example about accumulation of accumulations. The image has infinitely many accumulation points marked by dots outside the unit circle, which in turns accumulate at $`1`$.
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# Broadcasting of entanglement in three-particle GHZ state via quantum copying
## Abstract
We introduce entanglement measures to describe entanglement in a three-particle system and apply it to studying broadcasting of entanglement in three-particle GHZ state. We show that entanglement of three-qubit GHZ state can be partially broadcasted with the help of local or non-local copying processes. It is found that non-local cloning is much more efficient than local cloning for the broadcasting of entanglement.
PACS number(s): 03.75.Fi, 03.65Bz; 32.80.Pj, 74.20.De
Quantum entanglement, first noted by Einstein-Podolsky-Rosen and Schrödinger , is one of the essential features of quantum mechanics. It has been well known that quntum entanglement plays a key role in many such applications like quantum teleportation , super-dense coding , quantum error correction , and quantum computational speedups . Recently, maximally entangled states of three particles, i.e., three-particle GHZ states, have been produced experimentally . Usually, it is difficult to obtain ideal entangled multi-particle state. An interesting question is: whether there are ways which can broadcast the entanglement of correlated systems? The answer is ok. Masiak and Knight have shown that copies of entangled pair of qubits can be genetated through using a universal quantum cloning machine (UQCM) , although the degree of entanglement of the resultant copies is substantially reduced due to a residual entanglement between the copied output and the copying machine. On the other hand, up to now there is not an appropriate measure to describe quantatively the entanglement of three and more subsystems due to the high complexity of entanglement in multi-particle system . The purpose of this paper is to propose measures of entanglement in three qubit system in terms of the entanglement tensor approach , and use these measures to investigate broadcasting of entanglement in three-particle GHZ state.
Consider a system consisting of three qubits. The density operator of the system can be expanded as a sum of tensor products in terms of Puali matrices,
$`\widehat{\rho }`$ $`=`$ $`{\displaystyle \frac{1}{8}}[\widehat{1}\widehat{1}\widehat{1}+{\displaystyle \underset{i=1}{\overset{3}{}}}\lambda _i(1)(\widehat{\sigma _i}\widehat{1}\widehat{1})`$ (5)
$`+{\displaystyle \underset{j=1}{\overset{3}{}}}\lambda _j(2)(\widehat{1}\widehat{\sigma _j}\widehat{1})+{\displaystyle \underset{k=1}{\overset{3}{}}}\lambda _k(3)(\widehat{1}\widehat{1}\widehat{\sigma _k})`$
$`+{\displaystyle \underset{i,k=1}{\overset{3}{}}}K_{ik}(1,3)(\widehat{\sigma _i}\widehat{1}\widehat{\sigma _k})+{\displaystyle \underset{i,j=1}{\overset{3}{}}}K_{ij}(1,2)`$
$`\times (\widehat{\sigma _i}\widehat{\sigma _j}\widehat{1})+{\displaystyle \underset{j,k=1}{\overset{3}{}}}K_{jk}(2,3)(\widehat{1}\widehat{\sigma _j}\widehat{\sigma _k})`$
$`+{\displaystyle \underset{i,j,k=1}{\overset{3}{}}}K_{ijk}(1,2,3)(\widehat{\sigma _i}\widehat{\sigma _j}\widehat{\sigma _k})].`$
Here $`\{\widehat{\sigma }_i,i=1,2,3\}`$ are the Pauli matrice. $`\lambda (1)`$,$`\lambda (2)`$ and $`\lambda (3)`$ are the three coherence vectors beloning to qubit 1, 2 and 3, respectively, which determine the properties of the individual particles. $`K_{ij}(1,2)`$, $`K_{jk}(2,3)`$, and $`K_{ik}(1,3)`$) are the second-rank correlation tensors which describe the correlation between qubit 1 and 2 (qubit 2 and 3, qubit 1 and 3) resprctively. $`K_{ijk}(1,2,3)`$ are the three-qubit correlation tensor.
Making use of properties of Pauli matrices $`tr\{\widehat{\sigma _i}\}=0`$, and $`tr\{\widehat{\sigma _i}\widehat{\sigma _j}\}=2\delta _{ij}`$, we get the following relations:
$`\lambda _i(1)`$ $`=`$ $`tr(\widehat{\rho }\widehat{\sigma _i}\widehat{1}\widehat{1}),`$ (6)
$`\lambda _j(2)`$ $`=`$ $`tr(\widehat{\rho }\widehat{1}\widehat{\sigma _j}\widehat{1}),`$ (7)
$`\lambda _k(3)`$ $`=`$ $`tr(\widehat{\rho }\widehat{1}\widehat{1}\widehat{\sigma _k}),`$ (8)
$`K_{ij}(1,2)`$ $`=`$ $`tr(\widehat{\rho }\widehat{\sigma _i}\widehat{\sigma _j}\widehat{1}),`$ (9)
$`K_{ik}(1,3)`$ $`=`$ $`tr(\widehat{\rho }\widehat{\sigma _i}\widehat{1}\widehat{\sigma _k}),`$ (10)
$`K_{jk}(2,3)`$ $`=`$ $`tr(\widehat{\rho }\widehat{1}\widehat{\sigma _j}\widehat{\sigma _k}),`$ (11)
$`K_{ijk}(1,2,3)`$ $`=`$ $`tr(\widehat{\rho }\widehat{\sigma _i}\widehat{\sigma _j}\widehat{\sigma _k}).`$ (12)
After performing the partial trace over the second qubit and the third qubit, from Eq.(1) we obtain the reduced density operator for the first qubit:
$$\widehat{\rho }^{(1)}=tr_{2,3}(\widehat{\rho })=\frac{1}{2}(\widehat{1}+\underset{i}{\overset{3}{}}\lambda _i(1)\widehat{\sigma }_i).$$
(13)
Simmilarly, we can calculate the reduced density operators $`\widehat{\rho }^{(2)}`$, $`\widehat{\rho }^{(3)}`$ for qubit 2 and 3.
Comparing the direct product $`\widehat{\rho }^{(1)}\widehat{\rho }^{(2)}\widehat{\rho }^{(3)}`$ with Eq.(1), we can identify the difference by tensors $`M(m,n)(1m<n3)`$, $`M(1,2,3)`$ defined by
$$M_{ij}(m,n)=K_{ij}(m,n)\lambda _i(m)\lambda _j(n),$$
(14)
$`M_{ijk}(1,2,3)`$ $`=`$ $`K_{ijk}(1,2,3)\lambda _i(1)M_{jk}(2,3)`$ (17)
$`\lambda _j(2)M_{ik}(1,3)\lambda _k(3)M_{ij}(1,2)`$
$`\lambda _i(1)\lambda _j(2)\lambda _k(3),`$
which indicates that $`M(1,2,3)=0`$ for any product state of three qubits. From the above, we also find that a three-qubit entanglement state necessarilly involves entanglement between any two qubits. Hence entanglement measures in a three-qubit system should involve both an inter-three-qubit entanglement measure and an inter-two-qubit entanglement measure.
Based on $`M(1,2,3)`$, $`M(m,n)(1m<n3)`$, we introduce an inter-three-qubit entanglement measure $`E_3`$ and an inter-two-entanglement measure $`E_2`$ in the following form,
$`E_3={\displaystyle \frac{1}{4}}{\displaystyle \underset{i,j,k=1}{\overset{3}{}}}M_{ijk}(1,2,3)M_{ijk}(1,2,3),`$ (18)
$`E_2(m,n)={\displaystyle \frac{1}{3}}{\displaystyle \underset{i,j=1}{\overset{3}{}}}M_{ij}(m,n)M_{ij}(m,n).`$ (19)
It is easy to check that $`E_2`$ and $`E_3`$ obey all conditions as entanglement measures indicated in Ref.. These measures are invariant under local unitary transformations of the subsystems and varies between $`0`$ (product states) and $`1`$ (maximum entangled states). $`E_3`$ quantifies the three-qubit entanglement . The larger $`E_3`$ is, the stronger the three-qubit entanglement is. And $`E_2(m,n)`$ quatitfies the entanglement between any two qubits $`m,n`$ in the three qubit system. The larger $`E_2(m,n)`$ is, the stronger the entanglement between two qubit m,n is.
As an example of three-qubit entanglement, let us consider the GHZ state $`\psi =(111+000)/\sqrt{2}`$. Making use of Eqs.(2)-(9), from Eqs.(10) and (11) we can get the nonzero entanglement tensors, $`M_{xxx}=1`$, $`M_{xyy}=M_{yxy}=M_{yxx}=1`$, $`M_{zz}(1,2)=M_{zz}(2,3)=M_{zz}(1,3)=1`$. All other entanglement tensors vanish as well as the coherence vectors. Then from Eqs.(12)and (13) we find that
$$E_3=1,E_2(1,2)=E_2(2,3)=E_2(1,3)=\frac{1}{3},$$
(20)
which indicate that the GHZ state is the maximally entangled inter-three-qubit state and contains i nter-two-qubit entanglement, as expected.
In what follows we investigate broadcasting of entanglement in a three-particle GHZ state via quantum copying. We shallconsider the two cases of local cloning and non-local cloning, and assume that the three qubits are prepared in the GHZ state:
$$\psi =\frac{1}{\sqrt{2}}(111_{1_02_03_0}+000_{1_02_03_0}).$$
(21)
Firstly, we consider the case of local cloning. In this case, the three qubit $`1_0`$, $`2_0`$ and $`3_0`$ is copied by the UQCM denoted by the local unitary transformations :
$`U_10_{a_0}0_{a_1}X_x`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}00_{a_0a_1}_x`$ (23)
$`+\sqrt{{\displaystyle \frac{1}{3}}}+_{a_0a_1}_x,`$
$`U_11_{a_0}0_{a_1}X_x`$ $`=`$ $`\sqrt{{\displaystyle \frac{2}{3}}}11_{a_0a_1}_x`$ (25)
$`+\sqrt{{\displaystyle \frac{1}{3}}}+_{a_0a_1}_x`$
where $`+_{a_0a_1}=(10_{a_0a_1}+01_{a_0a_1})/\sqrt{2}`$. The system labelled by $`a_0`$ is the original (input) qubit, while the other system $`a_1`$ represents the target qubit onto which the information is copied. The states of the copying machine are labelled by $`x`$. The state space of the copying machine is two dimensional and we assume that it is always in the same state $`X_x`$ initially.
Suppose that each of the three original qubits $`1_0`$, $`2_0`$ and $`3_0`$ is cloned separately by three distant local cloners $`X_1`$, $`X_2`$ and $`X_3`$. The cloner $`X_1`$ ($`X_2`$,$`X_3`$) generates out of qubit $`1_0`$ ($`2_0`$,$`3_0`$) two $`1_0`$ and $`1_1`$ ($`2_0`$ and $`2_1`$, $`3_0`$ and $`3_1`$). After cloning, we get the following total output state
$$\psi _{total}^{(out)}=\frac{1}{\sqrt{2}}\underset{i=0}{\overset{1}{}}\underset{m=1}{\overset{3}{}}[U_1(m)i_{m_0}0_{m_1}X_m_x].$$
(26)
After performing trace over the three cloners we obtain a six-qubit density operator $`\widehat{\rho }_{1_01_12_02_13_03_1}^{(out)}`$ which also describes two nonlocal three-qubit systems $`\widehat{\rho }_{1_02_03_0}`$ and $`\widehat{\rho }_{1_12_13_1}`$. The two three-qubit systems are the clones of the original three-qubit GHZ state in Eq.(15) and they are described by the density operators:
$`\widehat{\rho }_{1_02_03_0}^{(out)}`$ $`=`$ $`{\displaystyle \frac{7}{24}}(111111+000000)`$ (31)
$`+{\displaystyle \frac{7}{54}}(111000+000111)`$
$`+{\displaystyle \frac{5}{72}}(110110+011011`$
$`+101101+100100`$
$`+010010+001001).`$
Now we check whether entanglement is broadcasted. Making use of Eq.(2)-(8), we can obey nonzero entanglement tensors in the output state $`M_{xxx}=7/27`$, $`M_{xyy}=M_{yxy}=M_{yyx}=7/27`$ and $`M_{zz}(1,2)=M_{zz}(1,3)=M_{zz}(2,3)=4/9`$ while all other entanglement tensors vanish as well as the coherence vectors. The values $`E_3`$ and $`E_2(m,n)`$ are then given by
$$E_3=\frac{49}{729},E_2(1,2)=E_2(2,3)=E_2(1,3)=\frac{16}{243}.$$
(32)
which indicate that both the inter-three-qubit entanglement and inter-two-qubit entanglement in a tree-qubit systems charactrized by $`E_3`$ and $`E_2(m,m)`$, respectively, can be broadcasted via local quantum cloners, although both the inter-three-qubit entanglement and inter-two-qubit entanglement of the copied state are less than those of the original state which are given by Eq.(14).
Secondly, we consider the case of non-local cloning. In this case, the entangled state of the three-qubits is treated as a state in a larger Hilbert space and cloned as a whole. The non-local quantum copying machine is an $`N`$ dimensional quantum system, and we shall let $`X_i_x`$ ($`i=1,\mathrm{},N`$) be a set of orthonormal basis of the copying machine Hilbert space. This copier is initially prepared in a particular state $`X_x`$. The action of the cloning transformation can be specified by a unitary transformation acting on the basis vectors of the tensor product space of the original quantum system $`\varphi _i_{a_0}`$, the copier, and an additional $`N`$-dimensional system which is to become the copy (which is initially prepared in an arbitrary state $`0_{a_1}`$). The corresponding transformation $`U_2`$ is given by
$`U_2\varphi _i_{a_0}0_{a_1}X_x`$ $`=`$ $`c\varphi _i_{a_0}\varphi _i_{a_1}X_i_x`$ (35)
$`+d{\displaystyle \underset{ji}{\overset{N}{}}}(\varphi _i_{a_0}\varphi _j_{a_1}`$
$`+\varphi _j_{a_0}\varphi _i_{a_1})X_j_x,`$
where $`i=1,\mathrm{},N`$, $`c^2=2/(N+1)`$, and $`d^2=1/2(N+1)`$.
For a three-qubit system, we have eight basis vectors $`\varphi _1=000`$, $`\varphi _2=001`$, $`\varphi _3=010`$, $`\varphi _4=011`$, $`\varphi _5=100`$, $`\varphi _6=101`$, $`\varphi _7=110`$, and $`\varphi _8=111`$. So that the original three-qubit GHZ state (15) can be simply expressed as $`\psi =(\varphi _1+\varphi _8)/\sqrt{2}`$. The copying is now performed by the transformation (21) with $`N=8`$. After cloning, we get the total output state:
$$\psi _{total}^{(out)}=\frac{1}{\sqrt{2}}U_2[(\varphi _1_{a_0}+\varphi _8_{a_0})0_{a_1}X_x].$$
(36)
Then performing trace over the cloner, from Eq.(22) we obtain a six-qubit density operator $`\widehat{\rho }_{1_01_12_02_13_03_1}^{(out)}`$ which also describes two nonlocal three-qubit systems $`\widehat{\rho }_{1_02_03_0}`$ and $`\widehat{\rho }_{1_12_13_1}`$. These two three-qubit systems are the clones of the original three-qubit GHZ state in Eq.(4.1) and they are described by the same density operator as:
$`\widehat{\rho }_{1_02_03_0}^{(out)}`$ $`=`$ $`{\displaystyle \frac{1}{3}}(111111+000000)+{\displaystyle \frac{5}{18}}(111000`$ (40)
$`+000111)+{\displaystyle \frac{1}{18}}(110110+011011`$
$`+101101+100100+010010`$
$`+001001)`$
For the output state (23) we find the nonzero entanglement tensors $`M_{xxx}=5/9`$, $`M_{xyy}=M_{yxy}=M_{yyx}=5/9`$, and $`M_{zz}(1,2)=M_{zz}(1,3)=M_{zz}(2,3)=5/9`$, While all other entanglement tensors vanish as well as the coherence vectors. Then the values the entanglement measures $`E_3`$ and $`E_2(m,n)`$ are are found to be
$$E_3=\frac{25}{81},E_2(1,2)=E_2(2,3)=E_2(1,3)=\frac{25}{243},$$
(41)
which implies that entanglements are broadcasted through nonlocal quantum copying, but entanglement of the copied state is less than that of the original state. Comparing Eq.(24) with Eq.(20), we see that non-local cloning is much efficient than the local copying process for broadcasting entanglement.
Finally, it is interesting to compare the fidelity $`F_1`$ of the output density operator after local copying relative to $`\psi `$ with the fidelity $`F_2`$ of the output density operator after non-local copying relative to $`\psi `$. The fidelity of a density matrix $`\rho `$ relative to $`\psi `$ is defined by $`F=\psi \rho \psi `$. Through calculating, we find that $`F_1=91/216`$ and $`F_2=11/18`$. Therefore, we can see that the output state after non-local copying is closer to the original state than the output state after local copying.
In conclusion, we have proposed entanglement measures for a three-particle system and applied them to the study of broadcasting of entanglement for a thrr-qubit GHZ state. By using local and non-local cloning transformations, we have shown that entanglement of the three-qubit GHZ state can be locally or no-locally copied with the help of local quantum copiers or non-local quantum copiers, and that the degree of entanglement of the resultant copies is reduced. And we have found that non-local cloning is much more efficient than local cloning for the broadcasting of entanglement.
ACKNOWLEDGMENTS
This work was supported in part NSF of China, the Excellent Young-Teacher Foundation of the Educational Commission of China, ECF and STF of Hunan Province.
Email: lmkuang@sparc2.hunnu.edu.cn
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# Singlet-aided infinite resource reduction in the comparison of distant fields
## Abstract
We present a task which can be faithfully solved with finite resources only when aided by particles prepared in a particular entangled state: the singlet state. The task consists of identifying the mutual parallelity or orthogonality of weak distant magnetic fields whose absolute directions are completely unknown.
It is well known that quantum mechanics helps to reduce the resources required to accomplish certain tasks . Some problems can be solved with exponentially less resources when aided by quantum mechanics, as featured in Shor’s factorization algorithm . Other problems lead to a quadratic speed-up, such as Grover’s search algorithm . Moreover, the use of quantum entanglement can result in resource reduction in a variety of communication associated tasks. For example, in quantum dense coding , prior entanglement is used to increase the classical information capacity of a quantum bit by a factor of two. Sharing entanglement can also reduce the amount of communication needed to evaluate certain functions of distributed inputs . In this paper, we present a task which illustrates the superiority of an entanglement-based strategy in very radical terms. We show that there is in fact an infinite gap in the resources required for accomplishing the task with or without the use of a certain entangled quantum state (a singlet). Without sharing a singlet state, the task requires infinitely many qubits for error free operation, while the use of shared singlets reduces the resource requirement to at most four qubits.
Consider the situation depicted in Figure 1., where two spatially separated and disconnected regions are occupied by distant partners Alice and Bob. A third person, Eve, who has access to both separated zones, may subject these two regions to two weak uniform magnetic fields of unit strength but otherwise random direction. However, she gives Alice and Bob an important promise: the two fields are either parallel or orthogonal. In other words, if she had chosen the direction $`\widehat{n}`$ for the field applied on Alice’s side, she chooses either $`\widehat{n}`$ or any direction $`\widehat{n}_{}`$ orthogonal to $`\widehat{n}`$ for the field applied on Bob’s side. We also make the assumption that these fields are sufficiently weak so that they cannot be determined classically. The only way to determine the field direction is by means of detecting their action on quantum states. Alice and Bob are given the task of faultlessly identifying (i.e. with unit probability of success), Eve’s choice among the two alternative relative orientations of the fields. Note that this is strictly a ”quantum task” in contrast to existing examples in which quantum mechanics is used to reduce the resources required to accomplish a ”classical task”. Here, because of the weakness of the magnetic fields, there is no hope to accomplish the task classically. But, as we will show, even within the available quantum protocols, using entanglement leads to an infinite resource reduction.
We will first consider the case when Alice and Bob do not share any entanglement. Suppose Alice has a qubit $`A`$ in an initial state $`|\psi _A`$ and Bob has a qubit $`B`$ be in an initial state $`|\psi _B`$. Let $`\stackrel{}{\sigma }^{(A)}=(\sigma _x^{(A)},\sigma _y^{(A)},\sigma _z^{(A)})`$ and $`\stackrel{}{\sigma }^{(B)}=(\sigma _x^{(B)},\sigma _y^{(B)},\sigma _z^{(B)})`$, where $`\sigma _i^{(A/B)}`$ denotes the Pauli matrices of $`A/B`$. To distinguish between the two alternatives locally would thus require
$$\psi |\stackrel{}{\sigma }^{(B)}.\widehat{n}\stackrel{}{\sigma }^{(B)}.\widehat{n}_{}|\psi _B=0.$$
(1)
The only solution for this is for $`|\psi _B`$ to be an eigenstate of the operator $`\stackrel{}{\sigma }^{(B)}.\widehat{n}`$, where the direction $`\widehat{n}`$ is completely unknown. Therefore, in order to account for error free detection, Bob will need to have a number of qubits, each in an eigenstate of $`\stackrel{}{\sigma }^{(B)}.\widehat{n}`$ corresponding to a different $`\widehat{n}`$. As there are an infinite number of choices of $`\widehat{n}`$, an error free detection scheme requires Bob to hold an infinite number of qubits. It is important to note that having any or all of Alice’s or Bob’s qubits in classically correlated mixed states of the type $`_ip_i|\psi ^i_A\psi ^i|_A|\psi ^i_B\psi ^i|_B`$ will also not help in perfect discrimination of $`\stackrel{}{\sigma }^{(A)}.\widehat{n}\stackrel{}{\sigma }^{(B)}.\widehat{n}`$ and $`\stackrel{}{\sigma }^{(A)}.\widehat{n}\stackrel{}{\sigma }^{(B)}.\widehat{n}_{}`$. If that were the case, then one would have been able to choose three mutually perpendicular directions and perform quantum dense coding of capacity $`\mathrm{log}_23`$ bits per qubit. But this is not possible with a disentangled state, as shown in Ref..
Let us now describe a strategy where Alice and Bob initially share entanglement. Imagine that the qubits $`A`$ and $`B`$ possessed by Alice and Bob are prepared in a singlet state
$$|\psi ^{}=\frac{1}{\sqrt{2}}(|01|10).$$
(2)
Now Eve subjects Alice’s and Bob’s qubits to her chosen unitary transformations. Suppose she chose the pair $`\{\stackrel{}{\sigma }^{(A)}.\widehat{n},\stackrel{}{\sigma }^{(B)}.\widehat{n}\}`$, (i.e. parallel fields). Then the state shared by Alice and Bob evolves to
$`|\psi `$ $`=`$ $`\stackrel{}{\sigma }^{(A)}.\widehat{n}\stackrel{}{\sigma }^{(B)}.\widehat{n}|\psi ^{}`$ (3)
$`=`$ $`|\psi ^{},`$ (4)
where we have used the fact that a singlet state is invariant under operations $`U^{(A)}U^{(B)}`$ (i.e. when the same unitary operation $`U`$ is applied to both qubits). On the other hand, if Eve decided to apply the pair $`\{\stackrel{}{\sigma }^{(A)}.\widehat{n},\stackrel{}{\sigma }^{(B)}.\widehat{n}_{}\}`$ (i.e. perpendicular fields) the singlet will evolve to a coherent superposition of the three triplet states $`|\psi ^+,|\mathrm{\Phi }^+`$, and $`|\mathrm{\Phi }^{}`$. This can easily be seen from the fact there is always a unitary transformation $`U(\widehat{n})`$ such that
$$\stackrel{}{\sigma }^{(A)}.\widehat{n}\stackrel{}{\sigma }^{(B)}.\widehat{n}_{}=U(\widehat{n})\sigma _x^{(A)}U(\widehat{n})^{}U(\widehat{n})\sigma _{y/z}^{(B)}U(\widehat{n})^{}$$
(5)
and therefore
$$\psi ^{}|\stackrel{}{\sigma }^{(A)}.\widehat{n}\stackrel{}{\sigma }^{(B)}.\widehat{n}_{}|\psi ^{}=\psi ^{}|\sigma _x^{(A)}\sigma _{y/z}^{(B)}|\psi ^{}=0.$$
(6)
As a result, Alice and Bob can now easily check which of the two possible relative orientations Eve has chosen. If the parallel configuration $`\{\stackrel{}{\sigma }^{(A)}.\widehat{n},\stackrel{}{\sigma }^{(B)}.\widehat{n}\}`$ was applied, Alice and Bob still share a singlet state. On the contrary, if they were subject to the orthogonal configuration $`\{\stackrel{}{\sigma }^{(A)}.\widehat{n},\stackrel{}{\sigma }^{(B)}.\widehat{n}_{}\}`$, Alice and Bob now hold a state that is orthogonal to the singlet state.
For the determination which state Alice and Bob are holding, two scenaria are possible. In one, we may assume that Alice and Bob are allowed to send each other quantum particles. In this case Alice simply sends Bob her particle, and Bob then measures the projection operator onto the singlet space. If this projection is successful, then he knows that the fields were parallel, if the projection was not successful then the fields were orthogonal. However, one may also demand that Alice and Bob only share some initial entanglement in the form of singlets and that no further quantum communication is possible. In that case Alice and Bob need altogether two pairs of singlet states to complete the task. The first pair is treated as outlined above, while the second one is kept isolated from Eve and will be needed to determine whether the first pair is in a singlet state or not. This can clearly be done, as one singlet state is enough to implement a controlled-NOT gate remotely . The quantum circuit required for the (local) discrimination of the shared entangled state is shown in Figure 2. First a remote controlled-NOT gate with Alice’s qubit acting the control bit is applied. This process consumes an e-bit of entanglement. Subsequently a Hadamard transformation onto Alices qubit takes the state $`|0`$ into $`|0|1`$ and the state $`|1`$ into $`|0+|1`$. As a result of this protocol, the four Bell states are mapped into product states as follows
$`|00+|11`$ $``$ $`|00`$ (7)
$`|00|11`$ $``$ $`|10`$ (8)
$`|01+|10`$ $``$ $`|01`$ (9)
$`|01|10`$ $``$ $`|11.`$ (10)
The state discrimination is completed when Alice and Bob measure their first pair. If they both find the state $`|1`$, then they initially shared a singlet state and the magnetic fields were parallel. Therefore, even without the use of quantum communication, Alice and Bob require only two singlets to accomplish the task of determining the relative orientation of the two fields without error.
We will now prove that the singlet is the only state which allows Alice and Bob to achieve error free discrimination. The state $`|\psi `$ which Alice and Bob must share in order to accomplish the required task needs to satisfy
$$\psi |\stackrel{}{\sigma }^{(A)}.\widehat{n}\stackrel{}{\sigma }^{(B)}.\widehat{n}\stackrel{}{\sigma }^{(A)}.\widehat{n}^{^{}}\stackrel{}{\sigma }^{(B)}.\widehat{n}_{}^{^{}}|\psi =0,$$
(11)
where $`\widehat{n}`$ and $`\widehat{n}^{^{}}`$ are two completely arbitrary directions. This is because Alice and Bob need to distinguish all cases of parallel field directions from all cases of perpendicular field directions. In other words, Eve could have picked one absolute direction $`\widehat{n}`$ to apply parallel magnetic fields and a different absolute direction $`\widehat{n}^{^{}}`$ to apply the perpendicular fields. These two cases need to be perfectly distinguished. We first find out the restrictions on $`|\psi `$ which already arise from considering the special case $`\widehat{n}=\widehat{n}^{^{}}`$. In that case Eq.(11) simplifies to
$$\psi |I^A\stackrel{}{\sigma }^{(B)}.\widehat{n}^{^{\prime \prime }}|\psi =0,$$
(12)
where $`I^A`$ is the identity operator for qubit $`A`$ and $`\widehat{n}^{^{\prime \prime }}=\widehat{n}\times \widehat{n}_{}`$ (i.e. $`\widehat{n}^{^{\prime \prime }}`$ is arbitrary as $`\widehat{n}`$ itself is arbitrary). Eq.(12) restricts the class of allowed $`|\psi `$ to maximally entangled states of the qubits $`A`$ and $`B`$. If we put a maximally entangled state $`|\psi _{\text{max}}`$ in Eq.(11) and simplify, we get the condition
$`(\widehat{n}.\widehat{n}^{^{}})(\widehat{n}.\widehat{n}_{}^{^{}})`$ $``$ $`\psi |\stackrel{}{\sigma }^{(A)}.(\widehat{n}\times \widehat{n}^{^{}})`$ (13)
$`\stackrel{}{\sigma }^{(B)}.(\widehat{n}`$ $`\times `$ $`\widehat{n}_{}^{^{}})|\psi _{\text{max}}=0.`$ (14)
Substituting $`|\psi _{\text{max}}`$ in the above equation by its expansion $`c_1|\psi ^++c_2|\psi ^{}+c_3|\varphi ^++c_4|\varphi ^{}`$ in terms of Bell states gives the unique solution $`|c_2|^2=1`$. This proves that the only state which satisfies Eq.(11) is the singlet state. Thus the singlet becomes the only feasible state for error free discrimination.
Let us now briefly point out how our scheme differs from those schemes which appear to be related. The fact that our scheme can only be carried out with singlets makes it different from quantum dense coding and precision magnetic field determination , which would both work for any maximally entangled state. It is also different from the standard inability to distinguish specific entangled (and also some unentangled ) states locally by finite resources. Here, we are not given any prior sets of unknown states to distinguish, but some unknown relative orientations of fields to discriminate. We have identified the state which works best for this purpose (namely the singlet).
The task of determining the relative orientation of two magnetic fields can be generalized in many directions. First of all one may allow for more possible relative directions, and ask Alice and Bob to determine the angle between the two directions. In this case both the entanglement based as well as the disentangled strategy are unable to deliver error free answers, but it can be expected that the entangled strategy will deliver the better overall precision or the lower error rate. A further generalization of the problem could also allow for a variable strength of the magnetic field. Again it can be expected that the entanglement based strategy will be superior. It should be noted that this problem is related to that of atomic frequency standards , which can be mapped onto a problem where a field of known orientation but unknown strength should be detected with the best possible resolution. It should also be noted that Eve could have given a different promise leading to a similar gap between the entanglement based and the disentangled strategy. If Eve promises that the fields are either parallel or anti-parallel, then either the singlet state remains invariant, or it is converted (after a suitable waiting time) into the triplet state $`|01+|10`$, which in turn allows to determine the relative orientation of the fields. Summarizing, we have presented a task that be solved efficiently using shared entanglement in the form of singlet states and demonstrated that the associated cost in resources represents an infinite gap as compared to applying a classical strategy. Error free performance requires that Alice and Bob hold an infinite number of disentangled particles while an entanglement-based strategy uses either one singlet pair, if subsequent quantum communication is allowed, or, at most, two singlet pairs if only local quantum operations and the exchange of classical communication is allowed. It is quite interesting to note that very recently yet another application of entanglement which uses the $`UU`$ invariance of a singlet in an essential way has been proposed .
This work has been supported by The Leverhulme Trust, the EQUIP project of the European Union, the UK Engineering and Physical Sciences Research Council (EPSRC) and DGICYT Project No. PB-98-0191 (Spain).
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# Screened Coulomb potentials for astrophysical nuclear fusion reactions
## I Introduction
At astrophysical energies of a few $`keV`$ corresponding to stellar temperatures of several millions degrees kelvin the cross section $`\sigma \left(E\right)`$ of the predominant $`s`$-wave fusion reactions is given by
$$\sigma \left(E\right)=\frac{S\left(E\right)}{E}P\left(E\right)$$
(1)
where the astrophysical factor $`S\left(E\right)`$ embodies all the nuclear effects of the reaction and for non-resonant cases is a slowly varying function of the center-of-mass energy $`E.`$ On the other hand , the penetrability factor $`P\left(E\right)`$ embodies all atomic effects of the reaction and when the electron cloud around the fusing nuclei is ignored it is given by $`P\left(E\right)=\mathrm{exp}\left(2\pi n\right)`$ where $`n`$ is the Sommerfeld parameter.
As the astrophysical factor varies slowly with energy we usually replace it with a truncated Taylor series which will be studied extensively in the present paper
$$S\left(E\right)=S\left(0\right)+S^{^{}}\left(0\right)E+0.5S^{^{\prime \prime }}\left(0\right)E^2$$
(2)
Any error in the zero-energy astrophysical factor $`S\left(0\right)`$ is actually an error in the corresponding reaction rate in the stellar plasma, which in turn reflects linearly on the energy production rate.
In the past years there have been exhaustive efforts to extend measurements of the $`S\left(E\right)`$ towards even lower energies in order to obtain a reliable value for $`S\left(0\right)`$. This is necessary as extrapolating higher energy data to zero energies introduces an inevitable numerical error. However, at such low energies, the electron cloud that screens the fusing nuclei enhances the fusion reaction by lowering the Coulomb barrier. Consequently, disregarding its presence leads to an overestimation of $`S\left(0\right)`$. Unfortunately, even very recent experiments cannot explain the screening enhancement which exceeds all the available theoretical predictions as was recently admitted , .
Various authors have studied the influence of the atomic cloud on the cross section of low energy nuclear reaction. A qualitative study, which parametrized various atomic processes such as molecular formation, excitation and ionization, yielded a fair approximation for the possible contributions of the electronic degrees of freedom in the nuclear collision experiment. Moreover, by assuming a constant charge density around the target nucleus, a subsequent model predicted a screening shift which was compatible with the experimental data . However that assumption is an oversimplification which will be amended in the present paper. The most sophisticated approach has been a few-body treatment which established a lower (sudden) and a higher (adiabatic) limit for the screening energy transferred into the relative nuclear motion. Although more studies followed, which also extended the calculations to molecular fusion reactions, despite their mathematical rigor they could not explain the discrepancy between experimental and theoretical screening energies.
In this work there is presented a mean-field model for the study of screened nuclear reactions at astrophysical energies in the laboratory . That model agrees well with the available experimental data, thus enabling us to improve the accuracy of the associated astrophysical factor. Moreover, by means of the proposed model the effect of a superstrong magnetic field on laboratory Hydrogen fusion reactions is also investigated for the first time, yielding the associated magnetic accelerating factor. Notably, the present method gives a sufficiently high screening energy for Hydrogen fusion reactions so that the spectator nucleus take-away energy can also be taken into account.
## II Screened Coulomb potentials
After the pioneering work that established the importance of atomic effects in low energy nuclear reactions various authors have tried to create models that account for the observed enhancement. A simple model, suggested at an early stage, assumed that the electronic charge density around the target nucleus is constant, thus predicting for the nucleus-atom reaction between the atomic target $`Z_1e`$and the projectile $`Z_2e`$ a screening energy $`U_e=\left(3/2\right)Z_1Z_2e^2a^1.`$ In order to take into account the dependence of the screening radius on the charge state of the reaction participants, that model used a screening radius taken from scattering experiments so that
$$a=0.8853a_0\left(Z_1^{2/3}+Z_2^{2/3}\right)^{1/2}$$
(3)
where $`a_0`$ the Bohr radius. Although that screening energy is larger than the one predicted by the simple formula $`U_e=Z_1Z_2e^2\left(a_0/Z_1\right)^1`$ it has some very obvious defects. The assumption that the charge density is constant leads to an unnaturally sharp cut-off at a distance $`r=a`$from the center of the target nuclei, which is not born out either by theory or experiment. Moreover, atomic excitations and deformations of the target atom are totally disregarded. On the other hand normalizing the charge distribution so that the total charge is $`Z_1e`$ gives a charge density
$$\rho _0=\frac{3}{4}\frac{Z_1e}{\pi a^3}$$
(4)
In order to assess the validity of that density we can consider the hydrogen-like atom $`Z_1e`$ which will also be used in this section . The charge density at the center of the cloud of such an atom (when the electron is in its ground state) is $`\rho _0^H=`$ $`e\left(Z_1/a_0\right)^3/\pi .`$ It is obvious that for $`Z_1=Z_2=1`$ we obtain $`\rho _0\left(e/a_0^3\right)`$ and $`\rho _0^H=\left(e/a_0^3\right)/\pi ,`$ that is the simplified model in question overestimates the central density by a factor of $`\pi .`$
Consequently it is obvious that if low energy nuclear reactions are to be treated by means of a mean-field potential a more sophisticated treatment is necessary.
As a first step we consider a more plausible charge distribution:
$$\rho \left(r\right)=\rho _0\left(1\frac{r^2}{a^2}\right)$$
(5)
which takes into account the depletion of charge with respect to distance from the center. The radius $`a`$ is the screening radius given by Eq. $`\left(\text{3}\right)`$ and the charge density $`\rho _0`$ at the center of the cloud can be found by means of the normalization condition $`:`$
$$_0^a\rho \left(r\right)4\pi r^4𝑑r=Z_1e$$
(6)
This integral yields a central value of
$$\rho _0=\frac{15}{8}\frac{Z_1e}{\pi a^3}$$
(7)
Note that for a collision $`Z_1=Z_2=1`$ we have a central charge density $`\rho _0=7.68\left(e/a_0^3\right)/\pi `$ which gives an even larger core density than the constant density assumption. An alternative approach would be to consider the value $`\rho _0`$ equal to the corresponding hydrogen-like one and then calculate the screening radius using Eq. $`\left(\text{6}\right).`$ The latter treatment gives a screening radius
$$a=\left(\frac{15}{8Z_1^2\pi }\right)^{1/3}a_0$$
(8)
which is independent of the charge of the projectile. For hydrogen isotopes Eq. $`\left(\text{8}\right)`$ gives a radius of $`a=0.842a_0`$
We can calculate the electrostatic energy by solving the equation of Poisson for the above charge distribution with the appropriate boundary conditions, so that
$$\mathrm{\Phi }\left(r\right)=\frac{15}{12}\frac{Z_1e}{a}\left[\frac{3}{2}\left(\frac{r}{a}\right)^2+\frac{3}{10}\left(\frac{r}{a}\right)^4\right]$$
(9)
Whenever a bare nucleus $`Z_2e`$ impinges on the target nuclei surrounded by the electron cloud of Eq.$`\left(\text{5}\right)`$the total interaction potential in the atom-nucleus reaction channel is
$$V\left(r\right)=\frac{Z_1Z_2e^2}{r}\frac{15}{12}\frac{Z_1Z_2e^2}{a}\left[\frac{3}{2}\left(\frac{r}{a}\right)^2+\frac{3}{10}\left(\frac{r}{a}\right)^4\right]$$
(10)
Although the above potential energy is more plausible than the constant charge density one, a more reliable charge distribution should be considered which could account for various other atomic effects as well as for the atom-atom reaction channel.
Let us consider a hydrogen-like atom with atomic number $`Z_1`$. When the wave function of the electron is given by $`\mathrm{\Psi }_{nl}(r,\theta )`$ then the charge density around the point-like nucleus is
$$\rho (r,\theta )=e\left|\mathrm{\Psi }_{nl}(r,\theta )\right|^2$$
(11)
by which it is obvious that both the previous screening model and that of Ref. are imperfect. If we solve the equation of Poisson for hydrogen atoms (or hydrogen-like ions) whose electron is in its ground $`\left(1s\right)`$ state we obtain
$$\mathrm{\Phi }_{00}\left(r\right)=\frac{e}{r}+\frac{e}{r}\left(1+\frac{r}{2r_0}\right)\mathrm{exp}\left(r/r_0\right)$$
(12)
where the screening radius is
$$r_0=\frac{a_0}{2Z_1}$$
(13)
If a positive projectile $`Z_2e`$ interacts with the above screened nucleus then the total potential energy is
$$V_{00}\left(r\right)=\frac{Z_1Z_2e^2}{r}\frac{Z_2e^2}{r}+\frac{Z_2e^2}{r}\left(1+\frac{r}{2r_0}\right)\mathrm{exp}\left(\frac{r}{r_0}\right)$$
(14)
On the other hand if we assume that the electron is in an excited state $`\left(2s\right)`$ then the potential energy is found to be:
$$V_{10}\left(r\right)=\frac{Z_1Z_2e^2}{r}\frac{Z_2e^2}{r}+\frac{Z_2e^2}{r}\left(1+\frac{3}{8}\frac{r}{r_0}+\frac{r^2}{16r_0^2}+\frac{r^3}{64r_0^3}\right)\mathrm{exp}\left(\frac{r}{2r_0}\right)$$
(15)
It should be emphasized that in the derivation of the above potentials we have assumed an unperturbed wavefunction of the target nuclei, throughout the tunnelling process. In fact at astrophysical energies the electron cloud responds rapidly and by the time tunneling begins the nuclei are so close that the wavefunction is actually that of a hydrogen-like atom with charge $`Z_1^{}=\left(Z_1+Z_2\right)`$ and a screening radius $`r_0^{}=a_0/2Z_1^{}`$
## III Nuclear reactions at astrophysical energies
At astrophysical energies reactions between light nuclei take place via $`s`$-interactions, thus enabling us to investigate them by means of the WKB.
If we assume that a bare nucleus $`Z_2e`$ collides at very low energy $`E`$ with a screened nucleus whose electron is in its ground state then the tunneling probability according to the WKB method is:
$$P\left(E\right)=\mathrm{exp}\left[\frac{2\sqrt{2\mu }}{\mathrm{}}_R^{r_c\left(E\right)}\sqrt{V_{00}\left(r\right)E}𝑑r\right]$$
(16)
We can assume that the lower limit of the WKB integral is given in terms of the mass number $`A`$ of the reacting nuclei : $`R=1.4\left(A_1^{1/3}+A_2^{1/3}\right).`$ For most practical purposes this lower bound is set equal to zero as all the nuclear effects of the fusion reaction are included in the cross section factor.
The classical turning point can be obtained by equating the relative collision energy $`E`$ with the potential energy of the interaction. The collision energy is set equal to the Gamow peak of the corresponding reaction in the plasma so that:
$$V_{00}\left(r_c\right)=1.220\left(Z_1^2Z_2^2AT_6^2\right)^{1/3}keV$$
(17)
where $`A`$ the reduced mass number and $`T_6`$ the temperature in million degrees kelvin. For a wide range of light nuclei we have performed extensive numerical solutions for Eq. $`\left(\text{17}\right)`$ as well as numerical integrations of Eq. $`\left(\text{16}\right).`$ At astrophysical energies, just as is the case with the Debye-Hückel model in plasma conditions, the results indicate that throughout the potential barrier the potential energy $`V_{00}\left(r\right)`$ of Eq. $`\left(\text{14}\right)`$ can be safely replaced by the much simpler formula:
$$V_{00}\left(r\right)\frac{Z_1Z_2e^2}{r}\frac{Z_1^{}Z_2e^2}{a_0}$$
(18)
Therefore the WKB penetration factor can be written as:
$$P\left(E\right)=\mathrm{exp}\left[\frac{2\sqrt{2\mu }}{\mathrm{}}_R^{r_c\left(E\right)}\sqrt{\frac{Z_1Z_2e^2}{r}\frac{Z_1^{}Z_2e^2}{a_0}E}𝑑r\right]$$
(19)
The equation for the classical turning point is modified accordingly:
$$\frac{Z_1Z_2e^2}{r_c}=1.220\left(Z_1^2Z_2^2AT_6^2\right)^{1/3}keV$$
(20)
where we have ignored the screening shift given by:
$$U_e=\frac{Z_1^{}Z_2e^2}{a_0}$$
(21)
It is now obvious that the relative energy of the reaction has been increased by $`U_e`$. In that case the penetration factor can be easily found to be:
$$f_{1s}\left(E\right)\mathrm{exp}\left[\pi n\left(E\right)\frac{U_e}{E}\right]$$
(22)
where the subscripts indicate the excitation state of the target atom. If we follow the same methodology for the $`2s`$ state we obtain
$$f_{2s}\left(E\right)\mathrm{exp}\left[\pi n\left(E\right)\frac{U_e}{4E}\right]$$
(23)
The much simpler potential model of Eq. $`\left(\text{10}\right)`$ gives a screening factor:
$$f_0\left(E\right)\mathrm{exp}\left[\pi n\left(E\right)\frac{\stackrel{~}{U_e}}{E}\right]$$
(24)
with an energy shift of
$$\stackrel{~}{U_e}=\frac{15}{8}\frac{Z_1Z_2e^2}{a}$$
(25)
where $`a`$ is given either from Eq. $`\left(\text{3}\right)`$ or Eq. $`\left(\text{8}\right)`$
## IV Magnetically catalyzed screening
By now it is obvious that any shift $`U_e<<E`$ of the interaction potential energy $`V\left(r\right)`$
$$V\left(r\right)=\frac{Z_1Z_2e^2}{r}U_e$$
(26)
accelerates the fusion cross section of hydrogen isotopes by a factor $`f_{1s}\left(E\right)`$ given by Eq. $`\left(\text{22}\right).`$ That observation will prove very useful in the study of the effects of a superstrong magnetic field on laboratory hydrogen fusion reactions which follows.
As a matter of fact under such extreme conditions the electron-screening cloud is deformed in the sense that it becomes compressed perpendicular and parallel to the magnetic field so that the screening potential energy for the strongly magnetized hydrogen atom is
$$U_e(\rho ,z;\alpha )=\frac{e^2}{\widehat{\rho }}\frac{1}{\sqrt{2\pi }}_0^{\mathrm{}}\frac{\mathrm{exp}\left[\frac{1}{2}\left(\frac{\overline{\rho }^2}{1+u}+\frac{\overline{z}^2}{\alpha ^2+u}\right)\right]}{\left(1+u\right)\sqrt{\alpha ^2+u}}𝑑u$$
(27)
where $`\rho ,z`$ are the coordinates in a cylindrical frame of reference whose origin coincides with the point-like nucleus of the hydrogen atom.
The natural length unit in the above formula is of course the cyclotron radius so that $`\overline{\rho }=\rho /\widehat{\rho },\overline{z}=z/\widehat{\rho }`$ , and $`\alpha `$ is a parameter which depends on the magnetic field and is determined by the variational method. The above formula was shown to be reliable for very strong fields whereas it becomes inaccurate below the threshold of the intense magnetic field regime given by:
$$B_{IMF}=4.7\times 10^9G.$$
(28)
In Ref. potential $`\left(\text{27}\right)`$ was applied at zero relative energies in order to obtain the mean-life times of hydrogen isotopes in neutron star surfaces. However, a more recent work used that potential in a problem where the relative energies were of the order of $`keV`$ showing that for energies $`E>0.5keV`$ and fields of the order of $`B_{12}=0.047`$ ($`B_{12}`$ being the field measured in $`10^{12}G)`$ the classical turning point is so deep inside the cloud that the screening shift can be considered constant and equal to the value of the potential at the center of the cloud given in Ref.
$$U_e(0,0;\alpha )=\frac{e^2}{\widehat{\rho }}\frac{2}{\sqrt{2\pi }}\frac{\mathrm{ln}\left(\alpha +\sqrt{\alpha ^21}\right)}{\sqrt{\alpha ^21}}$$
(29)
In the present work that approximation has been tested for various other fields and energies. The results show that for fields as high as $`B_{12}=4.7`$ and interaction energies $`E>0.5keV`$ the screening effect is independent of the angle at which the projectile enters the electron cloud and can be considered equal to Eq. $`\left(\text{29}\right).`$
Therefore if the target hydrogen nuclei are in such a magnetic field the reaction is going to be accelerated by a factor
$$f_{1s}\left(E\right)\mathrm{exp}\left[\pi n\left(E\right)\frac{U_e(0,0;\alpha )}{E}\right]$$
(30)
Figures 1 and 2 depict the acceleration of the $`pp`$ and $`dd`$ reactions respectively for various magnetic fields and interaction energies. Especially for the $`pp`$ reaction it is obvious that even in such a strong field the cross section is still significantly small. Namely, as the corresponding zero energy astrophysical factor is $`S_{pp}\left(0\right)4\times 10^{22}keVbarns,`$ the screening effect in a superstrong field $`B_{12}=4.7`$ can only increase $`S_{pp}\left(0\right)`$ by roughly one order of magnitude compared to the unmagnetized case.
The $`dd`$ reaction, on the other hand, can be significantly affected by such a magnetic field as it is already much faster than the $`pp`$ one. At very low energies the increase can be as high as two orders of magnitude compared to the unmagnetized case.
## V The astrophysical factor of $`dD`$ nuclear reactions.
Despite the fact that the reactions $`H^2(d,p)H^3,H^2(d,n)He^3`$ have been investigated since the early days of accelerators, the effect of screening on the associated astrophysical $`S\left(E\right),`$ which will eventually be used in theoretical calculations, is still under investigation. In the discussion that follows we will show that our model is compatible with the experimental data of that reaction.
The appropriate treatment of a low-energy experiment should take into account screening effects in order to calculate the respective values of $`S\left(E\right).`$ As a matter of fact once a screening model and the associated screening energy $`U_e`$ are adopted the corrected bare-nucleus astrophysical factor of the experiment is actually given by
$$S_b\left(E\right)=E\sigma \left(E\right)\mathrm{exp}\left(2\pi n\right)\mathrm{exp}\left(\pi n\frac{U_e}{E}\right)$$
(31)
Then Eq. $`\left(\text{2}\right)`$ is fitted to the data corrected through Eq. $`\left(\text{31}\right)`$ in order to obtain the zero-energy coefficient $`S\left(0\right).`$
Any effort to extrapolate from higher-energy data or fit all the uncorrected data with formula $`\left(\text{2}\right)`$is bound to induce errors.
There are three different ways to analyze low energy fusion data which must of course be consistent with each other. We will apply those methods on the available data for $`dd`$ reactions ($`E>2keV)`$ and compare them with the analytic model proposed in the present paper. First we note that for energies $`E>20keV`$ any screening correction is meaningless since the exponential term of Eq. $`\left(\text{31}\right)`$ is very close to unity at such high energies. Therefore we can obtain the asymptotic behavior of the astrophysical factor by using the available high-precision experimental data for higher energies which yielded
$$S_b\left(E\right)=55.49\left(0.46\right)+0.094\left(0.0054\right)E$$
(32)
We can now reasonably assume that this should be a fair approximation of the bare-nucleus astrophysical provided its use consistently describes the low-energy experimental data. In fact the screened value of $`S\left(E\right)`$ will now be given by
$$S\left(E\right)=\left(55.49+0.094E\right)\mathrm{exp}\left(\pi n\frac{U_e^{as}}{E}\right)$$
(33)
where the screening energy $`U_e^{as}`$ is determined by fitting Eq. $`\left(\text{33}\right)`$ to the uncorrected data of Ref. , so that $`U_e^{as}=0.019\left(0.003\right)keV`$ with $`\chi ^2=0.028.`$
The second method which will corroborate the validity of the proposed models entails fitting all four parameters $`S\left(0\right),S^{^{}}\left(0\right),S^{^{\prime \prime }}\left(0\right),U_e`$ simultaneously to the uncorrected experimental data. Thus we obtain a screening energy of $`U_e^{all}=0.017\left(0.003\right)`$$`keV`$and a bare nucleus astrophysical factor:
$$S_b\left(E\right)=54.54\left(1.39\right)+0.608\left(0.265\right)E0.026\left(0.026\right)$$
(34)
with $`\chi ^2=0.011.`$ Obviously, the two previous approaches give results which are compatible with each other as expected. Figure 3 shows that both the previous two fits provide a satisfactory description of the screening effect.
The third method is a straightforward application of the theoretical models derived in the present paper. However, in order to apply those models on the experimental data we have to take into account that the data refer to a molecular target while our models refer to atomic ones. Hence, we have to allow for the energy which will be carried away by the spectator nuclei plus the reduction due to the molecular binding energy. Although this assumption has been argued against, the actual energy reduction for a deuteron molecular target has been calculated by a Coulomb explosion process to be of the order of $`44eV.`$ Therefore modifying our models for a molecular deuteron target we derive a screening energy $`U_e=0.010keV`$ (Eq. $`\left(\text{21}\right))`$ and $`\stackrel{~}{U_e}=0.016keV`$ (Eq. $`\left(\text{25}\right))`$ which are in reasonably good agreement with the experimentally obtained values. We can now fit the formula
$$S\left(E\right)=\left[S\left(0\right)+S^{^{}}\left(0\right)E+0.5S^{^{\prime \prime }}\left(0\right)E^2\right]\mathrm{exp}\left(\pi n\frac{U_e}{E}\right)$$
(35)
by using the screening shift of our models. The results are as follows
$`U_e=0.010`$
$$S_b\left(E\right)=57.3\left(0.41\right)+0.160\left(0.125\right)E0.0056\left(0.002\right)E^2$$
(36)
with $`\chi ^2=0.013`$ and
$`\stackrel{~}{U_e}=0.016`$
$$S_b\left(E\right)=54.93\left(0.38\right)+0.537\left(0.1149\right)E0.0225\left(0.007\right)E^2$$
(37)
with $`\chi ^2=0.011`$
Although our models are fairly compatible with the experiment there is an inevitably degree of uncertainty in the associated astrophysical factors due to the actual amount of energy that is carried away by the spectator nuclei of the molecular target. In any case the models proposed here turn out to provide a simple and effective way of describing fusion reactions between hydrogen-like atoms.
## VI Conclusions
This work proposes a simple and efficient model for the study of the screening enhancing effect on low-energy nuclear fusion reactions. In that model, the fusing atoms are considered hydrogen-like atoms whose electron probability density is used in Poisson’s equation in order to derive the corresponding screened Coulomb potential energy. This way atomic excitations and deformations of the reaction participants can be taken into account. The derived mean-field potentials are then treated semiclassically, by means of the WKB, in order to derive the screening enhancement factor which is shown to be compatible with the experimentally obtained one for the $`H^2(d,p)H^3`$ reaction, although some ambiguity remains regarding the molecular nature of the deuteron target. Moreover, by means of the proposed model the effect of a superstrong magnetic field on laboratory Hydrogen fusion reactions is investigated for the first time showing that despite the remarkable increase in the cross section of the $`dd`$ reaction, the $`pp`$ reaction is still too slow to justify experimentation.
ACKNOWLEDGMENTS
This work was financially supported by the Greek State Grants Foundation (IKY) under contract #135/2000. The revised version was written at ECT during a nuclear physics fellowship. The author would like to thank the director of ECT Prof.Malfliet for his kind hospitality and support.
Prof. C.Rolfs’ advice during the Bologna-2000 conference has proved invaluable to the completion of this paper.
FIGURE CAPTIONS
Figure 1. The screening (acceleration) factor $`f_{1s}`$ with respect to the relative interaction energy of two fusing protons for various superstrong magnetic fields (in units of $`10^{12}G)`$
Figure 2. The screening (acceleration) factor $`f_{1s}`$ with respect to the relative interaction energy of two fusing deuterons for various superstrong magnetic fields (in units of $`10^{12}G)`$
Figure 3. The $`H^2(d,p)H^3`$ astrophysical factor $`S\left(E\right)`$ measured in keV-b with respect to the center of mass interaction energy $`E_{cm}`$ $`\left(keV\right).`$ The data (squares) are taken from Ref. $`.`$ The solid curve represents Eq. $`\left(\text{33}\right)`$, which makes use of the asymptotic form given in Ref. . The dashed curve represents Eq. $`\left(\text{34}\right)`$ where all four parameters $`S\left(0\right),S^{^{}}\left(0\right),S^{^{\prime \prime }}\left(0\right),U_e`$ are fitted simultaneously. The dotted curve is obtained by adopting as a screening energy the value given by Eq. $`\left(\text{25}\right),`$ while the dash-dotted curve stands for the astrophysical factor obtained by using Eq. $`\left(\text{21}\right).`$
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# The status of NLL BFKL aafootnote aTalk presented at the XXXVth Rencontres de Moriond, QCD and High Energy Hadronic Interactions, Les Arcs, France, March 18-25, 2000.
## 1 Introduction and total cross sections
Understanding perturbative high-energy hadronic scattering is one of the fundamental problems of QCD, and one which has seen considerable progress in the past couple of years.
The reaction under consideration is that of the scattering of objects with some transverse scale $`Q^2`$ at a centre-of-mass energy $`\sqrt{s}`$, in the limit $`sQ^2\mathrm{\Lambda }_{\mathrm{QCD}}^2`$. This was first examined in the mid 1970’s with the resummation of the leading logarithmic (LL) terms $`(\alpha _\mathrm{s}\mathrm{ln}s)^n`$ to give the result that the cross section $`\sigma s^{4\mathrm{ln}2\alpha _\mathrm{s}N_C/\pi }`$. This is known as the BFKL pomeron.$`^\mathrm{?}`$ This result however seems to be in contradiction with numerous experimental results, which invariably indicate that while there is a rise of the cross section at large $`s`$, the power is somewhat lower, of the order of $`0.3`$ as opposed to the $`0.5`$ expected from the LL calculation (for $`\alpha _\mathrm{s}0.2`$).
It was expected that this discrepancy would be resolved by the inclusion of the next-to-leading (NLL) corrections $`\alpha _\mathrm{s}(\alpha _\mathrm{s}\mathrm{ln}s)^n`$, but it turns out that they modify the power $`\omega `$ as follows: $`^\mathrm{?}`$
$$\omega =2.65\alpha _\mathrm{s}16.3\alpha _\mathrm{s}^20.12\text{(for }\alpha _\mathrm{s}=0.2\text{)},$$
(1)
which is just as incompatible with the data as the LL result. Closer inspection reveals that the NLL corrections, taken literally, also imply negative cross sections for the scattering of two objects of different transverse sizes $`^\mathrm{?}`$ (essentially, in such a collinear limit the convergence becomes even worse).
There have been several significantly different attempts to explain the large size of the corrections and to thus estimate yet higher-order corrections with the aim of obtaining stable predictions. One involves the idea of a minimum rapidity gap between emissions,$`^\mathrm{?}`$ while another argues that a meaningful answer can be obtained using a more ‘natural’ renormalisation scheme coupled with BLM resummation.$`^\mathrm{?}`$ Both suffer from instabilities (with respect to the choice of the rapidity gap and the renormalisation scheme, respectively) and have difficulty solving the problem of the negative cross sections, essentially because they do not take into account the worsening of the convergence in the collinear limit.
An alternative approach indeed starts from a consideration of the collinear limit, and observes that a large part of the NLL corrections (even for non-collinear scattering) come precisely from terms that are enhanced in the collinear limit. These can be resummed by requiring that the cross section satisfy the renormalisation group properties in both collinear limits (there are two collinear limits corresponding to which of the scattering objects is the smaller).$`^\mathrm{?}`$ One finds that the problem of negative cross sections goes away completely and that the power of the cross section growth is relatively stable with respect changes of scheme. The result for the power as a function of $`\alpha _\mathrm{s}`$ is shown in fig. 1.
The cleanest (most inclusive) measurement of the BFKL power comes from $`\gamma ^{}\gamma ^{}`$ scattering, by the L3 collaboration $`^\mathrm{?}`$ at scales $`Q^2=3.5`$ and $`14.5`$ GeV<sup>2</sup>. They quote a value for the power of $`\omega =0.37\pm 0.04`$. The quoted error is probably an underestimate: in the fit procedure the normalisation is fixed to be the LL value, whereas one should fit also for the normalisation because it too is liable to be subject to significant (but as yet uncalculated) higher-order corrections. A more realistic error is about three times larger. There are also systematic errors associated with the functional form used to fit for the power. So the measurement is roughly in accord with the theoretical expectation which goes from $`0.29`$ to $`0.33`$ for this range of scales (for comparison the LL power is in the range $`0.6`$ to $`0.8`$), but a satisfactory comparison requires, on the experimental front, higher precision data and larger energies (e.g. at the NLC), and on the theoretical front a calculation of the expected normalisation (i.e. the impact factors to NLL order).
## 2 Scaling violations
So far we have discussed the total cross section in a context where both scattering objects are perturbative. An important situation is that in which only one of the objects is perturbative, namely deep inelastic scattering. Since the other object is non-perturbative the calculation of the total cross section is beyond perturbation theory. When the $`s`$ is of the same order as $`Q^2`$ (i.e. $`x`$ not too small) we know however that the cross section can be factorised into non-perturbative parton distributions convoluted with perturbative coefficient functions, and the scaling violations of the parton distributions can be calculated perturbatively through the DGLAP equations.$`^\mathrm{?}`$ The intuitive justification for this picture is that emissions are ordered in transverse scale, so that the non-perturbative parton distribution contains emissions only below a given scale $`Q_0^2`$, while the coefficient function (and parton distribution evolution) involves only subsequent emissions, which are all above $`Q_0^2`$.
However in high-energy scattering (i.e. at small $`x`$) it has long been known that emissions are not ordered in transverse momentum. This led to the suggestion $`^\mathrm{?}`$ that the factorisation which holds at moderate $`x`$ might break down at small $`x`$. (Recently, on this basis, it has been argued that small-$`x`$ splitting functions are beyond calculation and can at best be fitted.$`^\mathrm{?}`$)
Last summer an explicit counterexample was presented,$`^\mathrm{?}`$ namely a toy model which retained the fundamental property of BFKL evolution, namely emissions which are not ordered in transverse scale, but in which the property of factorisation could be demonstrated for all $`x`$. Recently, numerical methods have been developed $`^\mathrm{?}`$ which allow the study of factorisation in the full BFKL equation (for the time being at LL, with running coupling). In this approach one solves the BFKL equation to obtain a gluon distribution as a function of $`x`$ and $`Q^2`$. This distribution depends intrinsically on the regularisation of $`\alpha _\mathrm{s}`$ in the infra-red (IR) and the initial transverse scale. One then performs a deconvolution of the scaling violations from the gluon distribution, to obtain an effective splitting function.
The results of this procedure are shown in fig. 2, which shows the effective splitting function for three different sets of infra-red regularisations. One sees that the splitting function is independent of the IR regularisation, except at very small $`x`$, where a non-perturbative component takes over. More detailed investigation reveals that this second component is higher-twist, and despite its much stronger $`x`$ dependence, the scaling violations themselves can be predicted to within a relative higher-twist term which is not $`x`$ enhanced, thus confirming that factorisation holds even at small $`x`$.
A rough explanation of why factorisation holds is the following: the fact that the evolution is not ordered in transverse momentum means that the separation into a parton distribution and coefficient or splitting functions is less trivial. If one chooses to separate the two at some scale $`Q_0`$, then the parton distribution now depends on emissions both below and above $`Q_0^2`$. But significantly, the splitting and coefficient functions turn out to still only depend on emissions above $`Q_0^2`$.
This brings us to the question of what exactly the small-$`x`$ splitting functions look like. A fundamental property is that the power-growth only sets in at relatively small $`x`$. It is also significantly weaker than what would be observed in total cross sections at the same scale. This is the reason for the presence of two powers in fig. 1: $`\omega _s`$ is that which applies to total cross sections, while $`\omega _c`$ is that which is relevant for splitting functions. These two points (which have been noted also by Thorne $`^\mathrm{?}`$) mean that for the $`x`$ values that are relevant at HERA, the small-$`x`$ resummed splitting functions are probably quite similar to the DGLAP splitting functions, explaining the unexpected success of the latter in reproducing the observed scaling violations.
## Acknowledgements
I wish to thank Marcello Ciafaloni and Dimitri Colferai (in collaboration with whom many of the results presented here have been obtained) for discussions. This work was supported in part by the E.U. QCDNET contract FMRX-CT98-0194.
## References
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# Manifestation of anisotropy persistence in the hierarchies of MHD scaling exponents
## Abstract
The first example of a turbulent system where the failure of the hypothesis of small-scale isotropy restoration is detectable both in the ‘flattening’ of the inertial-range scaling exponent hierarchy, and in the behavior of odd-order dimensionless ratios, e.g., skewness and hyperskewness, is presented. Specifically, within the kinematic approximation in magnetohydrodynamical turbulence, we show that for compressible flows, the isotropic contribution to the scaling of magnetic correlation functions and the first anisotropic ones may become practically indistinguishable. Moreover, skewness factor now diverges as the Péclet number goes to infinity, a further indication of small-scale anisotropy.
A wide interest has recently been devoted to the possible occurrence of small-scale isotropy restoration for scalar (see e.g., Ref. and references therein), Navier-Stokes and magnetohydrodynamical (MHD) turbulence . The scenario can be summarized as follows. In the presence of anisotropic large-scale injection mechanisms, the inertial-range statistics is characterized by an infinite hierarchy of scaling exponents; however the leading contribution to scaling comes from the isotropic component. From this point of view, one might argue that large-scale anisotropy does not affect inertial-range scaling properties. Actually, focussing on a larger set of observables, small-scale anisotropies become manifest. It turns out that the behavior of odd-order dimensionless ratios (e.g., skewness and hyperskewness) is completely different from the case of small-scale isotropy restoration. Such indicators go to zero down to the inertial range much slower than predicted by dimensional considerations or, more dramatically, they diverge at the smallest scales (see also ).
The main aim of this Letter is to present a model of MHD turbulence where, by varying the degree of compressibility of the velocity field, anisotropic persistence is now detectable both from the ‘flattening’ of the hierarchy of inertial-range scaling exponents (the isotropic component and the first anisotropic ones may become practically indistinguishable) and the divergence of skewness factor with the Péclet number. We give the basic ideas and results; longer and more exhaustive technical discussions will be presented in a forthcoming paper.
In the presence of a mean component $`𝑩^o`$ (actually varying on a very large scale $`L`$) and for the compressible velocity field $`𝒗`$, the kinematic MHD equations describing the evolution of the fluctuating (divergence-free) part $`𝑩`$ of the magnetic field are :
$`_tB_\alpha +𝒗\mathbf{}B_\alpha =(B_\alpha +B_\alpha ^o)\mathbf{}𝒗+𝑩\mathbf{}v_\alpha +`$ (1)
$`𝑩^o\mathbf{}v_\alpha +\kappa _0^2B_\alpha ,\alpha =1,\mathrm{},d,`$ (2)
where $`\kappa _0`$ is the magnetic diffusivity. The field $`𝑩^o`$ plays the same role as an external forcing driving the system and being also a source of anisotropy for the magnetic field statistics.
Our choice for the velocity statistics generalizes that of the well-known kinematic Kazantsev–Kraichnan model for the compressible case: $`𝒗`$ is a Gaussian process of zero average, homogeneous, isotropic and white in time. It is self-similar and defined by the two-point correlation function:
$$v_\alpha (t,𝒙)v_\beta (t^{},𝒙^{})=\delta (tt^{})[d_{\alpha \beta }^0S_{\alpha \beta }(𝒙𝒙^{})],$$
(3)
where $`d_{\alpha \beta }^0=\mathrm{const}\delta _{\alpha \beta }`$ and $`S_{\alpha \beta }(𝒙𝒙^{})`$ is fixed by isotropy and scaling:
$$S_{\alpha \beta }(𝒓)=r^\xi \left[𝒳\delta _{\alpha \beta }+𝒴\frac{r_\alpha r_\beta }{r^2}\right],r\left|𝒙𝒙^{}\right|,$$
(4)
with the coefficients
$$𝒳=\frac{𝒮^2(d+\xi 1)\xi 𝒞^2}{(d+\xi )(d1)\xi },𝒴=\frac{d𝒞^2𝒮^2}{(d+\xi )(d1)}.$$
(5)
The degree of compressibility is thus controlled by the ratio $`\mathrm{}𝒞^2/𝒮^2`$, with $`𝒮^2(_iv_k_iv_k)`$ and $`𝒞^2(\mathbf{}𝒗)^2`$. It satisfies the inequality $`0\mathrm{}1`$; $`\mathrm{}=0`$ and 1 corresponding to the purely solenoidal and potential velocity fields, respectively.
In the present Letter our attention will be focused on the inertial-range behavior of magnetic correlation functions, where power laws are expected in their decompositions on a set of orthonormal functions $`P`$
$$B_{}^n(t,𝒙)B_{}^q(t,𝒙^{})=_{j=0}^{\mathrm{}}P_j(\mathrm{cos}\varphi )r^{\zeta _j^{n,q}},$$
(6)
$`B_{}`$ being some component of $`𝑩`$, e.g., its projection along the direction $`\widehat{𝒓}𝒓/r`$ or $`\widehat{𝑩}^o𝑩^o/B^o`$, and $`\varphi `$ is the angle between $`𝒓`$ and $`𝑩^o`$. Notice that, due to the anisotropic injection mechanism, the inertial-range statistics is now characterized by an infinite hierarchy of exponents ($`j`$ denotes the $`j`$-th anisotropic sector), rather than just one exponent as in the isotropic case.
If the Kolmogorov 1941 isotropization hypothesis holds, the contribution from the anisotropic sectors (i.e. $`j0`$ ) to the scaling of correlation functions should be negligible with respect to the isotropic component. Such picture indeed holds for even correlation functions and solenoidal velocity , but it breaks down for compressibility strong enough.
In order to prove this fact, consider first the pair correlation function, $`C_{\alpha \beta }(t,𝒓)B_\alpha (t,𝒙)B_\beta (t,𝒙^{})`$. Here, exploiting the zero-mode technique , the complete set of scaling exponents $`\zeta _j^{1,1}`$ can be found nonperturbatively. We give the basic ideas of the strategy.
The first key point is that exact, closed equation for $`C_{\alpha \beta }`$ can be found due to the time decorrelation of the velocity field:
$`_tC_{\alpha \beta }=S_{ij}_i_jC_{\alpha \beta }\left(_jS_{i\beta }\right)_iC_{\alpha j}`$ (7)
$`\left(_jS_{\alpha i}\right)\left(_iC_{j\beta }\right)+\left(_i_jS_{\alpha \beta }\right)\left(C_{ij}+B_i^oB_j^o\right)+`$ (8)
$`2\kappa _0^2C_{\alpha \beta }B_\beta ^oB_j^o_i_jS_{\alpha i}B_\alpha ^oB_j^o_i_jS_{\beta i}+`$ (9)
$`B_\alpha ^oB_\beta ^o_i_jS_{ij}C_{\alpha j}(_i_jS_{\beta i})+C_{\beta j}_i_jS_{\alpha \beta }+`$ (10)
$`C_{\alpha \beta }_i_jS_{ij}+2(_iC_{\alpha \beta })_jS_{ij}`$ (11)
(see for the derivation in the incompressible case).
$`C_{\alpha \beta }`$ can be projected on a vector basis including the large scale magnetic field $`𝑩^o`$ , or, alternatively, on the purely inertial range basis that span the irreducible representations of SO($`d`$) . Deep in the inertial range, both descriptions yield a system of linear algebraic equations for the scaling law coefficients. The leading solutions are associated with the homogeneous solutions of such system, i.e. with zero modes. Their scaling exponents follow from the imposition that the determinant of the coefficients be zero. Schematically, the solutions can be expressed in the form:
$$\zeta _j^{1,1}=\alpha +\sqrt{\beta +\gamma \sqrt{\delta }}(j\text{even}),$$
(12)
where $`\alpha `$, $`\beta `$, $`\gamma `$ and $`\delta `$ are cumbersome functions of $`\xi `$, $`d`$ and $`j`$ and will not be reported here for the sake of brevity. Contributions to the scaling associated with odd $`j`$’s vanish due to the symmetry $`\varphi \varphi `$.
In particular, the expression derived in Ref. in the isotropic case has been recovered here for $`j=0`$.
For $`j=0`$ and $`j=2`$, the limit $`\xi 0`$ in Eq. (12) yields:
$`\zeta _0^{1,1}=(1+2\mathrm{}d\mathrm{})\xi +O(\xi ^2),`$ (13)
$`\zeta _2^{1,1}={\displaystyle \frac{2\mathrm{}[4+d(d2)(d+1)]}{(d1)(d+2)}}\xi +O(\xi ^2),`$ (14)
while for large $`d`$ we have:
$`\zeta _0^{1,1}={\displaystyle \frac{\mathrm{}\xi }{1+\mathrm{}\xi }}d{\displaystyle \frac{12\mathrm{}(1\xi )}{1+\mathrm{}\xi }}\xi +O(1/d),`$ (15)
$`\zeta _2^{1,1}={\displaystyle \frac{\mathrm{}\xi }{1+\mathrm{}\xi }}d{\displaystyle \frac{\mathrm{}\xi (\xi 2)}{1+\mathrm{}\xi }}+O(1/d).`$ (16)
In Fig. 1 we present the behavior of $`\zeta _j^{1,1}`$ obtained from the expression (12) for $`\mathrm{}=0`$ (thin lines) and $`\mathrm{}=1`$ (heavy lines). As one can see, for $`\mathrm{}=0`$ we have $`\zeta _0^{1,1}<0`$ but $`\zeta _2^{1,1}>0`$ and $`\zeta _4^{1,1}>0`$ so that, in particular, $`(r/L)^{\zeta _0^{1,1}}(r/L)^{\zeta _2^{1,1}}`$ in the inertial range of scales.
The situation changes for $`\mathrm{}>\mathrm{}_c`$ ($`\mathrm{}_c0.1274`$ for $`j=2`$ and $`d=3`$), when $`\zeta _2^{1,1}`$ becomes negative for all $`0\xi 1`$. This means that, for compressibility strong enough, $`(r/L)^{\zeta _0^{1,1}}`$ and $`(r/L)^{\zeta _2^{1,1}}`$ become very close (still, $`\zeta _j^{1,1}<\zeta _k^{1,1}`$ for $`j<k`$). As we can see from Eqs. (15), (16) the effect becomes dramatic for large dimensions and $`\mathrm{}0`$ when $`\zeta _0^{1,1}\zeta _2^{1,1}`$. The contribution to the scaling of the correlation functions in Eq. (6) coming from the sector $`j=2`$ thus becomes less and less subleading as $`\mathrm{}`$ and/or $`d`$ increase.
Further strong evidence of the crucial role of compressibility in the failure of the small-scale isotropy restoration can be obtained by looking at the higher-order correlation functions \[i.e. $`n`$ and/or $`q>1`$ in Eq. (6)\] and, in particular, at dimensionless ratios of odd-order moments. Were isotropy restored at small scales, such ratios would go to zero for large Péclet number. The latter is defined as Pe $`(L/\eta )^{1/\xi }`$, where $`L`$ is the integral scale and $`\eta \kappa _0^{1/\xi }`$ is the dissipation scale.
Non-zero values of such indicators are thus the signature of anisotropy persistence. As we are going to show, in contrast to the incompressible case, for $`\mathrm{}`$ large enough the skewness factor now diverges as $`\text{Pe}\mathrm{}`$. To be more specific, the leading contribution in the expression (6) can be written as:
$$B_{}^n(t,𝒙)B_{}^q(t,𝒙^{})(r/\eta )^{\alpha _0^{n,q}}(r/L)^{\beta _0^{n,q}}r^{\zeta _0^{n,q}},$$
(17)
where we have defined $`\zeta _0^{n,q}\alpha _0^{n,q}+\beta _0^{n,q}`$. The expressions for $`\zeta _0^{n,q}`$ have been obtained up to first order in $`\xi `$ by means of the field theoretic renormalization group and operator product expansion. A detailed presentation of these techniques for the case $`\mathrm{}=0`$ can be found in (see also Refs. for the scalar case and for general review); below we confine ourselves to only the necessary information.
The key role is played by the scaling dimensions $`\mathrm{\Delta }[n,j]`$ of the $`j`$-th rank composite operators $`B_{\alpha _1}(x)\mathrm{}B_{\alpha _j}(x)\left[B_\alpha (x)B_\alpha (x)\right]^l`$ ($`n2l+j`$ is the total number of $`𝑩`$’s). The first order in $`\xi `$ (one-loop approximation) yields
$$\mathrm{\Delta }[n,j]=\frac{\xi }{2(d1)(d+2)}\left\{n(n1)\left[2\mathrm{}(d^3d^22d+4)\right](nj)(n+j+d2)(d+12\mathrm{})\right\}+O(\xi ^2).$$
(18)
The exponents in Eq. (6) are related to the dimensions (18) as follows (see Ref. ):
$$\zeta _j^{n,q}=\mathrm{\Delta }[n+q,j]\mathrm{\Delta }[n,j_n]\mathrm{\Delta }[q,j_q],$$
where $`j_n=0`$ for $`n`$ even and 1 for $`n`$ odd. We thus conclude that the inequality $`\mathrm{\Delta }[n,j]/j0`$, which follows from Eq. (18) for all $`\mathrm{}`$ and $`d2`$, generalizes the hierarchy discussed above to the higher-order functions. It becomes flatter and flatter as $`\mathrm{}`$ grows ($`^2\mathrm{\Delta }[n,j]/j\mathrm{}0`$), while for $`d\mathrm{}`$ (and $`\mathrm{}0`$) the effect becomes even stronger: in the leading $`O(d)`$, expressions (18) and (19) are now independent of $`j`$.
The leading term of Eq. (6) in the inertial range is given by the contribution with the minimal $`j`$:
$$\zeta _0^{n,q}=\{\begin{array}{cc}\xi \frac{nq(1+\mathrm{}d^22\mathrm{})+(d+12\mathrm{})}{(d+2)}\hfill & \\ \frac{\xi nq}{(d+2)}(1+\mathrm{}d^22\mathrm{}),\hfill & \end{array}$$
(19)
where the first holds if both $`n`$ and $`q`$ are odd, and the second otherwise. For $`n=q=1`$, expression (13) is recovered. Knowing the exponents $`\alpha _0^{n,q}`$ and $`\beta _0^{n,q}`$, dimensionless ratios of the form $`R_{2n+1}(r)B_{||}^{2n}(𝒙)B_{||}(𝒙^{})/B_{||}(𝒙)B_{||}(𝒙^{})^{(2n+1)/2}`$ can be constructed and, as in Ref. , evaluated at the dissipative scale \[i.e. $`r=\eta `$ in Eq. (17)\]. When doing this, the explicit dependence on Pe appears and the final $`O(\xi )`$ expressions read:
$$R_{2n+1}(\text{Pe})\text{Pe}^{\sigma _{2n+1}},n=1,2,\mathrm{},$$
(20)
with the exponents
$$\sigma _{2n+1}(\mathrm{})\frac{2n^2(12\mathrm{}+\mathrm{}d^2)}{(d+2)}\frac{(12\mathrm{}+\mathrm{}d)}{2}.$$
(21)
It is easy to verify from expression (21) that for $`d2`$ and $`n1`$ we have $`\sigma _{2n+1}(\mathrm{})/\mathrm{}>0`$. Negative values of $`\sigma _{2n+1}(0)`$ may thus become positive due to $`\mathrm{}`$. In particular, for $`d=3`$ we obtain $`\sigma _3(\mathrm{})=(23\mathrm{}1)/10`$, which becomes positive for $`\mathrm{}>1/23`$. It then follows that $`R_3\mathrm{}`$ as $`\text{Pe}\mathrm{}`$, the footprint of the persistence of the small-scale anisotropy.
Some remarks on the limit of large space dimensions are worth. One immediately realizes from Eq. (19) that, for $`d\mathrm{}`$ the scaling exponents reduce to $`\zeta _0^{n,q}=\xi nqd\mathrm{}`$, that means $`\zeta _0^{n,q}=n\zeta _0^{1,q}=q\zeta _0^{n,1}`$, and thus the vanishing of intermittency for $`d\mathrm{}`$. Nevertheless, in this limit $`\zeta _0^{1,1}=\zeta _2^{1,1}`$ and, at the dissipative scale $`\eta `$, $`\sigma _{2n+1}(\mathrm{})=\mathrm{}d(2n^21/2)>0`$, two clear signatures of small-scale anisotropy persistence. The latter may thus occur also in the absence of intermittency.
Let us now examine the possible mechanisms at the origin of the small-scale anisotropy persistence in our problem. They can be easily grasped in two dimensions; our previous results being valid for all $`d2`$, it is reasonable to expect that the mechanisms we are going to show hold also for higher $`d`$’s. Let us start from the incompressible case assuming, without loss of generality, that $`\widehat{𝑩}^o`$ is oriented along the $`y`$-axis (i.e. $`\widehat{𝑩}^o𝒆_y`$). As we shall see, all our considerations will hold, a fortiori, in the compressible case.
It can be shown that the magnetic field can be represented by the (scalar) magnetic flux function in the form $`𝑩=\mathbf{}\psi \times 𝒆_z`$, where $`\psi `$ satisfies the passive scalar equation forced by a large scale gradient $`𝑮𝑩^o\times 𝒆_z`$. One finds
$$_t\mathrm{\Psi }+𝒗\mathbf{}\mathrm{\Psi }=\kappa _0^2\mathrm{\Psi }$$
(22)
where $`\mathrm{\Psi }(t,𝒙)\psi +𝑮𝒙`$. Notice that, fort $`d=2`$, expression (12) recovers the results of Ref. obtained directly for the scalar problem.
An interesting feature recently recognized for the passive scalar problem, for both synthetic advection by $`\delta `$-correlated velocity and Navier–Stokes advection , is related to the formation of “cliffs”, i.e. very steep scalar gradients within very short distances. The region where such gradients develop are separated by “plateaus” where the scalar depends smoothly on the position (i.e. regions of gradient expulsion).
The emergence of this ubiquitous pattern is explained by considering the action of the velocity derivative matrix $`_\alpha v_\beta `$. As emphasized in Refs. , scalar gradients are weak in the elliptic regions of the velocity field (in two dimensions they correspond to purely imaginary eigenvalues of the velocity derivative matrix) where its ‘rotational’ character inhibits the formation of strong scalar gradients. On the contrary, in the hyperbolic regions the flow is almost ‘irrotational’: the alignment of scalar gradients with the direction of the eigenvector corresponding to the most negative eigenvalue of the velocity derivative matrix (roughly, the direction corresponding to compression along one direction) is not discouraged, and actually observed , and strong scalar gradients develop.
When the scalar is forced by isotropic injection mechanisms, no preferential direction arises and the intense gradients are randomly oriented. The final result is a small scale isotropic statistics. The situation changes in the case of non-isotropic injection mechanisms like the one encountered here. In this case strong gradients are oriented along $`𝑮`$ and, as very recently shown in Refs. , small-scale isotropy is consequently not restored for the scalar field. Exploiting the relation between the magnetic field and the magnetic flux function, we can conclude that extreme magnetic fluctuations have a tendency to occur preferentially along the direction $`\widehat{𝑩}^o`$, the origin of the observed small-scale anisotropy.
In the compressible case both eigenvalues can be negative Compression may thus occur in both directions, enhancing the formation of fronts in the magnetic flux function.
In conclusion, we presented a simple model of MHD turbulence where, by varying the degree of compressibility of the velocity field, the persistence of anisotropy is detectable both from the hierarchy of inertial range exponents and from the divergence of skewness factor with the Péclet number. Although our results were obtained on the base of a specific model, they seem to be rather general: a similar behavior is also observed for a passive scalar advected by a compressible velocity field (however, the flattening is less pronounced for the latter).
The results presented here give the first evidence of a turbulent system displaying such twofold behavior.
We are deeply grateful to L. Ts. Adzhemyan, L. Biferale, A. Celani, A. Lanotte, and M. Vergassola for discussions throughout the present work. N. V. A. acknowledges the hospitality of the University of Helsinki. N. V. A. was supported by the Academy of Finland, GRACENAS Grant No. 97-0-14.1-30, and RFFI Grant No. 99-02-16783. A. M. was partially supported by the INFM PA project No. GEPAIGG01. P. M. G. was partially supported by grant ERB4001GT962476 from EU.
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# Thermostatistics of extensive and non–extensive systems using generalized entropies
## 1 Introduction
Non–extensive systems are those for which the thermodynamic potentials do not scale linearly with the system size. As a way of example, in some electric or magnetic systems with very long–range interactions the ground state energy per particle increases with the number of particles. In the absence of other effects, such as screening, these system are “genuinely” non–extensive. If we apply to them the standard Boltzmann–Gibbs formalism of the Statistical Mechanics, we find that the internal energy, Helmholtz free energy and other thermodynamic potentials are non–extensive as well. This standard formalism can be implemented by using the definition of the entropy $`S`$ in terms of the probabilities $`p_i`$ of the $`i=1,\mathrm{},W`$ possible microscopic configurations<sup>1</sup><sup>1</sup>1We use throughout the paper dimensionless units where the Boltzmann constant, $`k_B`$ is equal to $`1`$:
$$S=\underset{i}{}p_i\mathrm{ln}p_i$$
(1)
The actual calculation of the entropy assumes a set of probabilities $`\{p_i\}`$. These are computed by finding the maximum of the above expression when some extra conditions defining an ensemble (fixed number of particles and mean energy, for example) are imposed.
Even for systems in which the energy levels do scale with the system size, it is possible, by using generalized definitions of the entropy, to obtain non–extensive thermodynamic potentials. One of the most successful generalizations is that of Tsallis which in 1988 proposed the following alternative expression for the entropy:
$$S_q=\frac{1_ip_i^q}{q1}$$
(2)
where $`q`$ is an entropic index that characterizes the degree of non–extensivity. It is possible to show that the entropy of the composed system A+B satisfies the relation:
$$S_q(A+B)=S_q(A)+S_q(B)+(1q)S_q(A)S_q(B)$$
(3)
when A and B are independent systems in the sense that $`p_{ij}(A+B)=p_i(A)p_j(B)`$. We see that for $`q1`$ there is no additivity in the entropy, which also implies non–extensivity. The Boltzmann–Gibbs entropy, Eq. (1), and extensivity are recovered in the limit $`q1`$. Since the probabilities $`\{p_i\}`$ satisfy $`p_i^q>p_i`$ for $`q<1`$ and $`p_i^q<p_i`$ for $`q>1`$, the superextensive, $`q<1`$, and the subextensive, $`q>1`$, regimes will privilege the rare and frequent events respectively.
In the last years there have been many studies in which Tsallis non–extensive statistics has been applied to different situations (See for a review). In some cases, the systems considered are genuinely non–extensive (in the sense defined above) while in others the non–extensivity arises as a result of the application of the new statistics. In fact, and due to the intrinsic non–extensivity of the Tsallis statistics, it has been argued that its natural range of applicability should include systems with long–range interactions or long–range microscopic memory processes, as well as dynamical systems in which the space–time geometry has a multifractal–like structure, because those systems are in general genuinely non–extensive. Although most of the literature (including this paper) basically derives equilibrium properties starting from the generalized definition of the entropy, it has been conjectured recently , however, that the Tsallis entropy could be relevant instead in the study of non–equilibrium processes.
Due to the difficulty of deriving exact results, it is natural to use numerical methods to obtain the properties of a system with many degrees of freedom when studied under the rules of the new statistics. This is necessary in order to extract results that could be checked against experiments. However, these studies have been hampered by the failure of the typical Monte Carlo methods to adequately generate representative equilibrium configurations distributed according to Tsallis statistics. It is the purpose of this paper to explain in detail new methods that can be used to study the equilibrium properties of a many–particle system when it is considered under generalized statistics. Although our methods are quite general, we will illustrate their use by considering a prototypical genuinely non–extensive system: the Ising ferromagnet model with long–range interactions. We will also consider the short–range Ising model in order to test the simulation methods and to compare the results obtained from the use of Tsallis statistics in extensive and non–extensive systems.
In the remaining of the section, we will outline briefly which are the basic difficulties one encounters when trying to generalize the standard Monte Carlo methods (such as the Metropolis algorithm) to the study of Tsallis statistics. The main problem is that the probabilities $`\{p_i\}`$ can not be given an explicit expression, as we will see in the following discussion. Let us consider the canonical ensemble. The probabilities $`\{p_i\}`$ in this ensemble are found by solving the maximization problem of the entropy $`S_q`$ as given by Eq.(2) subject to the constrains of (i) positivity: $`p_i0`$, (ii) normalization: $`_ip_i=1`$ and (iii) a fixed mean value for the internal energy: $`=U`$, where $``$ is the Hamiltonian of the system and the mean value of any function $`O`$ of the microscopic configurations is computed according to the general rule:
$$O=\underset{i}{}O_iu(p_i),$$
(4)
$`O_i`$ is the value of $`O`$ at the configuration whose probability is $`p_i`$ and we have introduced a function $`u(p_i)`$ that allows the definition of generalized mean values. The standard mean values are recovered by taking $`u(p_i)=p_i`$. Although, initially, the choices $`u(p_i)=p_i`$ (first option), and $`u(p_i)=p_i^q`$ (second option) were considered, later, it was shown that a better choice, in the sense that it preserves the Legendre structure of the resulting thermodynamics formalism, is to consider $`u(p_i)=p_i^q/_jp_j^q`$ (third option) . We will use the following notation of the averages in this third option:
$$O_q=\underset{i}{}O_iP_i;P_i=\frac{p_i^q}{_jp_j^q}.$$
(5)
$`\{P_i\}`$ are known as the “escorts” probabilities . It is possible to recover the configuration probabilities $`p_i`$ from the escorts probabilities using:
$$p_i=\frac{P_i^{1/q}}{_jP_j^{1/q}}.$$
(6)
The entropy, in terms of the $`\{P_i\}`$’s is given by:
$$S_q=\frac{1(_iP_i^{1/q})^q}{q1}$$
(7)
Concerning the different definitions for the averages, it should be said that it has been shown recently that the standard mean values of the first option, $`u(p_i)=p_i`$, can be also made compatible with the Legendre structure of the thermodynamics and the resulting formalism also represents a thermodynamically stable description. In this paper, we follow mainly the formulation in terms of the mean values defined by Eq.(5), although in a later section we will show that the results obtained using the standard mean values can be mapped onto the ones obtained using Eq. (5).
The maximization problem for the unknown escort probabilities $`P_i`$ in the canonical ensemble with a given internal energy $`U_q`$ is:
$`{\displaystyle \frac{\delta }{\delta P_i}}\left[S_q\beta {\displaystyle \underset{i}{}}\epsilon _iP_i\alpha {\displaystyle \underset{i}{}}P_i\right]`$ $`=`$ $`0`$ (8)
$`P_i`$ $``$ $`0`$ (9)
$`{\displaystyle \underset{i}{}}P_i`$ $`=`$ $`1`$ (10)
$`{\displaystyle \underset{i}{}}\epsilon _iP_i`$ $`=`$ $`U_q`$ (11)
where $`\alpha `$, $`\beta `$ are Lagrange multipliers. $`\{\epsilon _i\}`$ are the energy levels of the system under consideration whose ground state energy will be denoted by $`E_0`$. Solving the problem Eq.(8-11) one obtains the probabilities for the canonical ensemble as :
$$P_i=\{\begin{array}{cc}0,\hfill & \hfill 1\frac{(1q)\beta (\epsilon _iU_q)}{(_jP_j^{1/q})^q}<0\\ \frac{[1(1q)\beta (\epsilon _iU_q)/(_jP_j^{1/q})^q]^{\frac{q}{1q}}}{_k[1(1q)\beta (\epsilon _kU_q)/(_jP_j^{1/q})^q]^{\frac{q}{1q}}},\hfill & \hfill \mathrm{otherwise}\end{array}$$
(12)
The probabilities defined in this way are real and non–negative. The condition giving the possible values of $`\beta `$ and $`\epsilon _i`$ for which $`P_i0`$ in Eq. (12) is called the cut–off condition. One can show that the probabilities Eq. (12) are invariant under a change in the energy levels $`\epsilon _i\epsilon _i+\mathrm{\Delta }\epsilon `$ (and the same change in the internal energy $`U_qU_q+\mathrm{\Delta }\epsilon `$) for arbitrary $`\mathrm{\Delta }\epsilon `$. By introducing the temperature $`T=1/\beta `$, it is possible to show also the validity of the relation
$$1/T=S_q/U_q$$
(13)
which reflects the Legendre structure of the thermodynamics obtained.
Notice that Eq.(12) does not give yet the actual values of the probabilities since the $`\{P_i\}`$’s appear in a non–trivial way in both sides of the equation (either explicitly or in the cut-off condition). This is different from the solution obtained in the usual Boltzmann–Gibbs canonical ensemble (recovered in the limit $`q1`$) in which the solution adopts the explicit form:
$$P_i=\frac{\mathrm{e}^{\beta \epsilon _i}}{_j\mathrm{e}^{\beta \epsilon _j}}$$
(14)
(although, of course, it is very difficult to compute the denominator of this expression, the partition function, for an interacting system). An iterative method to solve Eq.(12) has been used in . In this method, an initial guess for the probabilities is fed in the right hand side of (12) and this equation is used recursively until convergence is achieved. We will see, however, that for many–particle systems, it might very difficult to achieve convergence in some cases.
A convenient way of writing Eq. (12) is by using an auxiliary parameter $`\beta ^{}`$ defined as :
$$\beta ^{}=\frac{\beta }{(1q)\beta _j\epsilon _jP_j+(_jP_j^{1/q})^q}$$
(15)
Defining $`T^{}1/\beta ^{}`$, and using Eqs. (7) and (11) one can rewrite this equation as
$$T=\frac{T^{}(1q)U_q}{1+(1q)S_q}$$
(16)
In terms of $`\beta ^{}`$ the solution adopts a form similar to that of the standard canonical ensemble:
$$P_i=\{\begin{array}{cc}0,\hfill & \hfill 1(1q)\beta ^{}\epsilon _i<0\\ \frac{[1(1q)\beta ^{}\epsilon _i]^{\frac{q}{1q}}}{_j[1(1q)\beta ^{}\epsilon _j]^{\frac{q}{1q}}},\hfill & \hfill \mathrm{otherwise}\end{array}$$
(17)
One can then adopt the following practical procedure : choose a value for the parameter $`T^{}`$ and compute the probabilities $`P_i`$ as a function of $`T^{}`$ using Eq.(17). Compute the internal energy $`U_q`$, the entropy $`S_q`$ and the temperature $`T`$, always as a function of $`T^{}`$, using Eqs. (11), (7) and (16), respectively. Finally, vary $`T^{}`$ in order to make the parametric plots $`U_q(T)`$ and $`S_q(T)`$. Other thermodynamic potentials follow the usual definition. For instance, the Helmholtz free energy is $`F_q=U_qTS_q`$.
It is important to realize that although the probabilities $`P_i`$, when considered as a function of $`T`$, do not depend on an arbitrary shift $`\mathrm{\Delta }\epsilon `$ of the energy levels or, in other words, do not depend on the zero of energy, $`E_0`$, they do depend on $`E_0`$ when considered as a function of $`T^{}`$. This means that the averages as a function of $`T^{}`$ can not be physically relevant because they depend on the zero of energy. This is why $`T^{}`$ has to be interpreted only as an auxiliary parameter, not as an actual temperature. Of course, the relation $`T^{}T`$ depends also on the zero of energy in such a way that both dependences cancel and the averages as a function of $`T`$ are independent of a shift in the energy levels. An interesting question is to determine the range of values for the parameter $`T^{}`$ that should be used in order to obtain the usual range $`0T<\mathrm{}`$ for the temperature $`T`$. Using that, according to the definition (2), it is $`1+(1q)S_q>0`$ we obtain from Eq.(16) that $`T^{}`$ should vary in the range $`[(1q)E_0,\mathrm{})`$, where we have used that the energy at zero temperature is the ground state energy $`U_q(T=0)=E_0`$. Therefore, it is important to use the right range of values for $`T^{}`$ in order to reach all the possible values for $`T`$. In particular, $`T^{}`$ might need to take negative values either for $`q<1`$ or for $`q>1`$ unless one adopts $`E_0=0`$ as we will throughout this paper. To our understanding, it is not clear in the literature the fact that the averages as a function of the parameter $`T^{}`$ depend on the zero of energy and that it might be necessary to consider negative values for $`T^{}`$ in order to span the whole range of values for $`T`$.
As stated before, the main problem to perform Tsallis thermostatistics simulations at a given temperature $`T`$ is that there is not an explicit expression for the probabilities $`P_i`$, c.f. see Eq. (12). The practical procedure outlined above (compute $`P_i`$ as a function of $`T^{}`$ and then compute $`U_q`$, $`S_q`$ and $`T`$ as a function of $`T^{}`$ in order to make parametric plots by varying $`T^{}`$) is not straightforward to implement numerically since it is very difficult to use Eq.(7) to compute the entropy. This is because the usual Monte Carlo methods, i.e. the Metropolis algorithm, require only the probabilities $`P_i`$ up to a normalization factor where Eq.(7) requires the absolute, normalized values of $`P_i`$. It is the object of this paper to explain in detail some numerical methods of the Monte Carlo type that can be used to perform the necessary averages for generalized statistics, including Tsallis.
There have been previous attempts to perform numerical simulations of Tsallis statistics using Monte Carlo methods. An earlier work in this direction is that of T. Penna et. al. who extended the Metropolis acceptance procedure to include a dependence in the $`q`$ parameter. However, this method does not satisfy the detailed balance condition which is a key ingredient of Monte Carlo methods. Another interesting approach is that of I. Andricioaei et. al. , who performed a Metropolis Monte Carlo algorithm which does satisfy the detailed balance condition for the probability $`P_i`$ as a function of the parameter $`T^{}`$ but, since they do not make the temperature transformation $`T^{}T`$, they are unable to determine the actual temperature $`T`$ of the simulation. All these works considered the second version, $`u(p_i)=p_i^q`$, for the definition of averages, which, as discussed before, has proven afterwards not to be the optimal election . We have used also a similar sampling in the context of Simulated Annealing . A recent approach proposed by A. R. Lima et. al. uses the broad histogram Monte Carlo method, which determines the number of microstates using a balance equation between near neighbor energy levels. They are able to apply this method to the Ising Model with short–range interactions. This is a valid Monte Carlo simulation with full control of the temperature $`T`$ but its applicability is somewhat restricted. As we will show, the Ising model with long–range interactions can not be treated straightforwardly with this method because the spin flip dynamics does not produce transitions between near energy levels and, consequently, the broad histogram Monte Carlo method, in its present form, can not be used to study long–range interacting systems. We are able to overcome these difficulties by extending a method which had been developed some time ago to compute directly the number of states with a given energy and which does not depend on the definition of the entropy. In fact, our method can be used to study any statistics and applications will be shown both for the Boltzmann–Gibbs and the Tsallis statistics. We develop yet a second method which is devised specifically for the Tsallis statistics and that has the advantage that it uses the familiar Metropolis algorithm plus a numerical integration.
This paper is organized in the following form: in section 2 we use a simple and limited enumeration procedure valid only for small system sizes. However, this method is exact and can be used to check against some of the approximated methods we will introduce later. In this section we also introduce the short–range Ising model (SRIM) and the long–range Ising Model (LRIM). These two models will be employed in this paper in order to test the numerical methods described here and to compare the behavior of the non–extensive Tsallis thermostatistics in genuinely extensive and non–extensive systems. In section 3 we explain in some detail the Histogram by Overlapping Windows (HOW) method. We show that this method has a wide range of applicability since it can be used both for short–range and long–range systems as well as for any kind of statistics. Section 4 presents a Metropolis Monte Carlo type method specially devised to study systems in the Tsallis statistics. In section 5 we present some results concerning the validity of some scaling relations for the SRIM and LRIM in the Tsallis thermostatistics context. These scaling relations are an extension of the ones holding for the Boltzmann–Gibbs statistics . In section 6 we present results for these models using standard mean values instead of those defined by Eq. (5). In section 7 we discuss the results of using Tsallis statistics in the microcanonical ensemble. Finally, in section 8 we summarize the main conclusions of this work.
## 2 Exact enumeration
The problem to compute the probabilities $`\{P_i\}`$ using Eq.(17) is that the number of terms in the sum of the denominator of this equation, the number of microstate configurations $`W`$, is extremely large (typically scales exponentially with the system size). However, for small systems, it might be possible to enumerate completely the microstates and, therefore, to compute magnitudes of interest such as the internal energy $`U_q(T)`$, the entropy $`S_q(T)`$, the free energy $`F_q(T)`$, etc. We follow this approach in this section. Although the system sizes one can usually study with this method are very far from reaching a situation in which scaling laws with system size apply, we use the results as a bench-test in order to compare with other approximate methods that will be introduced in the following sections.
We will consider Ising type models with Hamiltonian:
$$=\underset{(i,j)}{}J_{ij}(1s_is_j)$$
(18)
where each of the $`N`$ spin variables $`s_i`$, $`i=1,\mathrm{},N`$ can take the values $`\pm 1`$. The sum $`_{(i,j)}`$ runs over all distinct pairs of sites on a $`d`$-dimensional regular lattice of linear size $`L=N^{1/d}`$ with periodic boundary conditions and lattice constant equal to $`1`$. $`J_{ij}`$ is the coupling parameter between spins $`i`$ and $`j`$. Note that for ferromagnetic couplings, $`J_{ij}0`$, the ground state is double degenerate and its energy is $`E_0=0`$. The usual, nearest–neighbors, or short–range Ising model (SRIM), is obtained taking $`J_{ij}=1`$, if $`r_{ij}=1`$ and $`J_{ij}=0`$, if $`r_{ij}>1`$. The long–range Ising model (LRIM) is defined by using
$$J_{ij}=1/r_{ij}^\alpha ,$$
(19)
where $`r_{ij}`$ is the distance between the spins $`i`$ and $`j`$, and the parameter $`\alpha `$ sets the interaction range. The SRIM is formally recovered by taking the limit $`\alpha \mathrm{}`$. Depending on the value of $`\alpha `$ and the space dimension $`d`$, the LRIM has two regimes: the extensive regime, $`\alpha >d`$, and the non–extensive regime, $`\alpha d`$. This can be seen by roughly estimating the mean energy per spin in an infinite system as $`_1^{\mathrm{}}𝑑rr^{d1}r^\alpha `$. We obtain a convergent integral if $`\alpha >d`$ (extensive behavior), and for $`\alpha d`$ the integral diverges (non–extensive). More precisely, a convenient scale for the mean energy per spin in a finite system of size $`N`$ is given by :
$$\stackrel{~}{N}=1+d_1^L𝑑rr^{d1}r^\alpha =\frac{N^{1\alpha /d}\alpha /d}{1\alpha /d}$$
(20)
The definition of $`\stackrel{~}{N}`$ is such that the limit $`\alpha d`$ is well defined. Again, we see that for $`\alpha >d`$ the internal energy per spin scales as a constant in the limit of large $`N`$, but for $`\alpha ad`$, it grows with the system size. The system is, in this latter case, genuinely non–extensive. The SRIM limit, $`\alpha \mathrm{}`$, gives the expected result $`\stackrel{~}{N}=1`$.
The number of configurations in the Ising model is $`W=2^N`$. We have made a complete enumeration of the $`i=1,\mathrm{},W`$ configurations and their energies $`\epsilon _i`$ for a linear chain, $`d=1`$, of sizes up to $`N=34`$. We have used these results to compute the probabilities $`P_i`$ and then the thermodynamic magnitudes of interest using Tsallis statistics with the third option for the averages. In Fig. 1, we plot the exact internal energy $`U_q`$ as a function of the temperature $`T`$, for several values of the parameter $`q`$ for the LRIM in a genuinely non–extensive situation, $`\alpha =0.8`$, Fig. 1.b, and for the SRIM, Fig. 1.a. In the $`q<1`$ case we observe that there is a range of temperatures for which the internal energy is not a single–valued function of the temperature and, for a given value of $`T`$, there are several possible values for $`U_q`$. This ambiguity is resolved by using a Maxwell–like construction that replaces the loop in the energy curves by a vertical straight line connecting the two points with the same free energy $`F_q(T)`$.
We stress that the loops in the energy curves appear as a result of the $`T^{}T`$ transformation and, therefore, will not be observed when plotting the energy as a function of $`T^{}`$. Typical $`T^{}T`$ transformations are shown in Fig. 2. We observe in this figure that for $`q<1`$ a same value of $`T^{}`$ can correspond, in some cases, to three or more values of $`T`$ giving rise to the observed multivalued behavior in the energy. Figure 2 helps us to understand the failure of the iterative method that has been proposed to solve the set of equations (12). Although each value of $`T^{}`$ defines a unique set $`\{P_i\}`$. We see Fig. 2 that for $`q>1`$ there are some intervals of $`T^{}`$ where the transformation $`T^{}T`$ is almost horizontal. Therefore, one value of $`T`$ corresponds nearly to a complete interval for $`T^{}`$ and hence, there are many possible solutions for $`\{P_i\}`$ very close to the real one. This is the main reason for the failure of iterative methods for $`q>1`$. The situation worsens for increasing system size $`N`$.
Whatever illuminating the method of exact enumeration is, its validity is limited to very small values for $`N`$. To our knowledge, the largest value ever considered in an exact enumeration scheme for a short–range Ising model is $`N=4^3=64`$ . Simulations at larger $`N`$ sizes require other methods, as the ones presented in the next two sections.
## 3 The Energy Histogram Method using Overlapping Windows
Although the number of possible microscopic configurations is in general very large, the range of possible energy values usually takes a much smaller value. For instance, for the one–dimensional SRIM introduced in the previous section, with $`N`$ spins we have $`W=2^N`$, but the number of possible energy values is $`N/2+1`$. Let $`M`$, in general, be the number of possible energy levels. We denote by $`\mathrm{\Omega }(E_k)`$ the number of microscopic states sharing the same energy $`E_k`$, $`k=0,\mathrm{},M1`$. Obviously $`_k\mathrm{\Omega }(E_k)=W`$. We rewrite all the sums in Eqs. (17,7,5) as:
$`P(E_k)`$ $`=`$ $`\{\begin{array}{cc}0,\hfill & \hfill 1(1q)\beta ^{}E_k<0\\ \frac{[1(1q)\beta ^{}E_k]^{\frac{q}{1q}}}{_n\mathrm{\Omega }(E_n)[1(1q)\beta ^{}E_n]^{\frac{q}{1q}}},\hfill & \hfill \mathrm{otherwise}\end{array}`$ (23)
$`S_q`$ $`=`$ $`{\displaystyle \frac{1(_k\mathrm{\Omega }(E_k)P(E_k)^{1/q})^q}{q1}}`$ (24)
$`O_q`$ $`=`$ $`{\displaystyle \underset{k}{}}\mathrm{\Omega }(E_k)O(E_k)P(E_k)`$ (25)
where the sums run over the $`M`$ energy levels.
Notice that, once the $`\mathrm{\Omega }(E_k)`$’s have been computed, any statistics can be performed upon the system. Whether we use Tsallis, Boltzmann–Gibbs or any other generalized statistics is simply a trivial change in the computational scheme. Moreover, it is also trivial to compute the averages for any value of the parameters, say $`T`$ or $`q`$. Therefore, although the calculation of the $`\mathrm{\Omega }_k`$’s might be time consuming, the pay–off is tremendous<sup>2</sup><sup>2</sup>2All the simulations in this paper have been performed using a Pentium-III processor at 550MHz.
In general, the $`\mathrm{\Omega }(E_k)`$ are very difficult to obtain exactly. An important exception that will be used throughout this paper is the SRIM in 1-d, for which the energy levels are given by $`E_k=4k`$ for $`k=0,\mathrm{},N/2`$ and whose degeneracy is:
$$\mathrm{\Omega }(E_k)=2\left(\genfrac{}{}{0pt}{}{N}{E_k/2}\right)$$
(26)
In the cases in which $`\mathrm{\Omega }(E_k)`$ is not known we need approximate numerical methods. The most naive way to find the $`\mathrm{\Omega }(E_k)`$’s is to generate different system configurations randomly and count how many times a configuration with energy $`E_k`$ appears. However, this approach fails because the complete set of $`\mathrm{\Omega }(E_k)`$ values span too many orders of magnitude. In general, two energy levels $`E_k`$ and $`E_n`$ could differ as much as $`\mathrm{\Omega }_k/\mathrm{\Omega }_n\mathrm{exp}N`$. This means that the range of variation of $`\mathrm{\Omega }(E_k)`$ over the $`M`$ different energy levels is very large and it is not possible to generate in a single run a histogram that covers all the energy levels, unless one generates a number of configurations of the order of the total number available to the system, $`W`$.
The Histogram by Overlapping Windows method (HOW) used here avoids this problem by generating system configurations within a restricted energy interval and estimating the relative weights of these energy levels from the number of times they appear in each interval. By generating enough intervals spanning the whole energy range, one is able to obtain good quality estimators of the numbers $`\mathrm{\Omega }(E_k)`$. An earlier account of the method has been given in and we explain now in some detail how the method works.
Let us consider first the SRIM in arbitrary dimension. In this case, the possible energy values are $`E_k=4k`$ for $`k=0,\mathrm{},dN/2`$. Following the original work , we consider the intervals (windows) $`\{E_0,E_1,E_2,E_3\}`$, $`\{E_3,E_4,E_5,E_6\}`$, $`\{E_6,E_7,E_8,E_9\}`$, etc. Each window consists of $`4`$ consecutive energy levels and the last energy value of one window is the first of the next one. The next step is to take one of the intervals and to generate configurations whose energy belongs to it. This is achieved, after preparing the system initially with one of the energies of the interval, by flipping spins chosen at random. A spin flip is accepted only if it leaves the system in one of the energy levels of the interval and it is rejected otherwise. The ratio of the number of generated states with energy $`E_k`$ to the number of generated states with energy $`E_n`$ is an unbiased estimator of $`\mathrm{\Omega }(E_k)/\mathrm{\Omega }(E_n)`$, for $`E_k`$ and $`E_n`$ within the energy window. The quality of the estimator increases with the number of generated configurations. From the overlap between windows one can compute $`\mathrm{\Omega }(E_k)`$ for the whole range of energies. The number of energy values in each window ($`4`$ in the previous example) is not important as far as it is not too large (such that the ratios $`\mathrm{\Omega }(E_k)/\mathrm{\Omega }(E_n)`$ are not extremely small or large) and it is not too small either. If the window is very small, most spin flips will yield an energy outside the range of allowed values and the number of accepted, i.e. independent, configurations will be very small. Moreover, the final algorithm must be ergodic: any energy value in a window should be obtained from any other value in the same window after a sufficient number of spin flips.
The same basic idea has been used in other short–range Hamiltonians . We are concerned now with the extension of this method to consider long–range interactions, in particular the LRIM introduced before. A modification needed is that the energy values $`E_k`$ will represent now a continuum set of energies with a bin size $`\delta E`$. The energy levels are then $`E_k=k\delta E`$ $`k=1,\mathrm{},M`$ and $`\mathrm{\Omega }(E_k)`$ counts all the configurations $`i`$ whose energy $`\epsilon _i`$ satisfies $`E_k\epsilon _i<E_k+\delta E`$. In other words, one makes the approximation of considering that all the energies lying in one bin count as one single level. This turns out not to be a critical point, although one has to check that the results are independent, within the simulation errors, of the magnitude of $`\delta E`$.
A more important point concerns the optimal size $`l`$ of the energy window $`\{E_k,E_{k+1},\mathrm{},E_{k+l}\}`$. Since, for a long–range interacting system, a single spin flip can produce a very large change in the energy, it is important not to choose $`l`$ too small. To make this point clear, we have calculated exactly the number of states $`\mathrm{\Omega }(E_k)`$ for a system with $`N=34`$ by using the complete enumeration of the $`W=2^{34}`$ possibles configurations, see Fig. 3. Using these exact results we study the energy changes that a single spin flip makes both in the SRIM and the LRIM cases. A typical situation is shown in Fig.4. In this figure we plot the histogram of the (exact) number of configurations using a value for the bin $`\delta E=4`$ for the SRIM, and $`\delta E=1`$ for the LRIM. We select several configurations belonging to one of the energy bins (marked black in the figure) and we dash all the levels that are obtained from these configurations by flipping one spin. We see that, as expected, the change in energy for the SRIM brings the system to one of its neighboring energy levels. However, for the LRIM the energy changes are very large and, in fact, the nearest–neighbor energy levels can not be reached by using the spin flip dynamics . A measure of the typical change in energy obtained when flipping one spin is estimated by considering the ferromagnetic ground state configuration with all the spins pointing in the same direction. One spin flip in this configuration produces a change in energy $`\mathrm{\Delta }E=2_{i=1}^Nr_{ij}^\alpha `$. The equivalent number of energy bins is $`\xi =\mathrm{\Delta }E/\delta E`$. We finally take the size of the energy windows $`l=3\xi `$. In order to make sure that ergodicity is satisfied, we adopt a large overlap of size $`2\xi `$ between the windows. This means that a window goes from $`E_k`$ to $`E_{k+3\xi }`$, but the next window goes from $`E_{k+\xi }`$ to $`E_{k+4\xi }`$ and so on. For the window $`\{E_k,E_{k+1},\mathrm{},E_{k+3\xi }\}`$ only those values in the interval $`\{E_{k+\xi },\mathrm{},E_{k+2\xi }\}`$ are considered for the evaluation of the ratios $`\mathrm{\Omega }(E_k)/\mathrm{\Omega }(E_n)`$. To summarize, for the LRIM, it is necessary that both the window size and the overlap between windows has the correct size, depending on $`\alpha `$. For $`\alpha =0.8`$, used in our simulations, we take $`\xi =4`$ independently of the system size and adjust $`\delta E`$ accordingly. We have checked that $`\xi =10`$ gives the same results within the numerical errors. The number of configurations necessary increases with the required accuracy. We have adopted in our simulations the criterion that the minimum number of counts for any energy bin within a window is $`100`$. The knowledge of the exact degeneration of the ground state $`\mathrm{\Omega }(E_0)=2`$ allows finally the calculation of all the $`\mathrm{\Omega }(E_k)`$.
In Fig. 3, we plot the number of states computed either exactly or by using the HOW method, for $`N=34`$ both for the SRIM and the LRIM. At the resolution of the figure, the exact results and the approximate ones are indistinguishable. This serves as a test that the HOW method, as implemented here, is capable or reproducing accurately the number of states in a known case. In Figs. 5, we show the number of states for sizes $`N=34,100,200,400,1000`$ as computed using the HOW method. The inserts of these figures show that this function scales as:
$$\mathrm{\Omega }(E_k)=\mathrm{}^{N\varphi (E_k/N\stackrel{~}{N})}$$
(27)
This is valid both for the SRIM and for the LRIM, if we recall that $`\stackrel{~}{N}=1`$ for the SRIM. The scaling function $`\varphi (x)`$ is different for the SRIM and the LRIM. For the SRIM, the result (26) leads to:
$$\varphi _{SRIM}(x)=\frac{x}{2}\mathrm{ln}\left(\frac{2}{x}1\right)\mathrm{ln}\left(1\frac{x}{2}\right)$$
(28)
No analytical expression is available for the LRIM. Fig. 6 compares the result for the internal energy $`U_q`$ obtained using the HOW method and the exact results obtained from the exact enumeration algorithm in the case $`N=34`$. Again, we can see that differences are too small to show up in this plot. Finally, in Fig. 7 we make a similar plot for larger values of $`N`$ obtained in this case by application of the HOW method.
## 4 The Monte Carlo method
We have mentioned already that the usual Monte Carlo algorithms of the Metropolis type can not be applied to Tsallis statistics, because the probabilities $`\{P_i\}`$ are not known as a function of the temperature $`T`$. However, it is possible to use them to compute averages as a function of the parameter $`T^{}`$ because the probabilities are known as a function of $`T^{}`$ except for a normalizing factor, which is irrelevant in the Monte Carlo methods. Those averages can be performed by using the Metropolis algorithm to generate configurations distributed according to the probabilities $`\{P_i\}`$ as follows: consider the configuration $`i`$ with energy $`\epsilon _i`$, flip a spin chosen at random to produce configuration $`j`$ with energy $`\epsilon _j`$, accept this change with a probability $`\mathrm{min}(1,P_j/P_i)`$. For the Boltzmann–Gibbs statistics the acceptance probability is the celebrated factor $`\mathrm{min}[1,\mathrm{exp}(\beta (\epsilon _j\epsilon _i))]`$. For Tsallis statistics, Eq. (17), it is instead:
$$P(ij)=\{\begin{array}{cc}0,\hfill & \hfill 1(1q)\beta ^{}\epsilon _j<0\\ \mathrm{min}[1,\frac{1(1q)\beta ^{}\epsilon _j}{1(1q)\beta ^{}\epsilon _i}]^{\frac{q}{1q}},\hfill & \hfill \mathrm{otherwise}\end{array}$$
(29)
This acceptance probability was used for the first time by I. Andricioaei et. al. . After generating $`𝒩`$ configurations, the averages $`O_q`$ at fixed $`T^{}`$ are obtained as the sum: $`O_q=𝒩^1_{s=1}^𝒩O(s)`$, where $`O(s)`$ is the value of the observable $`O`$ at the configuration $`s`$ in the sequence of configurations generated by the Monte Carlo method.
Still, we need to perform the $`T^{}T`$ transformation, in order to plot averages with respect to $`T`$, using Eq. (16). In this equation, we can use the Monte Carlo averages for the internal energy $`U_q=`$, but the entropy is yet unknown. In order to compute the entropy, we combine Eq. (16) and Eq. (13):
$$\frac{S_q}{U_q}=\frac{1+(1q)S_q}{T^{}(1q)U_q}$$
(30)
which can be integrated between arbitrary points $`A`$ and $`B`$:
$$\frac{1}{1q}\mathrm{ln}[1+(1q)S_q]|_A^B=\underset{U_q(A)}{\overset{U_q(B)}{}}\frac{dU_q}{T^{}(1q)U_q}$$
(31)
to obtain:
$$S_q(B)=\frac{[1+(1q)S_q(A)]\mathrm{}^{\left[(1q)\underset{U_q(A)}{\overset{U_q(B)}{}}\frac{dU_q}{T^{}(1q)U_q}\right]}1}{1q}$$
(32)
Finally, we need to know the value of the entropy, $`S_q(A)`$ at the initial integration point A. This depends on the particular system considered, but usually the extreme temperature cases are known. For the Ising models (both long–range and short–range), Eq. (18), the limits $`T0`$ and $`T\mathrm{}`$ are:
* $`S_q(T=0)=\frac{2^{(1q)}}{1q}`$.
* $`S_q(T=\mathrm{})=\frac{2^{N(1q)}}{1q}`$.
We have implemented this Monte Carlo method using a system size $`N=34`$. Fig. 8 shows that the function that needs to be integrated in order to perform the $`T^{}T`$ transformation is a smooth one. Fig. 9 compares the internal energy obtained by this Monte Carlo method with the exact results obtained by the exact enumeration procedure showing the validity of this Monte Carlo scheme.
The main disadvantage of this Monte Carlo method is that one might need to simulate a large range of values of $`T^{}`$ to be able to perform accurately the integration needed for the $`T^{}T`$ transformation. However, the use of extrapolation techniques, such as the multiple histogram reweighting , which permit to extend the results of a simulation at a value of the temperature to a continuum range of temperatures, allows to reduce dramatically the number of simulation points needed .
## 5 The scaling functions
For extensive systems, the internal energy per particle is just a function of the temperature, $`U(N,T)/N=u(T)`$. Clearly, by definition, this scaling relation does not hold for non–extensive systems and there has been some recent interest in finding the correct scaling laws that apply to non–extensive systems. A first result was to obtain the scaling laws that follow from the application of Boltzmann–Gibbs statistics to a genuinely non–extensive system such as the LRIM in the regime $`\alpha <d`$. The results for the internal energy, $`U`$, the magnetization $`M`$ (defined as $`M=|_{i=1}^Ns_i|`$ and computed using Eq. (25)), the Helmholtz free energy $`F`$ and the entropy $`S`$ can be summarized by the following relations :
$`U(N,T)`$ $`=`$ $`N\stackrel{~}{N}u(T/\stackrel{~}{N})`$ (33)
$`M(N,T)`$ $`=`$ $`Nm(T/\stackrel{~}{N})`$ (34)
$`F(N,T)`$ $`=`$ $`N\stackrel{~}{N}f(T/\stackrel{~}{N})`$ (35)
$`S(N,T)`$ $`=`$ $`Ns(T/\stackrel{~}{N})`$ (36)
where $`m`$,$`u`$,$`f`$,$`s`$ are the scaling functions. The argument justifying these scaling laws can be summarized as follows: the internal energy and the entropy appear in the definition of the free energy as $`F=UTS`$, therefore one expects that $`U`$ and $`TS`$ should have the same behavior for large $`N`$. Since $`U`$ scales as $`N\stackrel{~}{N}`$ and $`S`$ scales as $`N`$ one obtains that $`T`$ must scale as $`\stackrel{~}{N}`$ thus leading to the previous scaling ansatzs. Note that the SRIM case is recovered from the LRIM case in the limit $`\alpha \mathrm{}`$, when $`\stackrel{~}{N}1`$ and the scaling relations, Eq. (33-36), become the standard ones for extensive systems.
We present now the extensions of these scaling laws in the case that the models are considered under the rules of Tsallis entropy . In the case $`q1`$, the entropy is no longer an extensive quantity (this is true both for the SRIM and the LRIM). In order to generalize the argument of the previous paragraph giving the correct scale factor for the temperature, we derive from Eq. (3) the following general expression for the entropy $`A_q(N)`$ of a set of $`N`$ independent particles:
$$A_q(N)=S_q(N,T=\mathrm{})=\frac{[1+(1q)S_q(1)]^N1}{1q}$$
(37)
here $`S_q(1)`$ is the one particle entropy. In the Ising model case, one particle can be in any of the two states with probability $`1/2`$, yielding: $`S_q(1)=[12(1/2)^q]/(q1)=[2^{1q}1]/(1q)`$. After replacing in Eq. (37), we obtain:
$$A_q(N)=\frac{2^{N(1q)}1}{1q}$$
(38)
Assuming that Tsallis entropy will be scaled generically with $`A_q(N)`$, we now assume that $`U`$ and $`TS`$ scale in the same way as $`N\stackrel{~}{N}`$. Since $`TS/N\stackrel{~}{N}=[TA_q(N)/N\stackrel{~}{N}][S/A_q(N)]`$ we conjecture that the temperature has to be scaled with $`N^{}N\stackrel{~}{N}/A_q(N)`$. However, it turns out that the numerical results do not support this expression for the rescaling factors in the case $`q>1`$. Therefore, we write the scaling relations in the following more general form:
$`U_q(N,T)`$ $`=`$ $`N\stackrel{~}{N}u_q(T/N_U^{})`$ (39)
$`M_q(N,T)`$ $`=`$ $`Nm_q(T/N_U^{})`$ (40)
$`F_q(N,T)`$ $`=`$ $`N\stackrel{~}{N}f_q(T/N^{})`$ (41)
$`S_q(N,T)`$ $`=`$ $`A_q(N)s_q(T/N_S^{})`$ (42)
The previous argument, valid in the case $`q1`$, implies simply $`N_U^{}=N_S^{}=N^{}`$. For consistency in the notation, we define $`A_q^U(N)`$ and $`A_q^S(N)`$ by means of $`N_U^{}N\stackrel{~}{N}/A_q^U(N)`$ and $`N_S^{}N\stackrel{~}{N}/A_q^S(N)`$, respectively. Notice that for $`q=1`$ it is $`A_1(N)N`$ and the scaling laws (33-36) are recovered. In order to obtain a good scaling description in the case $`q>1`$ it is seen numerically that one needs to assume the limits $`A_q^U(N)2^{N(1q)}/(q1)`$ and $`A_q^S(N)2^{N(q1)}/(q1)`$. A unifying expression that reduces to the necessary ones for large $`N`$ and for all values of $`q`$ is:
$$A_q^S(N)=\frac{2^{N(1q)}}{q1},A_q^U(N)=\frac{A_q(N)^2}{A_q^S(N)}$$
(43)
Although we lack a satisfactory explanation for these relations, we note that similar scaling factors have been used previously to plot in the same scale curves for the specific heat in infinite–range Ising models and non–interacting ideal paramagnet .
In order to check the validity of these scaling relations, we have used the HOW method to simulate the one-dimensional SRIM and LRIM with $`\alpha =0.8`$, for system sizes $`N=34,100,200,400,800,1000`$, and several values for the non–extensive parameter $`q[0.1,1.9]`$. We test the proposed scaling relations, Eq. (39-42), by plotting the scaled results in Figs. 10,11,12. One can observe in the Fig. 10 that, for the same value of $`q`$, the collapse of curves of different sizes $`N`$. This is similar to what has been observed in the $`q=1`$ case .
It is more remarkable the fact that, with the previous choice for the scaling factors, all the $`q<1`$ data collapse in a single curve. The same thing occurs for the $`q>1`$ curves. Therefore, data can be described by just three universal scaling functions, corresponding to $`q<1`$, $`q=1`$, and $`q>1`$ regimes respectively (See Figs. 11,12).
One can see in Figs. 11,12 that the collapse in the entropy curves for $`q>1`$ is very poor. This is easily understood by noticing that the low temperature limit of the entropy for infinite system size is $`S_q(T=0)=(12^{1q})/(q1)`$ whereas the high temperature limit is $`S_q(T\mathrm{})=1/(q1)`$ and those two finite values can not be rescaled simultaneously. This is different of what happens for the internal energy and the magnetization for which the limits $`T0`$ and $`T\mathrm{}`$ coincide for different values of $`q`$. Finally, the scaling for the free energy follows directly from its definition $`F_q=U_qTS_q`$. For $`q1`$ it is $`f_q(x)=u_q(x)xs_q(x)`$, whereas for $`q>1`$ and in the limit of large $`N`$, the scaling function is given simply by $`f_q(x)=u_q(0)xs_q(\mathrm{})=x`$.
In summary, the scaling laws given by Eqs. (39-42) work for all values of $`q`$ when using the scaling factors given by Eqs.(43). Moreover, the scaling functions $`u_q`$, $`m_q`$ and $`f_q`$ adopt only three different forms for each magnitude corresponding to $`q>1`$, $`q=1`$ and $`q<1`$, both in the SRIM and LRIM cases.
## 6 Thermostatistics using standard mean values
It has been shown recently that the use of the standard rule for the calculation of the mean values in Tsallis statistics provides also a valid thermodynamical formalism. By “standard” rule we mean the use of the first option for the averages in which $`u(p_i)=p_i`$ is used in (4). Moreover, it has been argued that the results of using this first option coincide with the results of the third option (the one used up to here in this paper) with a trivial change in the parameter $`q1/q`$. In this section we show that it is possible indeed to map the results of one option into the results of the other, although the relation between them implies, besides the previous change in the parameter $`q`$, a non–trivial mapping for the temperature. Numerical results using the techniques developed in the previous sections, will allow us to plot the relation between the temperatures of the two options.
For the sake of clarity in the exposition we will use the subindexes “1” and “3” to denote the results one obtains in each option. Hence, the first option seeks the maximization of
$$S_1(q)=\frac{1_ip_i^q}{q1}$$
(44)
subject to the canonical ensemble constrains: $`p_i>0`$, $`_ip_i=1`$, $`_i\epsilon _ip_i=U`$. The third option, on the other hand, seeks the maximization of
$$S_3(q)=\frac{1(_ip_i^{1/q})^q}{q1}$$
(45)
subject to the same constrains. Of course, in the third option, the probabilities $`p_i`$ should be interpreted as “escort” probabilities, but this interpretation has no practical consequence whatsoever in the calculation of the averages. The key point now is that both entropies are related by:
$$S_1(1/q)=G_q[S_3(q)]$$
(46)
Where $`G_q(x)=\frac{q}{1q}[1(1+(1q)x)^{1/q}]`$ is a monotonically increasing function of $`x`$. This function satisfies the property: $`G_q^1(x)=G_{1/q}(x)`$, see Fig.(13). Hence, the same set of probabilities $`\{p_i\}`$ that maximize $`S_3(q)`$ for a given value of $`U`$ will maximize $`S_1(1/q)`$, for the same value for $`U`$. However, the fact that the probabilities coincide in both options does not mean that the averages computed using these probabilities coincide when they are plotted as a function of the temperature because it turns out that there is a non–trivial relation between the temperatures of both options. Let us denote by $`T_1`$ and $`T_3`$, respectively, the temperatures of the 1st and 3rd options. They can be defined as the (inverse of the) Lagrange multiplier needed to satisfy the constraint of fixed mean energy or, alternatively, they have been shown to satisfy the relations :
$$1/T_1(q)=\frac{S_1(q)}{U}1/T_3(q)=\frac{S_3(q)}{U}$$
(47)
Using (46) we find the desired relation between the two temperatures:
$$T_3(q)=T_1(1/q)/G_{1/q}^{}[S_1(1/q)]$$
(48)
$`G_q^{}(x)`$ is the derivative of $`G_q(x)`$. After substitution of Eq.(44), we find:
$$T_3(q)=T_1(1/q)(\underset{i}{}p_i^{1/q})^{q+1}$$
(49)
Therefore, it is true that the results of the third option at the value $`q`$ of the parameter can be obtained from those of the first one at the value $`1/q`$. However, the mapping requires a non–trivial rescaling of the temperature, as given by Eq.(49). Let us recall again that only the dependence with $`T`$ does not vary when changing the zero of energies and, hence, can be the only physically relevant one.
In order to give an alternative explanation of the relation between the temperatures of both options, let us write down the solutions for the probabilities using the $`\beta ^{}`$ parameter. For the third option, the solution is read directly from Eq.(17) that we rewrite using the notation of this section:
$$p_i=\{\begin{array}{cc}0,\hfill & \hfill 1(1q)\beta _3^{}\epsilon _i<0\\ \frac{[1(1q)\beta _3^{}\epsilon _i]^{\frac{q}{1q}}}{_j[1(1q)\beta _3^{}\epsilon _j]^{\frac{q}{1q}}},\hfill & \hfill \mathrm{otherwise}\end{array}$$
(50)
where the parameter $`\beta _3^{}=1/T_3^{}`$ is related to the temperature $`T_3`$ by Eq.(16) which reads:
$$T_3^{}(q)=T_3(q)[1+(1q)S_3(q)]+(1q)U$$
(51)
For the first option, it is possible to write the solution in the following form:
$$p_i=\{\begin{array}{cc}0,\hfill & \hfill 1(11/q)\beta _1^{}\epsilon _i<0\\ \frac{[1(11/q)\beta _1^{}\epsilon _i]^{\frac{1}{q1}}}{_j[1(11/q)\beta _1^{}\epsilon _j]^{\frac{1}{q1}}},\hfill & \hfill \mathrm{otherwise}\end{array}$$
(52)
where the parameter $`\beta _1^{}=1/T_1^{}`$ is related to the temperature $`T_1`$ by
$$T_1^{}(q)=T_1(q)[1+(1q)S_1(q)]+(11/q)U$$
(53)
Now, it is clear by comparing Eqs. (50) and (52) that the probabilities $`p_i`$ of the third option computed at $`q`$ are equal to the probabilities $`p_i`$ of the first option computed at $`1/q`$ provided that we choose the same values for the primed temperatures, $`T_3^{}(q)=T_1^{}(1/q)`$. After substitution of (51) and (53) and using (45), (44) we recover Eq.(49)
It is straightforward now to use the number of states $`\mathrm{\Omega }(E_k)`$ obtained using the HOW method to compute Eqs.(52,53,44) by replacing the sums over the configurations to sums over energy levels weighted by $`\mathrm{\Omega }(E_k)`$. In this way, we can perform the necessary averages implied in the first option as well as the temperature transformation factor needed in Eq. (49). In Fig. 14 we plot the internal energy $`U_q`$ as a function of the temperature using the standard averages of the first option. The most noticeable difference with the results of the third option, see figure (5) is that it is not necessary now to use the Maxwell construction because there are no loops with the temperature. In Fig. 15 we plot $`T_3(q)`$ vs. $`T_1(1/q)`$ in the LRIM and SRIM cases. Using these two results, it is possible to obtain the averages within the third option as a function of $`T_3`$. Of course, the results agree perfectly with those shown in Fig. 7. It is possible also to obtain from Eqs.(39-42) the scaling relations valid when the standard calculation of mean values is used for the calculation of thermodynamic quantitities.
## 7 Microcanonical Ensemble
As mentioned in the introduction, the third option can be formulated by using the entropic form Eq. (7) plus the standard rule for the calculation of mean values, Eq. (5) or, alternatively, by using the original entropic form Eq. (2), but with a mean value definition:
$$O_q=\underset{i}{}O_i\frac{p_i^q}{_jp_j^q}.$$
(54)
These two points of view are completely equivalent. The first option, as explained in the previous section, uses also the original entropic form but with standard mean values. We will consider in this section Tsallis original entropic form $`S(q)`$ in the context of the microcanonical formalism. The aim is to be able to derive the internal energy without any a priori assumption about the definition of averages. In the microcanonical ensemble we consider the maximization problem for the original entropic form $`S_q`$ given by Eq.(2) with the constrain of given energy $`E`$. The solution is the equiprobability,
$$p_i=\{\begin{array}{cc}\mathrm{\Omega }(E)^1,\hfill & \hfill \epsilon _i=E\\ 0,\hfill & \hfill \mathrm{otherwise}\end{array}$$
(55)
Where $`\mathrm{\Omega }(E)`$ is the number of configurations with energy $`E`$. The entropy as a function of the energy is:
$$S_q(E)=\frac{\mathrm{\Omega }(E)^{1q}1}{1q}$$
(56)
The temperature is defined by the thermodynamic relation Eq. (13), $`\frac{1}{T}=\frac{S_q}{E}`$ or
$$\frac{1}{T}=\mathrm{\Omega }(E)^q\frac{\mathrm{\Omega }}{E}$$
(57)
Inverting this relation, we obtain the energy as a function of the temperature, $`E(T)`$. In general, this relation needs to be inverted numerically. In terms of the scaling function $`\varphi (x)`$ defined in (27) we have:
$$T=\frac{\stackrel{~}{N}\mathrm{}^{(q1)N\varphi (E/N\stackrel{~}{N})}}{\varphi ^{}(E/N\stackrel{~}{N})}$$
(58)
where $`\varphi (x)`$ is known exactly for the SRIM and can be evaluated numerically using the HOW method for the LRIM (from the plot in the insert of Fig. 5b). Results are shown in Fig 16 where we plot the internal energy coming from the application of the microcanonical ensemble to both the SRIM and the LRIM. In the same figure we have also plotted the energy coming from the canonical ensemble using non standard mean values (third option). We can see that both approaches coincide for $`q1`$, and differ for $`q>1`$ in some temperature range. The ultimate reason for not having equivalence between the two ensembles is that fluctuations of the energy in the canonical ensemble can not be neglected. We have checked that this is indeed the case by computing the energy fluctuations $`\sigma ()=\sqrt{^2^2}`$ as a function of the system size. In the Fig. (17) we see that fluctuations, normalized by the scale of energy, $`N\stackrel{~}{N}`$, do not decay to zero for increasing $`N`$ in the range of temperatures for which the microcanonical and canonical ensemble do not agree. For $`q1`$ fluctuations do decay to zero with the system size in all the temperature range.
If we compare the microcanonical and the canonical ensemble using the standard mean values (the first option studied in the last section) we observe the coincidence of both ensembles for $`q1`$, and disagreement (in a given temperature range) for $`q<1`$. This turns out to be also consistent with non–vanishing energy fluctuations in the appropriate range. This is the expected result because of the mapping $`q1/q`$ implied in going from the third option to the first one.
## 8 Conclusions
In this paper, we have given details of two methods that can be used to perform numerical simulations for many particle systems that are governed by generalized statistics, such as the Tsallis one. The first method extends the Histogram by Overlapping Windows method, devised originally for short range Hamiltonians, to systems with very–long range interactions. The second method, devised specifically for the Tsallis thermostatistics, uses a typical Metropolis Monte Carlo updating scheme combined with a numerical integration. We have emphasized the need of using the right temperature definition if averages are to be independent of the zero of energy. We have applied our methods to the case of the Ising model with either short range (SRIM) or long range (LRIM) interactions. The latter case corresponds to a situation genuinely non–extensive in which the energy levels scale as $`N^\chi `$ with $`\chi >1`$, $`N`$ being the number of variables. We have compared the methods with some exact results available in the case of one–dimensional short range Ising model of arbitrary size and the long–range model for small system sizes.
We have shown that the internal energy, entropy, Helmholtz free–energy and magnetization follow non trivial scaling laws with the temperature $`T`$ and the number of variables $`N`$. We have justified these scaling laws by some heuristic arguments that, however, fail to reproduce the observed behavior for $`q>1`$. These scaling laws for $`q1`$ are non–extensive in the sense that the different thermodynamic potentials have to be scaled with a factor that depends in a non–trivial, i.e. non linear, way of $`N`$. The scaling laws hold for both the LRIM and the SRIM (with different scaling functions in each case), independently of the fact that the systems are genuinely non–extensive or extensive. This shows that the non–extensivity arises mainly because of the application of the Tsallis statistics.
We have discussed the differences between the use of standard (first option) and non–standard (third option) mean value definitions for the Tsallis Thermostatistics formalism. We show that, although the results of both definitions can be mapped onto each other by using the $`q1/q`$ transformations, this mapping requires as well a non–trivial change in the temperature. Finally, we have shown that the use of the microcanonical ensemble coincides with the results of the canonical ensemble in the third option only for $`q1`$. We interpret this result as the non–vanishing energy fluctuations that occur in the corresponding case.
An obvious extension of the results presented here is to consider the Ising models in spatial dimension greater than one, . The HOW method can be extended in any dimension for short range and long range interactions. For the SRIM in $`d=2`$ the exact results should be used.
We remark that the present work concerns equilibrium systems and there is no time dependence in our simulations. However, it has been recently conjectured that Tsallis statistics appears in some non–equilibrium systems such as the relaxation of the non–neutral plasma experiments in . To study these non–equilibrium systems within the Tsallis statistics formalism, it would be more appropriate to use Molecular Dynamics (MD) methods in which the evolution equations are solved as a function of time. We are currently working on a MD simulation valid for a Lennard–Jones system within the Tsallis statistics. This MD method uses a Kusnezov, Bulgac and Bauer thermostat where additionally the actual temperature has to be calculated using a relation similar to the Eq. (32).
Acknowledgments
We wish to thank A.R. Plastino for several discussions about the Tsallis statistics and for useful suggestions. We acknowledge financial support from DGES, grants PB94-1167 and PB97-0141-C02-01.
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# 1 Introduction
## 1 Introduction
The developments in the last two years have shown that the presence of a constant background $`B`$-field for open strings or D-branes lead to the noncommutativity of space-time coordinates (,,, ,,,). This can be equivalently realized by deforming the algebra of functions on the classical world volume. The operator product expansion for vertex operators is identified with the star (Moyal) product of functions on noncommutative spaces (,). In this respect it was shown that noncommutative U(N) Yang-Mills theory does arise in string theory.
The effective action in presence of a constant $`B`$-field background is
$$\frac{1}{4}Tr\left(F_{\mu \nu }F^{\mu \nu }\right)$$
where
$$F_{\mu \nu }=_\mu A_\nu _\nu A_\mu +iA_\mu A_\nu iA_\nu A_\mu $$
and the star product is defined by
$$f\left(x\right)g\left(x\right)=e^{\frac{i}{2}\theta ^{\mu \nu }\frac{}{\zeta ^\mu }\frac{}{\eta ^\nu }}f(x+\zeta )g(x+\eta )|_{\zeta =\eta =0}$$
This definition forces the gauge fields to become complex. Indeed the noncommutative Yang-Mills action is invariant under the gauge transformations
$$A_\mu ^g=gA_\mu g_{}^1_\mu gg_{}^1$$
where $`g_{}^1`$is the inverse of $`g`$ with respect to the star product:
$$gg_{}^1=g_{}^1g=1$$
The contributions of the terms $`i\theta ^{\mu \nu }`$ in the star product forces the gauge fields to be complex. Only conditions such as $`A_\mu ^{}=A_\mu `$ could be preserved under gauge transformations provided that $`g`$ is unitary: $`g^{}g=gg^{}=1.`$ It is not possible to restrict $`A_\mu `$ to be real or imaginary to get the orthogonal or symplectic gauge groups as these properties are not preserved by the star product (,). I will address the question of how is gravity modified in the low-energy effective theory of open strings in the presence of background fields. It has been shown that the metric of the target space gets modified by contributions of the $`B`$-field and that it becomes nonsymmetric (,). If we think of gravity as resulting from local gauge invariance under Lorentz transformations in the tangent manifold, then the previous reasoning would suggest that the vielbein and spin connection both get complexified with the star product. This seems inevitable as the star product appears in the operator product expansion of the string vertex operators.
We are therefore led to investigate whether gravity in D dimensions can be constructed by gauging the unitary group $`U(1,D1)`$. In this article we shall show that this is indeed possible and that one can construct a Hermitian action which governs the dynamics of a nonsymmetric complex metric. Once this is achieved, it is straightforward to give the necessary modifications to make the action noncommutative. The plan of this paper is as follows. In section two the action for nonsymmetric gravity based on gauging the group $`U(1,D1)`$ is given and the structure of the theory studied. In section three the equations of motion are solved to make connection with the second order formalism. In section four we give the generalization to noncommutative spaces. Section five is the conclusion.
## 2 Nonsymmetric gravity by gauging U(1,D-1)
Assume that we start with the $`U(1,D1)`$ gauge fields $`\omega _{\mu b}^a`$. The $`U(1,D1)`$ group of transformations is defined as the set of matrix transformations leaving the quadratic form
$$\left(Z^a\right)^{}\eta _b^aZ^b$$
invariant, where $`Z^a`$ are $`D`$ complex fields and
$$\eta _b^a=diag(1,1,\mathrm{},1)$$
with $`D1`$ positive entries. The gauge fields $`\omega _{\mu b}^a`$ must then satisfy the condition
$$\left(\omega _{\mu b}^a\right)^{}=\eta _c^b\omega _{\mu d}^c\eta _a^d$$
The curvature associated with this gauge field is
$$R_{\mu \nu b}^a=_\mu \omega _{\nu b}^a_\nu \omega _{\mu b}^a+\omega _{\mu c}^a\omega _{\nu b}^c\omega _{\nu c}^a\omega _{\mu b}^c$$
Under gauge transformations we have
$$\stackrel{~}{\omega }_{\mu b}^a=M_c^a\omega _{\mu d}^cM_b^{1d}M_c^a_\mu M_b^{1c}$$
where the matrices $`M`$ are subject to the condition:
$$\left(M_c^a\right)^{}\eta _b^aM_d^b=\eta _d^c$$
The curvature then transforms as
$$\stackrel{~}{R}_{\mu \nu b}^a=M_c^aR_{\mu \nu d}^cM_b^{1d}$$
Next we introduce the complex vielbein $`e_\mu ^a`$ and its inverse $`e_a^\mu `$ defined by
$`e_a^\nu e_\mu ^a`$ $`=`$ $`\delta _\mu ^\nu `$
$`e_\nu ^ae_b^\nu `$ $`=`$ $`\delta _b^a`$
which transform as
$`\stackrel{~}{e}_\mu ^a`$ $`=`$ $`M_b^ae_\mu ^b`$
$`\stackrel{~}{e}_a^\mu `$ $`=`$ $`\stackrel{~}{e}_b^\mu M_a^{1b}`$
It is also useful to define the complex conjugates
$`e_{\mu a}`$ $``$ $`\left(e_\mu ^a\right)^{}`$
$`e^{\mu a}`$ $``$ $`\left(e_a^\mu \right)^{}`$
With this, it is not difficult to see that
$$e_a^\mu R_{\mu \nu b}^a\eta _c^be^{\nu c}$$
transforms to
$$e_d^\mu M_a^{1d}M_e^aR_{\mu \nu f}^eM_b^{1f}\eta _c^b\left(M_c^{1l}\right)^{}e^{\nu l}$$
and is thus $`U(1,D1)`$ invariant. It is also Hermitian
$$\left(e_a^\mu R_{\mu \nu b}^a\eta _c^be^{\nu c}\right)^{}=e_c^\nu \eta _b^c\eta _e^bR_{\mu \nu f}^e\eta _a^fe^{\mu a}=e_a^\mu R_{\mu \nu b}^a\eta _c^be^{\nu c}$$
The metric is defined by
$$g_{\mu \nu }=\left(e_\mu ^a\right)^{}\eta _b^ae_\nu ^b$$
satisfy the property
$$g_{\mu \nu }^{}=g_{\nu \mu }$$
When the metric is decomposed into its real and imaginary parts:
$$g_{\mu \nu }=G_{\mu \nu }+iB_{\mu \nu }$$
the hermiticity property then implies the symmetries
$`G_{\mu \nu }`$ $`=`$ $`G_{\nu \mu }`$
$`B_{\mu \nu }`$ $`=`$ $`B_{\nu \mu }`$
The gauge invariant Hermitian action is given by
$$I=d^Dx\sqrt{G}e_a^\mu R_{\mu \nu b}^a\eta _c^be^{\nu c}$$
This action is analogous to the first order formulation of gravity obtained by gauging the group $`SO(1,D1)`$ One goes to the second order formalism by integrating out the spin connection and substituting for it its value in terms of the vielbein. The same structure is also present here and one can solve for $`\omega _{\mu b}^a`$ in terms of the complex fields $`e_\mu ^a`$ resulting in an action that depends only on the fields $`g_{\mu \nu }.`$ It is worthwhile to stress that the above action, unlike others proposed to describe nonsymmetric gravity is unique, except for the measure, and unambiguous. Similar ideas have been proposed in the past based on gauging the groups $`O(D,D)`$ and $`GL(D)`$ , in relation to string duality, but the results obtained there are different from what is presented here. The ordering of the terms in writing the action is done in a way that generalizes to the noncommutative case.
The infinitesimal gauge transformations for $`e_\mu ^a`$ is
$$\delta e_\mu ^a=\mathrm{\Lambda }_b^ae_\mu ^b$$
which can be decomposed into real and imaginary parts by writing $`e_\mu ^a=e_{0\mu }^a+ie_{1\mu }^a,`$ and $`\mathrm{\Lambda }_b^a=\mathrm{\Lambda }_{0b}^a+i\mathrm{\Lambda }_{1b}^a`$ to give
$`\delta e_{0\mu }^a`$ $`=`$ $`\mathrm{\Lambda }_{0b}^ae_{0\mu }^b\mathrm{\Lambda }_{1b}^ae_{1\mu }^b`$
$`\delta e_{1\mu }^a`$ $`=`$ $`\mathrm{\Lambda }_{1b}^ae_{0\mu }^b+\mathrm{\Lambda }_{0b}^ae_{1\mu }^b`$
The gauge parameters satisfy the constraints $`\left(\mathrm{\Lambda }_b^a\right)^{}=\eta _c^b\mathrm{\Lambda }_d^c\eta _a^d`$ which implies the two constraints
$`\left(\mathrm{\Lambda }_{0b}^a\right)^T`$ $`=`$ $`\eta _c^b\mathrm{\Lambda }_{0d}^c\eta _a^d`$
$`\left(\mathrm{\Lambda }_{1b}^a\right)^T`$ $`=`$ $`\eta _c^b\mathrm{\Lambda }_{1d}^c\eta _a^d`$
From the gauge transformations of $`e_{0\mu }^a`$ and $`e_{1\mu }^a`$ one can easily show that the gauge parameters $`\mathrm{\Lambda }_{0b}^a`$ and $`\mathrm{\Lambda }_{1b}^a`$ can be chosen to make $`e_{0\mu a}`$ symmetric in $`\mu `$ and $`a`$ and $`e_{1\mu a}`$ antisymmetric in $`\mu `$ and $`a`$. This is equivalent to the statement that the Lagrangian should be completely expressible in terms of $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$ only, after eliminating $`\omega _{\mu b}^a`$ through its equations of motion. In reality we have
$`G_{\mu \nu }`$ $`=`$ $`e_{0\mu }^ae_{0\nu }^b\eta _{ab}+e_{1\mu }^ae_{1\nu }^b\eta _{ab}`$
$`B_{\mu \nu }`$ $`=`$ $`e_{0\mu }^ae_{1\nu }^b\eta _{ab}e_{1\mu }^ae_{0\nu }^b\eta _{ab}`$
In this special gauge, where we define $`g_{0\mu \nu }=e_{0\mu }^ae_{0\nu }^b\eta _{ab}`$ , $`g_{0\mu \nu }g_0^{\nu \lambda }=\delta _\mu ^\lambda ,`$ and use $`e_{0\mu }^a`$ to raise and lower indices we get
$`B_{\mu \nu }`$ $`=`$ $`2e_{1\mu \nu }`$
$`G_{\mu \nu }`$ $`=`$ $`g_{0\mu \nu }{\displaystyle \frac{1}{4}}B_{\mu \kappa }B_{\lambda \nu }g_0^{\kappa \lambda }`$
The last formula appears in the metric of the effective action in open string theory .
## 3 Second Order Formulation
We can express the Lagrangian in terms of $`e_\mu ^a`$ only by solving the $`\omega _{\mu b}^a`$ equations of motion
$`e_a^\mu e^{\nu b}\omega _{\nu b}^c+e_b^\nu e^{\mu c}\omega _{\nu a}^be^{\mu b}e_a^\nu \omega _{\nu b}^ce_b^\mu e^{\nu c}\omega _{\nu a}^b`$ $`=`$
$`{\displaystyle \frac{1}{\sqrt{G}}}_\nu \left(\sqrt{G}\left(e_a^\nu e^{\mu c}e_a^\mu e^{\nu c}\right)\right)`$ $``$ $`X_a^{\mu c}`$
where $`X_a^{\mu c}`$ satisfy $`\left(X_a^{\mu c}\right)^{}=X_c^{\mu a}.`$ One has to be very careful in working with a nonsymmetric metric
$`g_{\mu \nu }`$ $`=`$ $`e_\mu ^ae_{\nu a}`$
$`g^{\mu \nu }`$ $`=`$ $`e^{\mu a}e_{\nu a}`$
$`g_{\mu \nu }g^{\nu \rho }`$ $`=`$ $`\delta _\mu ^\rho `$
but $`g_{\mu \nu }g^{\mu \rho }\delta _\mu ^\rho .`$ Care also should be taken when raising and lowering indices with the metric.
Before solving the $`\omega `$ equations, we point out that the trace part of $`\omega _{\mu b}^a`$ (corresponding to the $`U(1)`$ part in $`U(D)`$) must decouple from the other gauge fields. It is thus undetermined and decouples from the Lagrangian after substituting its equation of motion. It imposes a condition on the $`e_\mu ^a`$
$$\frac{1}{\sqrt{G}}_\nu \left(\sqrt{G}\left(e_a^\nu e^{\mu a}e_a^\mu e^{\nu a}\right)\right)X_a^{\mu a}=0$$
We can therefore assume, without any loss in generality, that $`\omega _{\mu b}^a`$ is traceless $`\left(\omega _{\mu a}^a=0\right).`$
Multiplying the $`\omega `$equation with $`e_a^\kappa e_c^\rho `$ we get
$$\delta _\kappa ^\mu \omega _{\nu \rho }^\nu +\delta _\rho ^\mu \omega _{\nu \kappa }^\nu \omega _{\kappa \rho }^\mu \omega _{\rho \kappa }^\mu =X_{\rho \kappa }^\mu $$
where
$`\omega _{\mu \nu }^\rho `$ $`=`$ $`e_{\nu a}e^{\rho b}\omega _{\mu b}^a`$
$`X_{\rho \kappa }^\mu `$ $`=`$ $`e_{\rho c}e_\kappa ^aX_a^{\mu c}`$
Contracting by first setting $`\mu =\kappa `$ then $`\mu =\rho `$ we get the two equations
$`3\omega _{\nu \rho }^\nu +\omega _{\nu \rho }^\nu `$ $`=`$ $`X_{\rho \mu }^\mu `$
$`\omega _{\nu \rho }^\nu +3\omega _{\nu \rho }^\nu `$ $`=`$ $`X_{\mu \rho }^\mu `$
These could be solved to give
$`\omega _{\nu \rho }^\nu `$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(3X_{\rho \mu }^\mu X_{\mu \rho }^\mu \right)`$
$`\omega _{\nu \rho }^\nu `$ $`=`$ $`{\displaystyle \frac{1}{8}}\left(X_{\rho \mu }^\mu +3X_{\mu \rho }^\mu \right)`$
Substituting these back into the $`\omega `$equation we get
$$\omega _{\kappa \rho }^\mu +\omega _{\rho \kappa }^\mu =\frac{1}{8}\delta _\kappa ^\mu \left(3X_{\rho \mu }^\mu X_{\mu \rho }^\mu \right)+\frac{1}{8}\delta _\rho ^\mu \left(X_{\kappa \mu }^\mu +3X_{\mu \kappa }^\mu \right)X_{\rho \kappa }^\mu Y_{\rho \kappa }^\mu $$
We can rewrite this equation after contracting with $`e_{\mu c}e_\sigma ^c`$ to get
$$\omega _{\kappa \rho \sigma }+e_a^\mu e_{\mu c}e_\sigma ^c\omega _{\rho \kappa }^a=g_{\sigma \mu }Y_{\rho \kappa }^\mu Y_{\sigma \rho \kappa }$$
By writing $`\omega _{\rho \kappa }^a=\omega _{\rho \nu \kappa }e^{\nu a}`$ we finally get
$$\left(\delta _\kappa ^\alpha \delta _\rho ^\beta \delta _\sigma ^\gamma +g^{\beta \mu }g_{\sigma \mu }\delta _\rho ^\alpha \delta _\kappa ^\gamma \right)\omega _{\alpha \beta \gamma }=Y_{\sigma \rho \kappa }$$
To solve this equation we have to invert the tensor
$$M_{\kappa \rho \sigma }^{\alpha \beta \gamma }=\delta _\kappa ^\alpha \delta _\rho ^\beta \delta _\sigma ^\gamma +g^{\beta \mu }g_{\sigma \mu }\delta _\rho ^\alpha \delta _\kappa ^\gamma $$
In the conventional case when all fields are real, the metric $`g_{\mu \nu }`$ is symmetric and $`g^{\beta \mu }g_{\sigma \mu }=\delta _\sigma ^\beta `$ so that the inverse of $`M_{\kappa \rho \sigma }^{\alpha \beta \gamma }`$ is simple. In the present case, because of the nonsymmetry of $`g_{\mu \nu }`$ this is fairly complicated and could only be solved by a perturbative expansion. Writing $`g_{\mu \nu }=G_{\mu \nu }+iB_{\mu \nu }`$ and from the definition $`g^{\mu \nu }g_{\nu \rho }=\delta _\mu ^\rho `$ we get
$$g^{\mu \nu }=a^{\mu \nu }+ib^{\mu \nu }$$
where
$`a^{\mu \nu }`$ $`=`$ $`\left(G_{\mu \nu }+B_{\mu \kappa }G^{\kappa \lambda }B_{\lambda \nu }\right)^1`$
$`=`$ $`G^{\mu \nu }G^{\mu \kappa }B_{\kappa \lambda }G^{\lambda \sigma }B_{\sigma \eta }G^{\eta \nu }+O(B^4)`$
$$b^{\mu \nu }=G^{\mu \kappa }B_{\kappa \lambda }G^{\lambda \nu }+G^{\mu \kappa }B_{\kappa \lambda }G^{\lambda \sigma }B_{\sigma \tau }G^{\tau \rho }B_{\rho \eta }G^{\eta \nu }+O(B^5)$$
We have defined $`G^{\mu \nu }G_{\nu \rho }=\delta _\rho ^\mu `$. This implies that
$`g^{\mu \alpha }g_{\nu \alpha }`$ $``$ $`\delta _\nu ^\mu +L_\nu ^\mu `$
$`L_\nu ^\mu `$ $`=`$ $`iG^{\mu \rho }B_{\rho \nu }2G^{\mu \rho }B_{\rho \sigma }G^{\sigma \alpha }B_{\alpha \nu }+O(B^3)`$
The inverse of $`M_{\kappa \rho \sigma }^{\alpha \beta \gamma }`$ defined by
$$N_{\alpha \beta \gamma }^{\sigma \rho \kappa }M_{\sigma \rho \kappa }^{\alpha ^{}\beta ^{}\gamma ^{}}=\delta _\alpha ^\alpha ^{}\delta _\beta ^\beta ^{}\delta _\gamma ^\gamma ^{}$$
is evaluated to give
$`N_{\alpha \beta \gamma }^{\sigma \rho \kappa }`$ $`=`$ $`{\displaystyle \frac{1}{2}}\left(\delta _\gamma ^\sigma \delta _\beta ^\rho \delta _\alpha ^\kappa +\delta _\beta ^\sigma \delta _\alpha ^\rho \delta _\gamma ^\kappa \delta _\alpha ^\sigma \delta _\gamma ^\rho \delta _\beta ^\kappa \right)`$
$`{\displaystyle \frac{1}{4}}\left(\delta _\beta ^\kappa \delta _\alpha ^\sigma L_\gamma ^\rho +\delta _\alpha ^\kappa \delta _\gamma ^\sigma L_\beta ^\rho \delta _\gamma ^\kappa \delta _\beta ^\sigma L_\alpha ^\rho \right)`$
$`+{\displaystyle \frac{1}{4}}\left(L_\gamma ^\kappa \delta _\beta ^\sigma \delta _\alpha ^\rho +L_\beta ^\kappa \delta _\alpha ^\sigma \delta _\gamma ^\rho L_\alpha ^\kappa \delta _\gamma ^\sigma \delta _\beta ^\rho \right)`$
$`{\displaystyle \frac{1}{4}}\left(\delta _\alpha ^\kappa L_\gamma ^\sigma \delta _\beta ^\rho +\delta _\gamma ^\kappa L_\beta ^\sigma \delta _\alpha ^\rho \delta _\beta ^\kappa L_\alpha ^\sigma \delta _\gamma ^\rho \right)+O(L^2)`$
This enables us to write
$$\omega _{\alpha \beta \gamma }=N_{\alpha \beta \gamma }^{\sigma \rho \kappa }Y_{\rho \sigma \kappa }$$
and finally
$$\omega _{\mu b}^a=e^{\beta a}e_b^\gamma \omega _{\mu \beta \gamma }$$
It is clear that the leading term reproduces the Einstein-Hilbert action plus contributions proportional to $`B_{\mu \nu }`$ and higher order terms. The most difficult task is to show that the Lagrangian is completely expressible in terms of $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$ only. The other components of $`e_{0\mu }^a`$ and $`e_{1\mu }^a`$ should disappear. We have argued from the view point of gauge invariance that this must happen, but it will be nice to verify this explicitly, to leading orders. We can check that in the flat approximation for gravity with $`G_{\mu \nu }`$ taken to be $`\delta _{\mu \nu }`$, the $`B_{\mu \nu }`$ field gets the correct kinetic terms. First we write
$`e_\mu ^a`$ $`=`$ $`\delta _\mu ^a{\displaystyle \frac{i}{2}}B_{\mu a}`$
$`e_{\mu a}`$ $`=`$ $`\delta _\mu ^a+{\displaystyle \frac{i}{2}}B_{\mu a}`$
and the inverses
$`e^{\mu a}`$ $`=`$ $`\delta _\mu ^a{\displaystyle \frac{i}{2}}B_{\mu a}`$
$`e_a^\mu `$ $`=`$ $`\delta _\mu ^a+{\displaystyle \frac{i}{2}}B_{\mu a}`$
The $`\omega _{\mu a}^a`$equation implies the constraint
$$X_a^{\mu a}=_\nu \left(e_a^\mu e^{\nu a}e_a^\nu e^{\mu a}\right)=0$$
This gives the gauge fixing condition
$$^\nu B_{\mu \nu }=0$$
We then evaluate
$$X_{\rho \kappa }^\mu =\frac{i}{2}\left(_\rho B_{\kappa \mu }+_\kappa B_{\rho \mu }\right)$$
This together with the gauge condition on $`B_{\mu \nu }`$ gives
$$Y_{\rho \kappa }^\mu =\frac{i}{2}\left(_\rho B_{\kappa \mu }+_\kappa B_{\rho \mu }\right)$$
and finally
$$\omega _{\mu \nu \rho }=\frac{i}{2}\left(_\mu B_{\nu \rho }+_\nu B_{\mu \rho }\right)$$
When the $`\omega _{\mu \nu \rho }`$ is substituted back into the Lagrangian, and after integration by parts one gets
$`L`$ $`=`$ $`\omega _{\mu \nu \rho }\omega ^{\nu \rho \mu }\omega _\mu ^{\mu \rho }\omega _{\nu \rho }^\nu `$
$`=`$ $`{\displaystyle \frac{1}{4}}B_{\mu \nu }^2B^{\mu \nu }`$
This is identical to the usual expression
$$\frac{1}{12}H_{\mu \nu \rho }H^{\mu \nu \rho }$$
where
$$H_{\mu \nu \rho }=_\mu B_{\nu \rho }+_\nu B_{\rho \mu }+_\rho B_{\mu \nu }$$
We have therefore shown that in $`D`$ dimensions one must start with $`2D^2`$ real components $`e_\mu ^a`$, subject to gauge transformations with $`D^2`$ real parameters. The resulting Lagrangian depends on $`D^2`$ fields, with $`\frac{D\left(D+1\right)}{2}`$ symmetric components $`G_{\mu \nu }`$ and $`\frac{D\left(D1\right)}{2}`$ antisymmetric components $`B_{\mu \nu }.`$
## 4 Noncommutative Gravity
At this stage, and having shown that it is perfectly legitimate to formulate a theory of gravity with nonsymmetric complex metric, based on the idea of gauge invariance of the group $`U(1,D1).`$ It is not difficult to generalize the steps that led us to the action for complex gravity to spaces where coordinates do not commute, or equivalently, where the usual products are replaced with star products.
First the gauge fields are subject to the gauge transformations
$$\stackrel{~}{\omega }_{\mu b}^a=M_c^a\omega _{\mu d}^cM_b^{1d}M_c^a_\mu M_b^{1c}$$
where $`M_a^{1b}`$ is the inverse of $`M_b^a`$ with respect to the star product. The curvature is now
$$R_{\mu \nu b}^a=_\mu \omega _{\nu b}^a_\nu \omega _{\mu b}^a+\omega _{\mu c}^a\omega _{\nu b}^c\omega _{\nu c}^a\omega _{\mu b}^c$$
which transforms according to
$$\stackrel{~}{R}_{\mu \nu b}^a=M_c^aR_{\mu \nu d}^cM_b^{1d}$$
Next we introduce the vielbeins $`e_\mu ^a`$ and their inverse defined by
$`e_a^\nu e_\mu ^a`$ $`=`$ $`\delta _\mu ^\nu `$
$`e_\nu ^ae_b^\nu `$ $`=`$ $`\delta _b^a`$
which transform to
$`\stackrel{~}{e}_\mu ^a`$ $`=`$ $`M_b^ae_\mu ^b`$
$`\stackrel{~}{e}_a^\mu `$ $`=`$ $`\stackrel{~}{e}_b^\mu M_a^{1b}`$
The complex conjugates for the vielbeins are defined by
$`e_{\mu a}`$ $``$ $`\left(e_\mu ^a\right)^{}`$
$`e_{}^{\mu a}`$ $``$ $`\left(e_a^\mu \right)^{}`$
Finally we define the metric
$$g_{\mu \nu }=\left(e_\mu ^a\right)^{}\eta _b^ae_\nu ^b$$
The $`U(1,D1)`$ gauge invariant Hermitian action is
$$I=d^Dx\sqrt{G}\left(e_a^\mu R_{\mu \nu b}^a\eta _c^be_{}^{\nu c}\right)$$
This action differs from the one considered in the commutative case by higher derivatives terms proportional to $`\theta ^{\mu \nu }`$. It would be very interesting to see whether these terms could be reabsorbed by redefining the field $`B_{\mu \nu }`$, or whether the Lagrangian reduces to a function of $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$ and their derivatives only.
The connection of this action to the gravity action derived for noncommutative spaces based on spectral triples (,,) remains to be made. In order to do this one must understand the structure of Dirac operators for spaces with deformed star products.
## 5 Conclusions
We have shown that it is possible to combine the tensors $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$ into a complexified theory of gravity in $`D`$ dimensions by gauging the group $`U(1,D1)`$. The Hermitian gauge invariant action is a direct generalization of the first order formulation of gravity obtained by gauging the Lorentz group $`SO(1,D1)`$. The Lagrangian obtained is a function of the complex fields $`e_\mu ^a`$ and reduces to a function of $`G_{\mu \nu }`$ and $`B_{\mu \nu }`$ only. This action is generalizable to noncommutative spaces where coordinates do not commute, or equivalently, where the usual products are deformed to star products. It is remarkable that the presence of a constant background field in open string theory implies that the metric of the target space becomes nonsymmetric and that the tangent manifold for space-time does not have only the Lorentz symmetry but the larger $`U(1,D1)`$ symmetry. The results shown here, can be improved by computing the second order action to include higher order terms in the $`B_{\mu \nu }`$ expansion and to see if this can be put in a compact form. Similarly the computation has to be repeated in the noncommutative case to see whether the $`\theta ^{\mu \nu }`$ contributions could be simplified. It is also important to determine a link between this formulation of noncommutative gravity and the Connes formulation based on the noncommutative geometry of spectral triples. To make such connection many points have to be clarified, especially the structure of the Dirac operator for such a space. This and other points will be explored in future publication
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# Separability of Density Matrices and Conditional Information Transmission
## Abstract
We give necessary and sufficient conditions under which a density matrix acting on a two-fold tensor product space is separable. Our conditions are given in terms of quantum conditional information transmission.
Ref. proposed using quantum conditional information transmission as a measure of entanglement. In its simplest case, this measure requires one speaker and two listeners. On the other hand, the simplest case of separability of density matrices is defined for two listeners but no speaker. Thus, it is not immediately apparent how quantum conditional information transmission is related to separability. And yet, they must be closely related since they are both closely related to the phenomenon of quantum entanglement. In this paper, we present a theorem that elucidates the hidden relationship between conditional information transmission and separability. The theorem gives necessary and sufficient conditions for the separability of density matrices acting on a two-fold tensor product space. The theorem can be easily generalized to the case of $`n`$-fold tensor products.
We will use $`_{\underset{¯}{a}},_{\underset{¯}{b}},\mathrm{}`$ to represent Hilbert spaces (finite dimensional ones for simplicity), and $`_{\underset{¯}{a},\underset{¯}{b}}`$ to represent $`_{\underset{¯}{a}}_{\underset{¯}{b}}`$, the tensor product of $`_{\underset{¯}{a}}`$ and $`_{\underset{¯}{b}}`$. $`dim()`$ will stand for the dimension of the Hilbert space $``$. The set of all density matrices acting on a Hilbert space $``$ will be denoted by $`𝒟()`$. If $`\rho _{\underset{¯}{x}\underset{¯}{y}}𝒟(_{\underset{¯}{x}\underset{¯}{y}})`$, we will denote the partial traces of $`\rho _{\underset{¯}{x}\underset{¯}{y}}`$ by $`\rho _{\underset{¯}{x}}=\mathrm{tr}_{\underset{¯}{y}}\rho _{\underset{¯}{x}\underset{¯}{y}}`$ and $`\rho _{\underset{¯}{y}}=\mathrm{tr}_{\underset{¯}{x}}\rho _{\underset{¯}{x}\underset{¯}{y}}`$. For any set $`S`$, we will use $`|S|`$ to represent the number of elements in $`S`$.
We will say $`\rho 𝒟(_{\underset{¯}{x}\underset{¯}{y}})`$ is separable (or more precisely, $`\underset{¯}{x},\underset{¯}{y}`$ separable) if $`\rho `$ can be expressed as
$$\rho =\underset{e}{}w_e\rho _{\underset{¯}{x}}^e\rho _{\underset{¯}{y}}^e,$$
(1)
where the $`w_e`$, called weights, are non-negative numbers that sum to 1, and where for all $`e`$, $`\rho _{\underset{¯}{x}}^e𝒟(_{\underset{¯}{x}})`$ and $`\rho _{\underset{¯}{y}}^e𝒟(_{\underset{¯}{y}})`$. Non-entangled $`\underset{¯}{x},\underset{¯}{y}`$ states are usually defined as those which are $`\underset{¯}{x},\underset{¯}{y}`$ separable.
We will say $`\rho 𝒟(_{\underset{¯}{x}\underset{¯}{y}\underset{¯}{e}})`$ is conditionally separable (or more precisely, $`\underset{¯}{x},\underset{¯}{y}|\underset{¯}{e}`$ separable) if $`\rho `$ can be expressed as
$$\rho =\underset{e}{}w_e|ee|\rho _{\underset{¯}{x}}^e\rho _{\underset{¯}{y}}^e,$$
(2)
where the $`w_e`$, called weights, are non-negative numbers that sum to 1, the states $`|e`$ are an orthonormal basis for $`_{\underset{¯}{e}}`$, and for all $`e`$, $`\rho _{\underset{¯}{x}}^e𝒟(_{\underset{¯}{x}})`$ and $`\rho _{\underset{¯}{y}}^e𝒟(_{\underset{¯}{y}})`$.
Suppose $`A`$ is a set of random variables. For example, $`A=\{\underset{¯}{a},\underset{¯}{b}\}`$. If $`\rho 𝒟(_A)`$ and $`A^{}A`$, then we will use $`S_\rho (A^{})`$ to represent $`S(\mathrm{tr}_{AA^{}}\rho )`$, where $`S()`$ is the von Neumann entropy. For example, if $`\rho 𝒟(_{\underset{¯}{a},\underset{¯}{b}})`$, then $`S_\rho (\underset{¯}{a})=S(\mathrm{tr}_{\underset{¯}{b}}\rho )`$. If $`\rho 𝒟(_{\underset{¯}{x}\underset{¯}{y}\underset{¯}{e}})`$, we define the quantum conditional mutual information, or conditional information transmission by
$$S_\rho [(\underset{¯}{x}:\underset{¯}{y})|\underset{¯}{e}]=S_\rho (\underset{¯}{x},\underset{¯}{e})+S_\rho (\underset{¯}{y},\underset{¯}{e})S_\rho (\underset{¯}{x},\underset{¯}{y},\underset{¯}{e})S_\rho (\underset{¯}{e}).$$
(3)
The classical counterpart of this is the classical conditional mutual information, which is defined, for random variables $`\underset{¯}{x},\underset{¯}{y},\underset{¯}{e}`$ with a joint distribution $`P(x,y,e)`$, by
$$H[(\underset{¯}{x}:\underset{¯}{y})|\underset{¯}{e}]=\underset{x,y,e}{}P(x,y,e)\mathrm{log}_2\frac{P(x,y|e)}{P(x|e)P(y|e)}.$$
(4)
See Ref. for a review of classical and quantum entropy presented in the same notation used in this paper.
For any $`\rho 𝒟(_{\underset{¯}{x}\underset{¯}{y}\underset{¯}{e}})`$,
$$S_\rho [(\underset{¯}{x}:\underset{¯}{y})|\underset{¯}{e}]0.$$
(5)
This is called the strong subadditivity inequality for quantum entropy. It was first proven by Lieb-Ruskai in Ref.. More recently, it has been shown that the strong subadditivity inequality becomes an equality (i.e., “is saturated”) if and only if $`\rho `$ satisfies
$$\mathrm{log}\rho =\mathrm{log}\rho _{\underset{¯}{x}\underset{¯}{e}}+\mathrm{log}\rho _{\underset{¯}{y}\underset{¯}{e}}\mathrm{log}\rho _{\underset{¯}{e}}.$$
(6)
Classical random variables $`\underset{¯}{x},\underset{¯}{y},\underset{¯}{e}`$ with joint distribution $`P(x,y,e)`$ satisfy
$$H[(\underset{¯}{x}:\underset{¯}{y})|\underset{¯}{e}]0,$$
(7)
which is the classical counterpart of Eq.(5). This inequality is saturated if and only if
$$P(x,y|e)=P(x|e)P(y|e)$$
(8)
for all $`x,y,e`$. When Eq.(8) is true, we say $`\underset{¯}{x},\underset{¯}{y}`$ are conditionally independent. Taking the logarithm of both sides of Eq.(8) yields
$$\mathrm{log}P(x,y,e)=\mathrm{log}P(x,e)+\mathrm{log}P(y,e)\mathrm{log}P(e),$$
(9)
which is the classical counterpart of Eq.(6).
Theorem 1: $`\rho 𝒟(_{\underset{¯}{x}\underset{¯}{y}})`$ is $`\underset{¯}{x},\underset{¯}{y}`$ separable if and only if there exists a Hilbert space $`_{\underset{¯}{e}}`$ and a density matrix $`\sigma 𝒟(_{\underset{¯}{x}\underset{¯}{y}\underset{¯}{e}})`$ such that
1. $`\rho =\mathrm{tr}_{\underset{¯}{e}}\sigma `$,
2. $`S_\sigma [(\underset{¯}{x}:\underset{¯}{y})|\underset{¯}{e}]=0`$,
3. $`\sigma _{\underset{¯}{y}\underset{¯}{e}},\sigma _{\underset{¯}{x}\underset{¯}{e}}`$ and $`\sigma _{\underset{¯}{e}}`$ commute pairwise,
4. the eigenvalues of $`\sigma _{\underset{¯}{e}}`$ are are non-zero and non-degenerate.
proof:
($``$) $`\rho `$ can be expanded as in Eq.(1). We can always choose the weights $`w_e`$ of the expansion to be non-zero and non-degenerate. Indeed, if $`w_e=0`$, we just eliminate that term from the expansion. If $`e_1e_2`$ and $`w_{e_1}=w_{e_2}`$, then we replace the $`e_1`$ and $`e_2`$ terms of the expansion by
$`w_{e_1}(\rho _{\underset{¯}{x}}^{e_1}\rho _{\underset{¯}{y}}^{e_1}+\rho _{\underset{¯}{x}}^{e_2}\rho _{\underset{¯}{y}}^{e_2})=`$ (10)
$`=w_{e_1}\rho _{\underset{¯}{x}}^{e_1}\rho _{\underset{¯}{y}}^{e_1}+({\displaystyle \frac{w_{e_1}}{2}}+ϵ)\rho _{\underset{¯}{x}}^{e_2}\rho _{\underset{¯}{y}}^{e_2}+({\displaystyle \frac{w_{e_1}}{2}}ϵ)\rho _{\underset{¯}{x}}^{e_2}\rho _{\underset{¯}{y}}^{e_2}=`$
$`={\displaystyle \underset{j=1}{\overset{3}{}}}w_{e_j}^{}\rho _{\underset{¯}{x}}^{{}_{}{}^{}e_{j}^{}}\rho _{\underset{¯}{y}}^{{}_{}{}^{}e_{j}^{}},`$
where we have define a new $`e`$ value called $`e_3`$ and we have set $`w_{e_1}^{}=w_{e_1}`$, $`w_{e_2}^{}=\frac{w_{e_1}}{2}+ϵ`$ and $`w_{e_3}^{}=\frac{w_{e_1}}{2}ϵ`$. For small enough $`ϵ>0`$, we achieve our goal of representing $`\rho `$ as in Eq.(1) with weights that are non-degenerate and non-zero. If $`E`$ is the new set of $`e`$ values, let $`_{\underset{¯}{e}}`$ be a Hilbert space of dimension $`|E|`$, and let $`|e`$ for $`eE`$ be an orthonormal basis for $`_{\underset{¯}{e}}`$. Define $`\sigma 𝒟(_{\underset{¯}{x}\underset{¯}{y}\underset{¯}{e}})`$ by
$$\sigma =\underset{eE}{}w_e|ee|\rho _{\underset{¯}{x}}^e\rho _{\underset{¯}{y}}^e.$$
(11)
Thus, $`\sigma `$ is $`\underset{¯}{x},\underset{¯}{y}|\underset{¯}{e}`$ separable. Clearly, $`\rho =\mathrm{tr}_{\underset{¯}{e}}\sigma `$. In Ref., it is shown by straightforward computation that any $`\underset{¯}{x},\underset{¯}{y}|\underset{¯}{e}`$ separable density matrix $`\sigma `$ satisfies $`S_\sigma [(\underset{¯}{x}:\underset{¯}{y})|\underset{¯}{e}]=0`$. $`\sigma `$ has the following partial traces:
$$\sigma _{\underset{¯}{e}}=\underset{e}{}w_e|ee|,$$
(12)
$$\sigma _{\underset{¯}{x}\underset{¯}{e}}=\underset{e}{}w_e|ee|\rho _{\underset{¯}{x}}^e,$$
(13)
$$\sigma _{\underset{¯}{y}\underset{¯}{e}}=\underset{e}{}w_e|ee|\rho _{\underset{¯}{y}}^e.$$
(14)
Clearly, $`\sigma _{\underset{¯}{y}\underset{¯}{e}},\sigma _{\underset{¯}{x}\underset{¯}{e}}`$ and $`\sigma _{\underset{¯}{e}}`$ commute pairwise. The eigenvalues of $`\sigma _{\underset{¯}{e}}`$ are the $`w_e`$, which are non-zero and non-degenerate.
($``$) $`S_\sigma [(\underset{¯}{x}:\underset{¯}{y})|\underset{¯}{e}]=0`$ so Eq.(6) is true for $`\sigma `$. In fact, since $`\sigma _{\underset{¯}{y}\underset{¯}{e}},\sigma _{\underset{¯}{x}\underset{¯}{e}}`$ and $`\sigma _{\underset{¯}{e}}`$ commute pairwise, and $`\rho _{\underset{¯}{e}}`$ has non-zero eigenvalues, we can combine the logarithms to obtain
$$\sigma =\sigma _{\underset{¯}{y}\underset{¯}{e}}\sigma _{\underset{¯}{x}\underset{¯}{e}}(\sigma _{\underset{¯}{e}})^1.$$
(15)
Since $`\sigma _{\underset{¯}{e}}`$ is a Hermitian matrix, it can be diagonalized:
$$\sigma _{\underset{¯}{e}}=\underset{e}{}w_e|ee|,$$
(16)
where $`w_e`$ and $`|e`$ for all $`e`$ are the eigenvalues and eigenvectors of $`\sigma _{\underset{¯}{e}}`$. One has that
$$\sigma _{\underset{¯}{e}}\sigma _{\underset{¯}{x}\underset{¯}{e}}=\sigma _{\underset{¯}{x}\underset{¯}{e}}\sigma _{\underset{¯}{e}}.$$
(17)
Thus,
$$w_ee|\sigma _{\underset{¯}{x}\underset{¯}{e}}|e^{}=e|\sigma _{\underset{¯}{x}\underset{¯}{e}}|e^{}w_e^{}$$
(18)
for all $`e,e^{}`$. Since the eigenvalues $`w_e`$ of $`\sigma _{\underset{¯}{e}}`$ are non-degenerate, $`ee^{}`$ implies $`w_ew_e^{}`$, and therefore $`e|\sigma _{\underset{¯}{x}\underset{¯}{e}}|e^{}=0`$. It follows that $`\sigma _{\underset{¯}{x}\underset{¯}{e}}`$ is diagonal in its $`_{\underset{¯}{e}}`$ sector:
$$\sigma _{\underset{¯}{x}\underset{¯}{e}}=\underset{x,x^{},e}{}A_{x,x^{}}^e|e,xe,x^{}|,$$
(19)
where for all $`x,x^{},e`$, $`A_{x,x^{}}^e`$ is a complex number, and where $`|x`$ for all $`x`$ is any orthonormal basis of $`_{\underset{¯}{x}}`$. If for each $`e`$, $`\rho _{\underset{¯}{x}}^e𝒟(_{\underset{¯}{x}})`$ is defined by
$$\rho _{\underset{¯}{x}}^e=\underset{x,x^{}}{}\frac{A_{x,x^{}}^e}{w_e}|xx^{}|,$$
(20)
then Eq.(19) can be rewritten as
$$\sigma _{\underset{¯}{x}\underset{¯}{e}}=\underset{e}{}w_e|ee|\rho _{\underset{¯}{x}}^e.$$
(21)
By a similar argument, $`\sigma _{\underset{¯}{y}\underset{¯}{e}}`$ is also diagonal in its $`_{\underset{¯}{e}}`$ sector and can be expressed as
$$\sigma _{\underset{¯}{y}\underset{¯}{e}}=\underset{e}{}w_e|ee|\rho _{\underset{¯}{y}}^e,$$
(22)
where for all $`e`$, $`\rho _{\underset{¯}{y}}^e𝒟(_{\underset{¯}{y}})`$. Our newly found, diagonal in the $`_{\underset{¯}{e}}`$ sector, expressions for $`\sigma _{\underset{¯}{y}\underset{¯}{e}},\sigma _{\underset{¯}{x}\underset{¯}{e}}`$ and $`\sigma _{\underset{¯}{e}}`$ can now be substituted into Eq.(15) to get
$$\sigma =\underset{e}{}w_e|ee|\rho _{\underset{¯}{x}}^e\rho _{\underset{¯}{y}}^e.$$
(23)
Thus, $`\sigma `$ is $`\underset{¯}{x},\underset{¯}{y}|\underset{¯}{e}`$ separable. Taking the $`\underset{¯}{e}`$ trace of this $`\sigma `$ to get $`\rho `$, we see that $`\rho `$ is $`\underset{¯}{x},\underset{¯}{y}`$ separable. QED
There probably exist certain $`\rho 𝒟(_{\underset{¯}{x}\underset{¯}{y}})`$ for which conditions 1 to 4 on the right hand side of Theorem 1 cannot be achieved for finite $`dim(_{\underset{¯}{e}})`$, but can be achieved in the limit $`dim(_{\underset{¯}{e}})\mathrm{}`$. Such $`\rho `$ could be described as being weakly separable.
Let $`𝒟_{insep}(_{\underset{¯}{x}\underset{¯}{y}})`$ be the set of all $`\rho 𝒟(_{\underset{¯}{x}\underset{¯}{y}})`$ which are not $`\underset{¯}{x},\underset{¯}{y}`$ separable. Let $`𝒟_{pos}(_{\underset{¯}{x}\underset{¯}{y}})`$ be the set of all $`\rho 𝒟(_{\underset{¯}{x}\underset{¯}{y}})`$ for which all extensions $`\sigma 𝒟(_{\underset{¯}{x}\underset{¯}{y}\underset{¯}{e}})`$ such that $`\rho =\mathrm{tr}_{\underset{¯}{e}}\sigma `$ satisfy $`S_\sigma [(\underset{¯}{x}:\underset{¯}{y})|\underset{¯}{e}]0`$. Then, by Theorem 1, $`𝒟_{pos}(_{\underset{¯}{x}\underset{¯}{y}})𝒟_{insep}(_{\underset{¯}{x}\underset{¯}{y}})`$. Density matrices in $`𝒟_{pos}(_{\underset{¯}{x}\underset{¯}{y}})`$ and those in $`𝒟_{insep}(_{\underset{¯}{x}\underset{¯}{y}})𝒟_{pos}(_{\underset{¯}{x}\underset{¯}{y}})`$ exhibit different kinds of entanglement.
Some goals for future research are: give concrete examples of Theorem 1; explore the connection between Theorem 1 and the necessary condition for separability given by Peres, and the bound entanglement discovered by Horodecki..
Acknowledgements: I thank M.A. Nielsen, D. Petz and M.B Ruskai for their generosity in communicating to me that the quantum strong subadditivity inequality is saturated iff Eq.(6).
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# REFERENCES
Consequences of parton’s saturation and string’s percolation
on the developments of cosmic ray showers
C. Pajares, D. Sousa, and R.A. Vázquez
Departamento de Física de Partículas, Universidade de Santiago
E-15706 Santiago de Compostela, Spain
Pacs Numbers 13.85.Tp, 13.85.-t, 96.40.De
## Abstract
At high gluon or string densities, gluons’ saturation or the strong interaction among strings, either forming colour ropes or giving rise to string’s percolation, induces a strong suppression in the particle multiplicities produced at high energy. This suppression implies important modifications on cosmic ray shower development. In particular, it is shown that it affects the depth of maximum, the elongation rate, and the behaviour of the number of muons at energies around $`10^{17}`$$`10^{18}`$ eV. The existing cosmic ray data point out in the same direction.
One of the most crucial astrophysical issues of the highest energy cosmic rays (above $`10^{17}`$ eV) is that of their composition. This problem is linked to the identification of the origin and possible sources of these cosmic rays. Current theoretical models expect a transition from galactic to extragalactic or galactic halo sources near the region of the ankle which leads to the usual expectation of the changing of composition from heavy to light elements.
Experimentally, measuring the composition at these energies is a challenging task. The very low fluxes involved imply that one has to rely on indirect measurements which depend on simulations of the development of cosmic ray cascades in the atmosphere. These, in turn, are model dependent and, specifically, depend on extrapolations of hadronic models to energies and kinematical regions never measured in the laboratory. There is, therefore, some degree of uncertainty in the shower development and one may ask what is the effect of this uncertainty in the reconstruction of shower parameters, mainly total energy and mass composition.
To avoid this problem, experimental groups have concentrated on observables which are expected to be more or less independent of the hadronic model used, or which have its dependence under theoretical control. These parameters include the maximum of the cascade, $`X_{\mathrm{max}}`$, the slope parameter, $`\beta =d\mathrm{log}(\rho _\mu (600))/d\mathrm{log}(E)`$, where $`\rho _\mu (600)`$ is the muon density at 600 m from the core, and the elongation rate, $`D_{10}=dX_{\mathrm{max}}/d\mathrm{log}_{10}(E)`$. Other parameters have been less frequently used, see Ref. for a general review.
Several experimental groups have measured the cosmic ray spectrum and mass composition in the ankle region and beyond using the above mentioned parameters see also. The results on mass composition are inconclusive. Fly’s Eye observe a change on the slope of $`X_{\mathrm{max}}`$ in the region around $`3\times 10^{17}`$ eV, which is interpreted as a change on the composition from heavy (iron) dominated to light (proton) dominated. However, AGASA measures a muon component and $`\beta `$ parameter consistent with iron on that region. Although some part of the discrepancy between AGASA and Fly’s Eye may be due to the use of different hadronic models , as pointed out by Nagano et al. the important issue is that AGASA sees no significant change on the muon component of showers all along the ankle region and beyond, from $`10^{16.5}`$ eV to $`10^{19.5}`$ eV, and thus no strong change on composition is inferred. HIRES and MIA collaborations have jointly measured both the $`X_{\mathrm{max}}`$ and the $`\beta `$ parameter. They observe a strong change of $`X_{\mathrm{max}}`$ with energy, which implies a large elongation rate, $`D_{10}=95`$ gr/cm<sup>2</sup> and on the other hand they see no change on the slope of the muonic component, $`\beta =0.73`$, measured which is broadly compatible with the AGASA observations. HIRES and MIA however have measured these parameters in a narrow range of energy, from $`10^{17}`$ to $`10^{18}`$ eV and with low statistics, only during a limited exposure.
In this paper we will show that under rather general conditions a change on the hadronic interactions at the energies of interest is expected, which may have important consequences for the interpretation of cosmic ray data. Whether this change is enough or not to produce the observed changes on the cosmic ray data we can not tell at present. On the other hand we can state the necessary conditions for this change to explain the observed data: i) There should be an abrupt change on the hadronic interactions at the observed energy $`E_{\mathrm{lab}}5\times 10^{17}`$ eV for Fe–Air collisions. This corresponds to a CM energy of $`4200`$ AGeV and a density of gluons of $`9`$ fm<sup>-3</sup> ii) At this energy the slope of the growth of the multiplicity with energy should vary from $`0.24`$ to a essentially flat $`0.09`$. If i) and ii) are verified then there is no additional need for a change on composition to explain the data.
It is important to point out that, although the change on the multiplicity may or may not be enough to produce the observed results, some effect should always be present and should be taken into account in any realistic simulation of cosmic ray showers. Currently no Monte Carlo code for cosmic ray showers has yet been implemented with these effects<sup>*</sup><sup>*</sup>*Sibyll version 2.0 incorporates some shadowing effect. However this was done for $`pp`$ collisions only and does not affect to our reasoning below.
In the last years, a wealth of data coming from HERA and the heavy ion SPS experiments have risen questions about the behaviour of the hadronic interactions at very high energy. We may consider perturbative, gluons, or non perturbative, strings, as the fundamental variables of our description, At high gluon density, the saturation of gluons and/or a strong jet shadowing are expected. In the case of high string density we expect the fusion of strings or colour rope, and probably, above a critical string density the percolation of strings and the formation of quark gluon plasma at the nuclear scale are expected.
One general feature of all these hadronic phenomena is the strong suppression of particle multiplicity compared to the multiplicity expected in their absence. Namely, for central Pb–Pb collisions the charged particle multiplicity expectations in the central rapidity region changes between 1500 (7500) for the relativistic heavy ion collider, RHIC, (the large hadron collider, LHC) for models that do not include these effects to 900 (3000) when they are included. As a framework we will use the quark gluon string model (QGS), a modified version of the Dual Parton Model. The model is based on the large $`N`$ expansion of QCD but it is largely phenomenological and describes most of the soft hadronic physics rahter well. Inclusion of hard, perturbative, physics has been done in various ways. In the quark–gluon string model, multiparticle production is related to the interchange of multiple strings which break and subsequently hadronize.
In this model one can most easily understand the expected changes on the behavior of hadronic collisions at high energies. It is more convenient to work in the plane transverse to the collision. In this plane, strings are seen as small circles of fixed radius, $`r`$. As the energy increases, the number of strings interchanged increases and the total area occupied by strings increase. At high energy, strings start to overlap and fuse together. For high enough string density, $`n_c`$, strings may percolate in a second order phase transition, i.e. continuous paths of strings are formed in the collision area. Since the number of independent strings is reduced after the fusion one expects a depletion on the number of particles produced, i.e. a reduction on the multiplicity.
In the QGS the multiplicity grows with energy as $`n(s)s^\mathrm{\Delta }`$, where $`\mathrm{\Delta }`$ is related to the intercept of the soft pomeronHere we consider minimum bias events, which are relevant for cosmic ray experiments. The parameter $`\mathrm{\Delta }`$ may depend on the centrality of the collision. In the case of percolation, the reduction of multiplicity is given by
$$n^{}(s)=n(s)\sqrt{(}F(\eta )),$$
(1)
where
$$F(\eta )=\frac{1e^\eta }{\eta },$$
(2)
and the parameter $`\eta `$ is the fraction of the total area occupied by strings
$$\eta =\frac{\pi r^2N_s}{\pi R^2}.$$
(3)
Here $`N_s`$ is the number of strings produced in the collision, $`r`$ is the string’s transverse size, and $`R`$ is the total collision area. $`N_s`$ grows with energy as $`N_ss^\mathrm{\Delta }^{}`$, where $`\mathrm{\Delta }^{}`$ is the intercept of the soft pomeron. Therefore at large $`\eta `$ the total multiplicity grows with energy as
$$n^{}(s)s^{\mathrm{\Delta }\mathrm{\Delta }^{}/2}.$$
(4)
This reduction of multiplicity is not exclusive of the percolation of strings. For instance, in perturbative QCD a reduction in the number of jets produced as the energy increases is also expected. At high energy the number of jets produced grows with energy as $`n(s,p_t^2)(s/(4p_t^2))^{\mathrm{\Delta }_H}`$, where $`\mathrm{\Delta }_H`$ is the intercept of the hard pomeron and $`p_t`$ is the transverse momentum of the jet. A (mini)jet occupies a transverse area of order $`\pi /p_t^2`$, since the number of jets increases rapidly with energy, saturation occurs when the area occupied by the jets equals the total transverse area :
$$\frac{n(s,p_t^2)\pi /p_t^2}{\pi R^2}=1,$$
(5)
which implies
$$n(s,p_t^2)s^{\frac{\mathrm{\Delta }_H}{1+\mathrm{\Delta }_H}}.$$
(6)
Both Eqs.(4,6) have been checked directly in Monte Carlo simulation . Surprisingly, for nucleus–nucleus collisions the reduction in the power of multiplicity growth with energy is of the same order both for the case of string fusion and of shadowing. The power changes from $`0.24`$ to $`0.19`$. In general parton saturation, shadowing, string fusion, or percolation will produce the effect of reduction of the multiplicity although we expect the degree of this reduction to be model dependent.
For cosmic ray showers, the rate of change of the multiplicity with energy is directly related to the change of the elongation rate. This has been known for a long time as the elongation rate theorem . The elongation rate theorem can be deduced easily, as it follows from a scaling argument. Let $`X_{\mathrm{max}}(E)`$ be the maximum depth of the shower produced by a primary of energy $`E`$. On average, the first interaction occurs at depth $`\lambda `$, the mean free path of the initial particle. In this first interaction the initial particle splits into $`n(E)`$ particles each carrying an average energy $`E/n(E)`$. Therefore, we have
$$X_{\mathrm{max}}(E)=\lambda +X_{\mathrm{max}}(E/n(E)).$$
(7)
Assuming that $`X_{\mathrm{max}}(E)`$ depends logarithmically on energy we get
$$X_{\mathrm{max}}(E)=A\mathrm{log}_{10}(E/n(E))+B,$$
(8)
where $`A=X_0\mathrm{ln}10`$ and $`B`$ are constants. $`X_0=37`$ gr/cm<sup>2</sup> is the electromagnetic radiation length. If we now assume that $`n(E)E^\mathrm{\Delta }`$, we get
$$X_{\mathrm{max}}(E)=A(1\mathrm{\Delta })\mathrm{log}_{10}(E)+B^{}.$$
(9)
This is the elongation rate theorem. We can now directly read the elongation rate from the above equation $`D_{10}=A(1\mathrm{\Delta })`$. As stated previously, a change in the multiplicity growth with energy implies a change in the elongation rate.
In Fig.(1) we show $`X_{\mathrm{max}}`$ as a function of energy for the Fly’s Eye and HIRES-MIA experiments. Data have been taken from references . The errors shown are only statistical. An additional systematic error of $`20`$ gr/cm<sup>2</sup> must be included in the data. The dash line represents our calculation for the slope parameter, $`D_{10}=65`$ gr/cm<sup>2</sup>,($`\mathrm{\Delta }=0.24`$ from our simulations) for Fe–Air collision without fusion, normalized with the data at $`6\times 10^{17}`$ eV. The dotted curve has a slope parameter $`D_{10}=78`$ gr/cm<sup>2</sup>, which would imply a maximum reduction in the slope of growing of multiplicities: from $`\mathrm{\Delta }0.24`$ to $`\mathrm{\Delta }0.09`$. The data from HIRES–MIA is not completely consistent with the Fly’s Eye data. Statistical uncertainties are larger. The elongation rate obtained by the HIRES-MIA collaboration is very large, $`95`$ gr/cm<sup>2</sup>. Notice that the elongation rate theorem predicts an elongation rate always less than 85 gr/cm<sup>2</sup>, an elongation rate larger would imply multiplicities decreasing with energy. Therefore, if the HIRES-MIA result is confirmed, a change of the composition is necessary to explain the data. In the figure we also show the energy region in which a phase transition is expected in Fe–Air collision using the string fusion model. This region corresponds to the range obtained from percolation theory, $`1.1\eta 1.2`$, where $`\eta `$ is given by Eq.(3).
A word of caution is necessary in the use of the elongation rate theorem. The elongation rate theorem is based on the assumption that the energy is equally shared between the secondaries in the hadronic interaction. From this assumption immediately follows the logarithmic dependence of $`X_{\mathrm{max}}`$ on the multiplicity. In realistics cases this assumption does not hold and one has to resort to simulations since no analytical formula is known for the $`X_{\mathrm{max}}`$. We have parameterized for a number of models the dependence of $`X_{\mathrm{max}}`$ on the change of multiplicity, see Ref. for details. The results of a full Monte Carlo agree with our qualitative discussion. With this in mind we can conclude that to be able to explain the change on the elongation rate we need a change on the slope of the multiplicity from $`\mathrm{\Delta }0.24`$ to an essentially flat $`\mathrm{\Delta }0.09`$. The energy region for such change must be around $`5\times 10^{17}`$ eV.
The experimental situation with the lateral distribution of muons $`\rho _\mu (r)`$ is clearer. Both in simulations and experiments it is found that the shape of the lateral distribution function for muons is rather independent of the primary’s energy and composition. Therefore, at fixed distance to the core, $`r_0`$, $`\rho _\mu (r_0)`$ is proportional to the total number of produced muons in the shower. Under rather general arguments this number scales with energy
$$\rho _\mu (r_0)N_\mu =AE^\beta ,$$
(10)
where $`A`$ is a normalization constant and $`\beta `$ is the slope parameter. As mentioned previously, the slope parameter is found to be constant over a wide range of energies. This result is consistent with Yakutsk, Haverah Park , and with all the lower energy experiments.
It is rather simple to calculate the slope parameter, $`\beta `$, for a pionic cascade from a scaling argument similar to the elongation rate theorem. The number of muons is proportional to the number of charged pions at the maximum. The number of pions, at maximum, produced by a primary of energy $`E_0`$ is given by
$$N_\pi (E_0)=f_\pi n_0^1𝑑xP(x)N_\pi (xE_0),$$
(11)
where $`f_\pi =2/3`$ is the charged pion fraction, $`n`$ the total pion multiplicity, and $`P(x)`$ is the probability of producing a pion with a fraction of energy $`x`$ of the primary energy. Assuming a scaling form, $`N_\pi =AE^\beta `$, we get
$$1=f_\pi n_0^1𝑑xP(x)x^\beta =f_\pi Z(\beta ),$$
(12)
where $`Z(\beta )`$ is the spectrum–weighted momentum. For a given $`P(x)=1/ndn/dx`$, the above equation gives an implicit equation for $`\beta `$. It reduces to the textbook’s expression if we assume energy equipartition, i.e. $`P(x)=\delta (x1/n)`$, which gives $`\beta =(1+\mathrm{log}(f_\pi ))/\mathrm{log}(n)0.82`$, for $`n10`$. For realistic models the slope parameter $`\beta `$ ranges between 0.7–0.9. In Eq.(12) the multiplicity enters explicitly in the left hand of the expression but also enters implicitly since the probability $`P(x)`$ must verify total probability and energy conservation. Since $`Z(\beta )`$ is a monotonically decreasing function of $`\beta `$ for reasonable choices of $`dn/dx`$, a reduction of multiplicity induces a reduction on $`\beta `$. Indeed this is what is observed for the $`dn/dx`$ calculated for the model with and without fusion. The DPM model gives a value of $`\beta 0.89`$ which agrees with that of the QGSJET model . In the presence of fusion the slope parameter is reduced and we get $`\beta `$ 0.72–0.77 depending on the specific implementation of the fusion model. This number was calculated both using Eq.(12) and by direct calculation with a Monte Carlo code.
In Fig.(2), we show de density of muons at 600 m, $`\rho _\mu (600)`$, as a function of the energy for the AGASA measurements given as a parameterization and the HIRES-MIA measurements, shown as triangles. Also shown are the QGSJET results for pure iron and proton. Notice that the slope parameter measured by HIRES-MIA agrees with our slope parameter for the case of fusion. Our result, a change of slope parameter from that of pure iron to 0.77 at the energy where percolation is expected, is shown for comparison only. We can see that again it is consistent with the data. The slope measured by AGASA is different from the one calculated for either proton or iron for the QGSJET model. A rapid change on composition, as suggested by the HIRES-MIA data on $`X_{\mathrm{max}}`$, would imply a kink in the data for $`\rho _\mu (600)`$ at the same energy which is not seen. Instead, our result points towards a mild change on the slope parameter, from 0.9–0.8 to 0.77 which would be hardly seen given the error bars in the data.
There are a number of additional predictions in our scenario. The average $`p_t`$ in hadronic collisions should increase in the case of string fusion by about 10 – 20 %. This would produce flatter lateral distribution for the muon densities which could be observed. A particularly well–suited place to look for this effect would be in inclined showers. Inclined showers are composed esentially of high energy muons, and therefore are more sensitive to changes on the first hadronic interactions. Given the current systematical and statistical errors we can not conclude that indeed cosmic ray experiments are observing a saturation of gluons or percolation of strings in hadronic interactions. High quality data with large statistics, as the expected from HIRES and the Pierre Auger observatories, are needed. However it is suggestive that all cosmic ray existing data are consistent with such interpretation. RHIC and LHC will measure the total multiplicity in the relevant energy region and ascertain whether a strong reduction of multiplicity takes place or not. But in any case, even if the change of composition is real, these effects must be taken into account in a complete simulation of cosmic ray showers.
###### Acknowledgements.
We thank N. Armesto, E.G. Ferreiro, C. Merino, and E. Zas for useful discussions. This work has been partially supported by CICYT (Spain), AEN99-0589-C02-02.
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# Parametric Quantum Resonances for Bose-Einstein Condensates
## I INTRODUCTION
There has been a blossoming of literature on the features of systems exhibiting Bose-Einstein condensation (BEC), triggered by its recent experimental realization. Intially experiments with $`10^3`$ to $`10^6`$ atoms of rubidium or sodium (later experiments have used lithium and eventually (spin-polarized) hydrogen), in harmonic or cigar shaped traps have demonstrated condensation to a “pseudo-macroscopic” level of occupancy of the ground state for $`nK`$ temperatures. Time of flight measurements, velocity distributions as well as spatial profiles have convincingly supported the physical picture of an abrupt transition in the behavior of the Bose gas, which has been interpreted as the signature of BEC.
Following these experiments, many theoretical studies were launched to characterize different aspects of Bose condensates such as hydrodynamic modes, collective excitations, the behavior of ideal quantum fluids, the fraction of noncondensate vs. condensate atoms, or the generation and stability of vortices. In turn, experimental studies have progressed to address some of the theoretical predictions and open up new questions.
Here, we concern ourselves with one aspect of these quantum fluids, namely parametric driving. For the purpose of this report, we will restrict ourselves to the framework of the mean-field or Hartree-Fock approximation. This approximation is rigorously justifiable only at $`T=0`$ but it is expected that the contribution of the non-condensate to the density is quite small. It is well-known that at this mean-field level the condensate wavefunction is governed by the Gross-Pitaevskii (GP) equation. An issue addressed after the original experiments achieving the condensation was the study of collective excitations. In these papers these excitations were induced by a harmonic trap weakly modulated in time with appropriate types of symmetry. More recently, it was demonstrated that extended parametric resonances can occur in a two-dimensional (2d) NLS equation with a harmonic trap. This result may or may not (for reasons to be explained below) be relevant for two dimensional studies of Bose gases. However, this naturally raises the question of whether a similar result can be deduced for the $`3d`$ case which is certainly of direct relevance to experimental studies.
The main question we will address is whether weak harmonic modulation of trapped $`3d`$ condensates can cause an anomalously large response in their wavefunction. Our answer, which will be in the affirmative, will be motivated by mathematical analysis using a moment method and verified by numerical simulation. We will briefly discuss the implications of these results and the suggestion of relevant experiments.
## II MOMENT METHODS
Considering a spherical trap, the dimensionless GP equation for the dynamics of the BEC condensate is
$`iu_t={\displaystyle \frac{1}{\zeta ^4}}^2u+\left(\lambda (t)r^2+\nu |u|^2\right)u.`$ (1)
Here, the subscript $`t`$ denotes time derivative and $`\zeta =(8\pi N|a|/a_{})^{1/5}`$ is a dimensionless parameter arising from the number $`N`$ of particles, the s-wave scattering length $`a`$ and from $`a_{}`$ characterizing the strength of the trap (see, Ref. ). In Eq.(1) $`\lambda (t)`$ is a dimensionless function allowing for time dependence of the trap and $`\nu =\text{sign}(a)`$, generalizes the equation to describe attractive ($`a<0)`$ as well as repulsive ($`a>0`$) interactions. Since, we restrict ourselves to spherical symmetry we only include the radial contribution in the Laplace operator. Although we are mainly interested in the full three-dimensional ($`3d`$) case we will in general consider the $`d`$-dimensional version of Eq.(1) so that
$$^2=\frac{1}{r^{d1}}\frac{}{r}\left(r^{d1}\frac{}{r}\right).$$
(2)
Similarly to Ref. , we define the following quantities
$`I_{2,a}^{(d)}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}r^a|u|^2r^{d1}𝑑r,`$ (3)
$`I_{3,a}^{(d)}`$ $`=`$ $`i{\displaystyle _0^{\mathrm{}}}r^a(uu_{r}^{}{}_{}{}^{}c.c.)r^{d1}dr,`$ (4)
$`I_{4,a}^{(d)}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}r^a\left|{\displaystyle \frac{u}{r}}\right|^2r^{d1}𝑑r,`$ (5)
$`I_{5,a}^{(d)}`$ $`=`$ $`{\displaystyle _0^{\mathrm{}}}r^a|u|^4r^{d1}𝑑r,`$ (6)
where $`(d)`$ indexes the dimension. This type of nonlinear Schrödinger (NLS) equation (1) has two conserved quantities: as a result of the phase invariance the norm corresponding to $`I_{2,0}^{(d)}`$ is conserved in any dimension $`d`$. Also, since Eq.(1) is a Hamiltonian system arising from
$`H={\displaystyle _0^{\mathrm{}}}\left[\zeta ^4|u|^2+{\displaystyle \frac{\nu }{2}}|u|^4+\lambda (t)r^2|u|^2\right]r^{d1}𝑑r`$ (7)
this quantity is conserved. It is useful to note that the Hamiltonian (or more appropriately energy functional), $`H`$, can be expressed in terms of the moments Eqs.(3)-(6)
$`H=\zeta ^4I_{4,0}^{(d)}+{\displaystyle \frac{\nu }{2}}I_{5,0}^{(d)}+\lambda (t)I_{2,2}^{(d)}.`$ (8)
The relevance of the moments Eqs.(3)-(6), is based on their time evolution and in the following we will derive the relations governing this dynamics. The physical rationale behind such an approach lies in the fact that the resulting equations can yield predictive diagnostics for the dynamics of the BEC. In particular, $`I_{2,2}^{(d)}`$ will essentially yield the width of the spatial profile of the wavefunction. If a condensation phenomenon is to take place even when starting from a spatially uniform distribution (at temperatures $`T<T_c`$), the width must evolve towards a constant non-zero value.
Using Eq.(1) and its complex conjugate one can derive
$`\dot{I}_{2,a}^{(d)}`$ $`=`$ $`\zeta ^4aI_{3,a1}^{(d)},`$ (9)
$`\dot{I}_{3,a}^{(d)}`$ $`=`$ $`4\lambda (t)I_{2,a+1}^{(d)}+4a\zeta ^4I_{4,a1}^{(d)}`$ (10)
$``$ $`(a+d1)(a1)(a3+d)\zeta ^4I_{2,a3}^{(d)}`$ (11)
$`+`$ $`\nu (a+d1)I_{5,a1}^{(d)}.`$ (12)
In deriving these we assume $`u`$ to vanish as $`r\mathrm{}`$.
Unfortunately, it is in general impossible to close this hierarchy of equations because the time derivative couples to the next order e.g. $`\dot{I}_{2,a}^{(d)}`$ couples to $`I_{3,a1}^{(d)}`$ and $`\dot{I}_{3,a}^{(d)}`$ couples to $`I_{5,a1}^{(d)}`$ and $`I_{4,a1}^{(d)}`$, and so on. However, combining Eqs. (9) and (11) provides some insight
$`\ddot{I}_{2,a}^{(d)}`$ $`=`$ $`\zeta ^4a[4\lambda (t)I_{2,a}^{(d)}+4(a1)\zeta ^4I_{4,a2}^{(d)}`$ (13)
$``$ $`(a+d2)(a2)(a4+d)\zeta ^4I_{2,a4}^{(d)}`$ (14)
$`+`$ $`\nu (a+d1)I_{5,a2}^{(d)}].`$ (15)
First, this confirms that the norm $`I_{2,0}^{(d)}`$ is conserved in any dimension. Secondly, this relation clearly suggests $`a=2`$ as good choice since the term involving $`I_{2,a4}^{(d)}`$ then vanishes irrespective of dimension. Also, this choice allows the use of Eq.(8) to reduce the expression (15) to
$$\ddot{I}_{2,2}^{(d)}=8\zeta ^4H16\zeta ^4\lambda (t)I_{2,2}^{(d)}+2\nu (d2)I_{5,0}^{(d)},$$
(16)
which corresponds to the relation commonly referred to as the virial theorem for the nonlinear Schödinger without a trap, $`\lambda (t)0`$.
Clearly, the $`2d`$ case is special as the structure of Eq.(16) is such that a closed time evolution can be prescribed. In addition, for time dependent modulation of the trap amplitude, we find (as was observed for this problem in Ref. and again in Ref. ) a Hill type equation which establishes parametric resonances for the behavior of the width (or amplitude) of the wavefunction. Important as this conclusion about the $`2d`$ behavior may be for general NLS-GP equations, it is not clear that it is relevant to BEC. Since BEC is not possible in spatial dimensions less that three ($`d<3`$) where a Kosterlitz-Thouless topological transition seems to be occurring instead, the applicability of the GP equation for $`d<3`$ is controversial. Our search for condensate instabilities is therefore most compelling in three dimensions where no such reservations exist.
Although Eq.(16) is not closed in $`3d`$, it is easily seen that the following inequalities hold for $`d2`$
$`\ddot{I}_{2,2}^{(d)}+16\zeta ^4\lambda (t)I_{2,2}^{(d)}`$ $``$ $`8\zeta ^4H\text{for}\nu <0,`$ (17)
$`\ddot{I}_{2,2}^{(d)}+16\zeta ^4\lambda (t)I_{2,2}^{(d)}`$ $``$ $`8\zeta ^4H\text{for}\nu >0,`$ (18)
where the equality applies to the two-dimensional case only (and in fact for the noninteracting $`3d`$ case $`\nu =0`$). The value of these inequalities lies in the predictions about the 3d case. Since we can resolve, or at least very well characterize, the 2d behavior, we are now able to extend this to quantitative predictions about of the 3d behavior. For the attractive case $`\nu <0`$ the possibility of collapse occurs as $`I_{2,2}^{(d)}`$ can become zero in finite time. In $`2d`$ and due to Eq.(18) also in $`3d`$, a sufficient condition for collapse is $`H<0`$, although, depending on the initial configurations, collapse can be achieved even for $`H>0`$. A more complete discussion is given in Ref. . For the more realistic case (in BEC contexts) of repulsive interaction $`\nu >0`$ we see for example that the parametric resonances that were demonstrated in the two-dimensional case will also be present in the three-dimensional case. So, for instance, for $`d=3`$ and $`\lambda (t)=\lambda _0^2(1+ϵ\mathrm{cos}(\omega t))`$ will exhibit resonance around $`\lambda _0=n\omega /2,n=1,2,3..`$ where the extent of the first resonance is determined by the inequality $`|1\omega ^2/(4\lambda _0^2)|<ϵ/2`$. However, additional resonances may be possible in $`3d`$ due to the nonlinear driving resulting from the repulsive interaction. Some understanding of the influence of the last term in Eq.(16) on the dynamics can be gained by assuming that the wavefunction can be approximated as
$`u=\sqrt{{\displaystyle \frac{I_{2,0}^{(3)}}{\sigma _1}}}B^{3/2}\psi (r/B),`$ (19)
where $`\sigma _1=_0^{\mathrm{}}r^2\psi (r)𝑑r`$ is a shape-dependent constant. Thus, assuming adiabatically the wavefunction does not alter its shape $`\psi `$ significantly as a result of the dynamics, $`u`$ as defined in Eq. (19) automatically satisfies the norm conservation. Utilizing this in Eq.(16) yields
$`{\displaystyle \frac{d^2B^2}{dt^2}}\lambda (t)B^2=Q_1+\nu Q_2B^3,`$ (20)
where $`Q_1`$ and $`Q_2`$ are constants determined by the shape $`\psi `$ and the initial value of $`B`$. Clearly, the last term in Eq.(20) will only influence the dynamics when $`B`$ becomes small. In the repulsive $`\nu >0`$ case however $`B`$ will generally not become small since there is no collapse. This simple analysis suggests that the parametric forcing of the experimentally realizable $`3d`$ case will result in a resonance picture analogous to that previously reported for the $`2d`$ problem. Our numerical simulations of the full Eq.(1) with $`\lambda (t)`$ as given above verifies the validity of this prediction. A typical example for $`\lambda _0=1`$, $`\omega =1`$, (i.e. $`n=2`$) and $`ϵ=0.05`$ is given in Fig.1. The response of the wavefunction (whose initial condition had $`\text{max}_x|\psi (x,t=0)|^2=1`$), corresponding to the parametric resonance, can be observed directly from the wave function Fig. 1(a) but even more clearly in the time evolution of $`I_{2,2}^{(3)}`$, as shown in Fig. 1(b).
## III CONCLUSION AND FUTURE CHALLENGES
In this paper, we have presented and extended the formalism of the moment method, used in Refs. and for the $`2d`$ GP equations, to the more relevant $`3d`$ case. We have commented on the special nature of the two-dimensional problem where the moment equations form a closed set of equations. We have also added a note of caution in considering the results of the GP analysis for $`d<3`$. It might well be that, analogous to mean-field analysis in statistical physics systems, the “critical dimension” for this system is, indeed, $`d_c=3`$ and for lower dimensionalities the predictions of the mean-field theory are unreliable. A satisfactory self-consistent first principles description of an interacting boson gas for $`d<3`$ presents a very challenging theoretical problem, and it remains an unresolved issue whether a transition is present (and if it is, what is its nature).
On the other hand, we have used the moment methods and have derived results for GP functionals in all dimensions of physical interest. Considering, in particular, the $`3d`$ case, where the validity of the GP approximation is clear, we have obtained a non-closed set of equations for the moments of the wavefunction. We have demonstrated that for a parametric time-disturbance of the trap amplitude, parametric resonances are possible. To date the experiments that have used parametric modulation have not observed such phenomena. These experiments have been performed in cigar-shaped traps (where the analysis is considerably more complicated even in the $`2d`$ problem). No fine tuning of frequencies and amplitudes was explored since the aim of the studies was to excite collective modes rather than to observe parametric resonances. Hence, we propose an experiment in a spherical trap using weak harmonic modulation of the condensate. Given the current experimental advances (see e.g. Ref. for a review), such an experiment seems feasible. Such a study would, apart from the validation of the theoretical prediction, also explore how such resonances might destabilize the condensate.
## ACKNOWLEDGMENTS
PGK gratefully acknowledges fellowship support from the “A.S. Onasis” Public Benefit Foundation and assistantship support from the Computational Chemodynamics Laboratory of Rutgers University. Research at Los Alamos National Laboratory is performed under the auspices of the US DOE.
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# The Boltzmann equation for colorless plasmons in hot QCD plasma. Semiclassical approximation
## 1 Introduction
For about three decades there has been an increasing interest in theoretical research into various dynamical properties of (ultra)relativistic many-particle systems. It is connected with manifold applications to various problems in astrophysical systems, modern cosmology, in multiparton processes in experiment with high energy heavy ion collisions etc. The kinetic phenomenons, having a purely collective character, are one of the more important aspects of complicated dynamics of many-particle systems under extreme conditions. Here, the basic element in the description of transport phenomena is the derivation of the corresponding kinetic equations which would take into account (depending on the character of the problem being studied) the presence of mean fields in the system, two - (and more) body collisions, the possible renormalization effects, the effects of quantum fluctuations (stochastic), pair production, etc. Here, we restrict our consideration to a brief review of the results of the derivation of relativistic kinetic equations essentially based on two-body collisions in hot gauge theories.
At present there are a few methods of construction of relativistic collision integrals. In particular, we mention the Zubarev’s method of the non-equilibrium statistical operator (use of this method on a relativistic systems of quark-gluon plasma (QGP) type can be found in Ref. ), and the method developed by Klimontovich for an ordinary nonrelativistic plasma (the so-called second momentum or polarization approximation) and expanded to relativistic (semi)classical systems - in . It should be stressed that the above-mentioned methods are particularly effective in the construction of collision integrals for relativistic (semi)classical systems, whose evaluation is described by so-called exact ”microscopic” dynamical equations arising from the classical equations of motion. However, the extension of these methods to the relativistic quantum systems encounters some difficulties and therefore is uneffective. For the latter systems, mention may be made of the method based on the use of the reduction formulae of quantum field theory given by de Groot et al . A more powerful and more convenient tool to derive the approximate relativistic kinetic equations from exact field Shwinger-Dyson equations, is the so-called closed-time-path (CTP) formalism . Examples of its relativistic field-theoretical generalization can be found in .
As is known, the cornerstone of derivation of the kinetic equations for a hot non-Abelian plasma is a fundamental separation of the momentum scale. The physical justification for such a separation is the fact that the collective excitations which develop at a particular energy scale $`gT`$ (or $`g^2T`$, where $`T`$ is the temperature, and $`g`$ is the coupling constant), are well separated, when $`g1`$, from the typical energies of plasma particles $`T`$ . Generally speaking one can define two types of kinetic equations: the equations for hard <sup>1</sup><sup>1</sup>1For the connection with our previous papers here, we shall follow the definitions of the momentum scales, accepted in : the hard scale, corresponding to momentum of order $`T`$, the soft scale $`gT`$, and the ultrasoft scale $`g^2T`$. particles - hard quarks($`q`$), antiquarks($`\overline{q}`$) and hard transverse gluons($`\mathrm{g}`$), and equation for soft (ultrasoft) collective modes (in the case of Bose excitations - transverse and longitudinal modes, the latter are called by plasmons). Most efforts were directed at the derivation of the first type of equations. However for considerably excited states, when a characteristic time of relaxation of the hard particle distributions $`f_s,s=q,\overline{q},\mathrm{g}`$ is commensurable with a characteristic time of relaxation of the soft oscillations or even significantly exceeds it, along with kinetic equations for hard particles it is necessary to use the kinetic equation for the soft modes.
The calculation of the collision term for the quark-gluon plasma was apparently first made in . Within the framework of the concepts developed in the theory of electron-ion plasma, the scattering probability of (anti)quarks among themselves through longitudinal and transverse virtual oscillations which account for dynamical screening, is derived. We note that in this paper the Boltzmann equation for hard quarks and antiquarks was supplemented by an equation characterized the relaxation of field excitations and was presented as the second type of kinetic equation mentioned above. In paper a similar Boltzmann equation for hard gluons was obtained. Here, the scattering probability was deduced within the framework of usual diagrammatic perturbation theory with inclusion of screening effects in the random-phase approximation (one-loop order). Although the relativistic Boltzmann equations constructed in these papers take into account such an important QGP property as screening, the range of their validity restricts their use to just colorless deviations from equilibrium distribution functions.
In within the framework of the (semi)classical representation of QGP the Balescu-Lenard-type collision terms for small color and singlet deviations of the distributions from the initial colorless equilibrium, were derived by the Klimontovich method. However, in these papers, the organizing role of the various momentum scales was not recognized, resulting in some inconsistent and complicated transport equations for hard particles.
Only in recent years it been possible to derive the Boltzmann equation for hard modes of hot non-Abelian plasma by a rigorous and self-consistent way using an expansion in the coupling constant, and clearly clarify the nature of the approximation involved, and thus fix its range of applicability. It was be shown that for longer wavelengths ($`\lambda 1/g^2T`$) of color excitations in the non-Abelian plasma, not only was consideration of interaction of hard particles with the soft degrees of freedom represented by mean fields essential, but also taking into account the collisions of hard particles among themselves. A similar Vlasov-Boltzmann equation reproduces exactly (at leading order in $`g`$) a large variety of the thermal results obtained by a more fundamental analysis of the diagrammatic perturbation theory and provides in some cases a letter description of phenomena do not yield to a perturbative analysis.
Here, one can extract three approaches to constructing an effective kinetic equation for hard particles with collision term. The first of them is connected with Bödeker’s effective theory for the ultrarelativistic field modes . Starting from the collisionless non-Abelian Vlasov equation, which is the result of integrating out the scale $`T`$ , Bödeker has shown how one can integrate out the scale $`gT`$ in an expansion in the gauge coupling $`g`$. At leading order in $`g`$, he has obtained the linearized Vlasov-Boltzmann equation for the hard field modes, which besides a collision term also contains a Gaussian noise. Subsequently, this equation was also proposed by Arnold et al who derive the relevant collision term on phenomenological grounds – by analyzing the scattering processes between hard particles in the plasma. The kinetic equation derived in has a non-trivial matrix structure, since, the distribution function that describes color fluctuations is not diagonal in color space.
Afterwards, an alternative derivation of the collision term of Balescu-Lenard-type was proposed by Litim and Manuel, and by Valle Basagoiti . The former authors used a classical transport theory, whereas the latter used the set of ”microscopic” dynamical equations coming from the HTL effective action describing the evolution of the collisionless plasma. In both cases the collision terms were derived by averaging the statistical fluctuations in the plasma on the basis of the method developed by Klimontovich .
Blaizot and Iancu suggested a detailed derivation of the Vlasov-Boltzmann equation, starting from the Kadanoff-Baym equations. The derivation is based on the method of gauge covariant gradient expansion, which was first proposed by them for the collective dynamics at the scale $`gT`$ . By using the given equation they obtain the effective amplitudes for the ultrasoft color fields, which generalize the HTL’s by including the effects of the collisions (see also Guerin ).
The paper of Bezzerides and DuBois devoted to the non-thermal QED plasma, is one of the first papers, in which the relativistic kinetic equation for the soft correlation function was considered. Over 60-70 years in connection with application to thermonuclear fusion in the theory of the nonrelativistic electron-ion plasma, a powerful perturbative method (the so-called weak turbulent approximation) was developed \[22-24\] for research into various nonlinear plasma processes of the following types: on-shell scattering of soft modes on hard particles (the nonlinear Landau damping), three- and four-wave decays, etc. In spite of the fact that this method is not able to describe the phenomena connected with strong turbulence; nevertheless, it enables one within the framework of a unified scheme to encompass a wide class of plasma phenomena. The paper is the attempt at the extending of the above-mentioned weak turbulent theory to the electron-positron-photon plasma governed by the quantum electrodynamics. As a basic tool for such an extension, the CTP-formalism was used. However, since the main effort here was directed into investigating the collision integrals for hard electron and positrons, the authors have restricted their derivation to the plasmon kinetic equation, only taking into account pair production and the (linear) Landau damping process. Properly high-order processes, which we are interested in and are responsible for the nonlinear interaction mechanisms of the plasma waves, were not considered at all (see also the discussion of one-loop computations below).
A similar kinetic equation for soft Bose-modes in a QGP was apparently first made by Heinz et al . In the context of the imaginary time formalism in the one-loop approximation the imaginary part of the complete color linear response function was deduced and it is shown that it can be expressed in the form of the Boltzmann-Nordheim collision term. On the basis of such a derived rate of decay $`(\mathrm{\Gamma }_\mathrm{d})`$ and the rate for regeneration of the perturbations $`(\mathrm{\Gamma }_i)`$, the kinetic equation defining the evolution of a phase space distribution $`N(𝐱,t;𝐤,\omega )`$ of soft electric perturbations of the momentum $`k=(\omega ,𝐤)`$ in the form proposed by Weldon , was written as
$$\frac{\mathrm{d}N}{\mathrm{d}t}=N\mathrm{\Gamma }_\mathrm{d}+(1+N)\mathrm{\Gamma }_i.$$
(1.1)
As in a previous case, the higher orders, when $`\mathrm{\Gamma }_\mathrm{d}`$ and $`\mathrm{\Gamma }_i`$ itself can functionally (in the general case, nonlinearly) depend on $`N`$ here, were not considered. However the derivation of such dependence becomes important if we take into account that all computations in () were performed with the ”rigid” one-loop approximation, with bare propagators of massless gluons and quarks (electrons). However, as is known, quarks (electrons) and gluons inside the loop are not massless, they acquire the effective temperature-induced masses. The consequence of this fact is a kinematic prohibition of a decay of the soft perturbations into physical states. Furthermore by virtue of the fact that the phase velocities of both transverse and longitudinal eigenmodes of the plasma exceed velocity of light, the linear Landau damping is also absent. By virtue of the above-mentioned, the rates of decay and regeneration are just zero in this approximation.
This paper is devoted to further study of the kinetic equation for soft modes of the non-Abelian plasma. The theoretical framework of this paper is derived from synthesis of two formal developments. The first one is the development of a nonlinear theory of plasma wave interactions in ordinary plasma - more exactly the weak turbulent approximation \[22-24\]. The second is the development of effective theory of hot QCD originally proposed by Braaten and Pisarski , Frenkel and Taylor , Jackiw and Nair , Blaizot and Iancu and then developed in the papers \[15-20\]. In our previous papers without resorting to a complicated diagrammatic technique within the framework of the semiclassical representations, the following term in the expansions of $`\mathrm{\Gamma }_\mathrm{d}`$ and $`\mathrm{\Gamma }_i`$, linear on a phase space distribution of soft perturbations, was derived. However, as was shown in (see also section 2), this approximation is not sufficient for a complete definition of the relaxation process of the soft modes in QGP. Here, we consider the next terms in the expansions of $`\mathrm{\Gamma }_\mathrm{d}`$ and $`\mathrm{\Gamma }_i`$, and show that the corresponding nonlinear equation (1.1) is of purely Boltzmann type, i.e. the collision term on the r.h.s. of this equation has a standard Boltzmann structure, with a gain term and a loss term.
The outline of the paper is as follows. In section 2 the preliminary comments, with regard to derivation of the Boltzmann equation, describing the plasmon-plasmon scattering are explained. In section 3 the essential features of the scheme, which we used previously in to derive the kinetic equation with allowance for the nonlinear Landau damping, are summarized. In section 4 we discuss the consistency with gauge symmetry of the approximation scheme used. Section 5 is devoted to the determination of the interacting fields in the form of the expansion in free fields with the necessary accuracy for further research. In section 6 we select all terms in the expansion of the color random current responsible for the four-plasmon decay and derive the intermediate kinetic equation which then in section 7 will be rewritten in the terms of HTL-amplitudes. Section 8 is devoted to deriving the probability of plasmon-plasmon scattering, which is the main result of this work. In the next section on the basis of the explicit form of the obtained collision integral, the expression for lifetimes of colorless plasmons is defined and an estimate for the leading order in the coupling at the soft momentum scale is deduced. Finally in section 10 we present our conclusions and future avenues of study.
## 2 Preliminary comments
We denote the localized number density of the plasmons by $`N^l(𝐤,x)N_𝐤^l`$, and the distribution function of hard thermal gluons by $`f(𝐩,x)f_𝐩`$. In this paper we consider processes with longitudinal oscillations only, propagating in a purely gluonic plasma, with no quarks. Besides, we suppose that there is no external color current and/or mean color field in the system, and the system is in the global equilibrium state, i.e.
$$f_𝐩^{ab}=\delta ^{ab}f_𝐩\delta ^{ab}\mathrm{\hspace{0.17em}2}\frac{1}{\mathrm{e}^{E_𝐩/T}1}.$$
(2.1)
Here, $`E_𝐩|𝐩|`$ for a massless hard gluon, the coefficient 2 takes into account that the hard gluon has two helicity states and $`a,b=1,\mathrm{},N_c^21`$ for the $`SU(N_c)`$ gauge group. The triviality of the color structure of the plasmon number density is a consequence of these restrictions (see section 4)
$$N_𝐤^{lab}=\delta ^{ab}N_𝐤^l.$$
The dispersion relation for plasmons $`\omega =\omega ^l(𝐤)\omega _𝐤^l`$ is defined from
$$\mathrm{Re}\epsilon ^l(\omega ,𝐤)=0,$$
(2.2)
where
$$\epsilon ^l(\omega ,𝐤)=1+\frac{3\omega _{pl}^2}{𝐤^2}\left[1F\left(\frac{\omega }{|𝐤|}\right)\right],F(x)\frac{x}{2}\left[\mathrm{ln}\left|\frac{1+x}{1x}\right|i\pi \theta (1|x|)\right]$$
(2.3)
is longitudinal color permeability and $`\omega _{pl}^2=g^2N_cT^2/9`$ is a plasma frequency.
We expect the time-space evolution of $`N_𝐤^l`$ to be described by
$$\frac{N_𝐤^l}{t}+𝐕_𝐤^l\frac{N_𝐤^l}{𝐱}=N_𝐤^l\mathrm{\Gamma }_\mathrm{d}[N_𝐤^l]+(1+N_𝐤^l)\mathrm{\Gamma }_i[N_𝐤^l],$$
(2.4)
where $`𝐕_𝐤^l=\omega _𝐤^l/𝐤`$ is the group velocity of the longitudinal oscillations<sup>2</sup><sup>2</sup>2Notice, that in the general case, when the distribution function of hard gluons is a slowly varying function in time and space, equation (2.2) is replaced by $`\mathrm{Re}\epsilon ^l(\omega ,𝐤;t,𝐱)=0`$. In this case the l.h.s. of equation (2.4) should be supplemented by $`(\omega ^l(𝐤,x)/𝐱)(N^l(𝐤,x)/𝐤)`$.. For a generalized decay rate $`\mathrm{\Gamma }_\mathrm{d}`$ and inverse decay rate $`\mathrm{\Gamma }_i`$ it is shown in an explicit form that in general case they are functionals dependent on the plasmon number density. Although the approach we shall use in the subsequent discussion is correct only within the framework of semiclassical approximation, it is convenient to interpret the terms entering into $`\mathrm{\Gamma }_\mathrm{d}`$ and $`\mathrm{\Gamma }_i`$, using a quantum language.
The Eq. (2.4) in general, describes two principal processes of the nonlinear wave-interaction. The first of them represents the process of the stimulated emission and absorption of the collective wave quanta by hard particles of plasma. In this case the more general expression for the decay rate $`\mathrm{\Gamma }_\mathrm{d}^{(𝒮)}`$ and the regeneration rate $`\mathrm{\Gamma }_i^{(𝒮)}`$ can be written in the following forms, respectively:
$$\mathrm{\Gamma }_\mathrm{d}^{(𝒮)}[N_𝐤^l]=\underset{n,m}{}\frac{\mathrm{d}𝐩}{(2\pi )^3}d𝒯_{nm}^{(𝒮)}w(𝐩|𝐤,𝐤_1,\mathrm{},𝐤_n;𝐤_1^{},\mathrm{},𝐤_m^{})N_{𝐤_1}^l\mathrm{}N_{𝐤_n}^l$$
(2.5)
$$\times (1+N_{𝐤_1^{}}^l)\mathrm{}(1+N_{𝐤_m^{}}^l)f_𝐩[1+f_𝐩^{}],$$
and
$$\mathrm{\Gamma }_i^{(𝒮)}[N_𝐤^l]=\underset{n,m}{}\frac{\mathrm{d}𝐩}{(2\pi )^3}d𝒯_{nm}^{(𝒮)}w(𝐩^{}|𝐤_1^{},\mathrm{},𝐤_m^{};𝐤,𝐤_1,\mathrm{},𝐤_n)N_{𝐤_1^{}}^l\mathrm{}N_{𝐤_m^{}}^l$$
(2.6)
$$\times (1+N_{𝐤_1}^l)\mathrm{}(1+N_{𝐤_n}^l)f_𝐩^{}[1+f_𝐩].$$
Here, the phase-space integration is
$$d𝒯_{nm}^{(𝒮)}\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\mathrm{}\frac{\mathrm{d}𝐤_n}{(2\pi )^3}\frac{\mathrm{d}𝐤_1^{}}{(2\pi )^3}\mathrm{}\frac{\mathrm{d}𝐤_m^{}}{(2\pi )^3}$$
(2.7)
$$\times (2\pi )\delta (E_𝐩+\omega _𝐤^l+\omega _{𝐤_1}^l+\mathrm{}+\omega _{𝐤_n}^lE_𝐩^{}\omega _{𝐤_1^{}}^l\mathrm{}\omega _{𝐤_m^{}}^l),$$
with the delta function expressing the energy conservation of the processes of stimulated emission and absorption of the plasmons. The function $`w(𝐩|𝐤,𝐤_1,\mathrm{},𝐤_n;𝐤_1^{},\mathrm{},𝐤_m^{})`$ is the probability of absorption of $`n+1`$ plasmons with the frequencies $`\omega _𝐤^l,\omega _{𝐤_1}^l,\mathrm{},\omega _{𝐤_n}^l`$ and the wavevectors $`𝐤,𝐤_1,\mathrm{},𝐤_n`$ by a hard gluon carrying of momentum $`𝐩`$ with consequent radiation of $`m`$ plasmons with frequencies $`\omega _{𝐤_1^{}}^l,\mathrm{},\omega _{𝐤_m^{}}^l`$ and the wavevectors $`𝐤_1^{},\mathrm{},𝐤_m^{}`$. The function $`w(𝐩^{}|𝐤_1^{},\mathrm{},𝐤_m^{};𝐤,𝐤_1,\mathrm{},𝐤_n)`$ is the probability of inverse process - the absorption of $`m`$ plasmons by a hard gluon with the momentum $`𝐩^{}𝐩+𝐤+𝐤_1+\mathrm{}+𝐤_n𝐤_1^{}\mathrm{}𝐤_m^{}`$ with consequent radiation of $`n+1`$ plasmons. Diagrammatically this corresponds to a Feynman graph with two hard external lines and an arbitrary number of $`n+m+1`$ soft external lines. By using the fact that $`|𝐩||𝐤|,|𝐤_1|,\mathrm{},|𝐤_n|,|𝐤_1^{}|,\mathrm{},|𝐤_m^{}|`$, the energy conservation law can be represented in the form of following ”generalized” resonance condition:
$$\omega _𝐤^l+\omega _{𝐤_1}^l+\mathrm{}+\omega _{𝐤_n}^l\omega _{𝐤_1^{}}^l\mathrm{}\omega _{𝐤_m^{}}^l𝐯(𝐤+𝐤_1+\mathrm{}+𝐤_n𝐤_1^{}\mathrm{}𝐤_m^{})=0,$$
(2.8)
where $`𝐯𝐩/|𝐩|`$ . Furthermore, one can approximate the distribution function of hard gluons on the r.h.s. of equations (2.5) and (2.6),
$$f_𝐩^{}f_𝐩+(𝐤+𝐤_1+\mathrm{}+𝐤_n𝐤_1^{}\mathrm{}𝐤_m^{})\frac{f_𝐩}{𝐩},$$
and set $`1+f_𝐩1+f_𝐩^{}1`$ by virtue of $`f_𝐩,f_𝐩^{}1`$.
The r.h.s. of (2.5) and (2.6) can be formally considered as an expansions of $`\mathrm{\Gamma }_\mathrm{d}^{(𝒮)}`$ and $`\mathrm{\Gamma }_i^{(𝒮)}`$ in the functional series in powers of the plasmon number density. The actual dimensionless parameter of expansion here, is (for classical statistic) the ratio of the energy of longitudinal plasma excitations to the averaged thermal energy per particle, i.e.
$$\epsilon =\left(\frac{\mathrm{d}𝐤}{(2\pi )^3}\omega _𝐤^lN_𝐤^l\right)/\left(\overline{n}\frac{\mathrm{d}𝐩}{(2\pi )^3}E_𝐩f_𝐩\right),$$
where $`\overline{n}`$ is the mean density. In conditions, when the excitations energy is a small quantity compared with the thermal energy of hard particles, we have
$$\epsilon 1.$$
(2.9)
The last inequality means that the fields of longitudinal oscillations are sufficiently small and they cannot essentially change such ”crude” equilibrium parameters of a plasma as particles density, temperature and thermal energy (this, in particular, justifies the choice of the distribution function of thermal gluons in the form of (2.1)).
On the other hand, however, we shall consider the energy of the plasma oscillations to be sufficiently large, i.e. greatly exceeding the energy of thermal fluctuations of the the color field in the plasma. The consequence of the last requirement is the inequality
$$\epsilon \delta ,$$
(2.10)
where $`\delta `$ is the plasma parameter
$$\delta =\frac{\overline{r}^3}{r_D^3}1.$$
Here, $`\overline{r}`$ is the inter-particle distance $`(\overline{n}^{1/3})`$, and $`r_D`$ is the Debye length
$$r_D^2=\frac{T}{4\pi \overline{n}g^2N_c}.$$
The condition (2.10) demonstrates the validity of ignoring the hard gluon collisions among themselves relative to their interactions with soft plasma modes.
Inequalities (2.9) and (2.10) correspond to the weak turbulent approximation, within the framework of which one can restrict the consideration to several first terms in a functional expansions of $`\mathrm{\Gamma }_\mathrm{d}^{(𝒮)}`$ and $`\mathrm{\Gamma }_i^{(𝒮)}`$. We note however, that when the energy level of the plasma excitations becomes comparable with the thermal energy of particles, e.g. as the result of development of a strong instability (strong turbulence), the perturbation theory here, is no longer applicable and the problem of summation of all the relevant contributions thus appears. The last situation can be really take place in the processes proceed in QGP emerging from the heavy ion collisions at higher energies. In this work we do not consider this very complicated problem, assuming that the inequality (2.9) is always fulfilled.
Let us discuss in more detail the first terms in the expansions of $`\mathrm{\Gamma }_\mathrm{d}^{(𝒮)}`$ and $`\mathrm{\Gamma }_i^{(𝒮)}`$. For $`n=m=0`$, the equation (2.8) results in the relation
$$\omega _𝐤^l\mathrm{𝐯𝐤}=0,$$
which is well-known as the Cherenkov resonance condition, which does not hold in the gluon plasma. Therefore $`\mathrm{\Gamma }_\alpha ^{(𝒮)}=\mathrm{\Gamma }_i^{(𝒮)}=0`$, when only $`O(\epsilon ^0)`$ terms on the r.h.s. of (2.5) and (2.6) are retained.
For $`n=1,m=0`$ we have
$$\omega _𝐤^l+\omega _{𝐤_1}^l𝐯(𝐤+𝐤_1)=0,$$
(2.11)
and for $`n=0,m=1`$, respectively (here, we replace $`𝐤_1^{}`$ by $`𝐤_1`$)
$$\omega _𝐤^l\omega _{𝐤_1}^l𝐯(𝐤𝐤_1)=0.$$
(2.12)
The first resonance condition (2.11) describes the simultaneous radiation (or absorption) of two plasmons with frequencies $`\omega _𝐤^l,\omega _{𝐤_1}^l`$ and wavevectors $`𝐤,𝐤_1`$. By virtue of the fact that the phase velocity of the longitudinal oscillations exceeds the velocity of light, this process is kinematically forbidden. The first non-trivial terms in the expansion of the functionals $`\mathrm{\Gamma }_\mathrm{d}^{(𝒮)}`$ and $`\mathrm{\Gamma }_i^{(𝒮)}`$ are defined by the second resonance condition (2.12). It is associated with the absorption of the plasmon by a hard gluon with frequency $`\omega _𝐤^l`$ and wavevector $`𝐤`$ with its consequent radiation with frequency $`\omega _{𝐤_1}^l`$ and wavevector $`𝐤_1`$ (and vice versa). Schematically this process can be represented as follows:
$$\mathrm{g}^{}+\mathrm{g}\mathrm{g}^{}+\mathrm{g},$$
(2.13)
where $`\mathrm{g}^{}`$ are the plasmon collective excitations and $`\mathrm{g}`$ are excitations with characteristic momenta of order $`T`$. In the theory of the ordinary plasma \[22-24\] this process is known as the nonlinear Landau damping. In the case of QGP it was studied in detail in . We have shown, that the nonlinear Landau damping rate
$$\gamma ^l(𝐤)(\mathrm{\Gamma }_\mathrm{d}^{(𝒮)}[N_𝐤^l]\mathrm{\Gamma }_i^{(𝒮)}[N_𝐤^l])|_{n=0,m=1}$$
(2.14)
defines two processes: the effective pumping of energy across the spectrum towards small wavenumbers with the conservation of excitation energy and properly nonlinear dissipation (damping) of the longitudinal plasma waves by hard particles, where the first process is crucial. The consequence of this fact is the inequality: $`\gamma ^l(0)<0`$, i.e. $`𝐤=0`$ \- mode is increased. The main conclusion, which we drew in is that the only process of the nonlinear Landau damping does not lead to the total relaxation of soft excitations in the gomogeneous isotropic plasma. At the scale of a small $`|𝐤|(|𝐤|gT)`$ it is necessary to consider the processes of higher-order in $`\epsilon `$, than (2.13), which lead to the suppression of increase of the $`𝐤=0`$-mode.
The following terms in the expansion of the decay rate (2.5) and the inverse decay rate (2.6), corresponding to $`n,m=1,2`$, are defined by
$$\omega _𝐤^l\pm \omega _{𝐤_1}^l\pm \omega _{𝐤_2}^l𝐯(𝐤\pm 𝐤_1\pm 𝐤_2)=0,$$
(2.15)
$$\omega _𝐤^l\pm \omega _{𝐤_1}^l\omega _{𝐤_2}^l𝐯(𝐤\pm 𝐤_1𝐤_2)=0.$$
Physically this corresponds to simultaneous absorption (radiation) of three plasmons by a thermal gluon, or simultaneous absorption (radiation) of two plasmons with consequent radiation (absorption) of one plasmon. At the long-wavelength range these processes are kinematically forbidden and therefore in this approximation $`\mathrm{\Gamma }_\mathrm{d}^{(𝒮)}`$ and $`\mathrm{\Gamma }_i^{(𝒮)}`$ vanish.
However there are contributions different from zero in the expansions of the generalized decay rate $`\mathrm{\Gamma }_\mathrm{d}`$ and the inverse $`\mathrm{\Gamma }_i`$ in the equation (2.4) which are of the same order as the processes (2.15), i.e. of order $`O(\epsilon ^2)`$. These contributions are concerned with the second type of the nonlinear processes defining the time-space evolution of $`N_𝐤^l`$ and going without exchange of energy between hard thermal gluons and plasmons. They represent the processes of decays, fusions of plasmons and their scattering off each other. Diagrammatically, this corresponds to Feynman graph, where all external lines are soft. The relevant decay rate $`\mathrm{\Gamma }_i^{(𝒫)}`$ and regenerating rate $`\mathrm{\Gamma }_\mathrm{d}^{(𝒫)}`$ can be formally represented in the form
$$\mathrm{\Gamma }_\mathrm{d}^{(𝒫)}[N_𝐤^l]=\underset{n,m}{}d𝒯_{nm}^{(𝒫)}w(𝐤,𝐤_1,\mathrm{},𝐤_n;𝐤_1^{},\mathrm{},𝐤_m^{})N_{𝐤_1}^l\mathrm{}N_{𝐤_n}^l(1+N_{𝐤_1^{}}^l)\mathrm{}(1+N_{𝐤_m^{}}^l),$$
(2.16)
and
$$\mathrm{\Gamma }_i^{(𝒫)}[N_𝐤^l]=\underset{n,m}{}d𝒯_{nm}^{(𝒫)}w(𝐤_1^{},\mathrm{},𝐤_m^{};𝐤,𝐤_1,\mathrm{},𝐤_n)N_{𝐤_1^{}}^l\mathrm{}N_{𝐤_m^{}}^l(1+N_{𝐤_1}^l)\mathrm{}(1+N_{𝐤_n}^l).$$
(2.17)
Here, the phase-space measure is
$$d𝒯_{nm}^{(𝒫)}\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\mathrm{}\frac{\mathrm{d}𝐤_n}{(2\pi )^3}\frac{\mathrm{d}𝐤_1^{}}{(2\pi )^3}\mathrm{}\frac{\mathrm{d}𝐤_m^{}}{(2\pi )^3}(2\pi )^4\delta ^{(4)}(k+k_1+\mathrm{}+k_nk_1^{}\mathrm{}k_m^{}),$$
(2.18)
with the delta functions expressing the energy and the momentum conservation of the decay processes, the fusion and the plasmon scattering among themselves. The decay rate (2.16) and the regenerating rate (2.17) are different from zero if the following ”resonance” conditions are obeyed
$$\omega _𝐤^l+\omega _{𝐤_1}^l+\mathrm{}+\omega _{𝐤_n}^l\omega _{𝐤_1^{}}^l\mathrm{}\omega _{𝐤_m^{}}^l=0,$$
$$𝐤+𝐤_1+\mathrm{}+𝐤_n𝐤_1^{}\mathrm{}𝐤_m^{}=0.$$
The first contribution is different from zero in $`\mathrm{\Gamma }_\mathrm{d}^{(𝒫)}`$ and $`\mathrm{\Gamma }_i^{(𝒫)}`$, which would arise in the case of $`n=m=1,2`$ as defined by the system of equations
$$\omega _𝐤^l\pm \omega _{𝐤_1}^l\omega _{𝐤_2}^l=0,$$
(2.19)
$$𝐤\pm 𝐤_1𝐤_2=0.$$
These conservation laws describe a decay of one plasmon into two plasmons and the reverse process of fusion of two plasmons into one plasmon,
$$\mathrm{g}^{}\mathrm{g}_1^{}+\mathrm{g}_2^{}.$$
(2.20)
Three-plasmon decay (2.20) is as important as the process of nonlinear Landau damping (2.13). However, the specific peculiarity of a dispersion law of the longitudinal oscillations in hot non-Abelian plasma is that resonance equations (2.19) have no solutions no matter what the values of wavevectors $`𝐤,𝐤_1`$ and $`𝐤_2`$ may be. The processes (2.20) are kinematically forbidden by the conservation laws, and it makes a contribution only in the second order over parameter $`\epsilon `$ of perturbation theory, for $`n,m=1,2,3`$.
In general case, with four plasmons, two different processes occur
$$\mathrm{g}^{}\mathrm{g}_1^{}+\mathrm{g}_2^{}+\mathrm{g}_3^{},$$
(2.21)
$$\mathrm{g}^{}+\mathrm{g}_1^{}\mathrm{g}_2^{}+\mathrm{g}_3^{}.$$
(2.22)
The first of them corresponds to the process of the decay of one plasmon $`\mathrm{g}^{}`$ into three plasmons $`\mathrm{g}_1^{},\mathrm{g}_2^{},\mathrm{g}_3^{}`$, and the reverse process of fusion of three plasmons into one plasmon $`\mathrm{g}^{}`$. The second process presents the plasmon scattering by plasmon. The last one is considered as the decay process and is interpretated as the process of the decay (fusion) of two plasmons $`\mathrm{g}^{}`$ and $`\mathrm{g}_1^{}`$ into two plasmons $`\mathrm{g}_2^{}`$ and $`\mathrm{g}_3^{}`$. For four-plasmon decay process (2.21), the following resonance conditions are obeyed:
$$\omega _𝐤^l\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l=0,$$
(2.23)
$$𝐤𝐤_1𝐤_2𝐤_3=0,$$
and for the scattering process (2.22) we have, respectively,
$$\omega _𝐤^l+\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l=0,$$
(2.24)
$$𝐤+𝐤_1𝐤_2𝐤_3=0.$$
It is not difficult to show that the conservation laws (2.23) and (2.24) kinematically forbid the processes of the direct decay of the plasmon into three (and vice versa) (2.21) and admit only the processes of (2.22) type. As shown in section 9, at the soft scale the latter process is suppressed by a power of $`g`$ relative to the process of nonlinear Landau damping (2.13). However at the ultrasoft scale, one would expect that the process of the elastic scattering of the plasmon by a plasmon may be as larger as the process (2.13) and thus plays an important role in the kinetics of the plasmons at larger wavelength (see the end of section 9).
Putting the expressions (2.16) and (2.17) into equation (2.4), where only $`O(\epsilon ^2)`$ relevant terms in the expansions are retained, we result in the Boltzmann equation, describing four-plasmon decays of (2.22) type
$$\frac{N_𝐤^l}{t}+𝐕_𝐤^l\frac{N_𝐤^l}{𝐱}=$$
(2.25)
$$=\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\frac{\mathrm{d}𝐤_2}{(2\pi )^3}\frac{\mathrm{d}𝐤_3}{(2\pi )^3}(2\pi )^4\delta (\omega _𝐤^l+\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l)\delta (𝐤+𝐤_1𝐤_2𝐤_3)$$
$$\times w(𝐤,𝐤_1;𝐤_2,𝐤_3)\{N_{𝐤_2}^lN_{𝐤_3}^l(1+N_𝐤^l)(1+N_{𝐤_1}^l)N_𝐤^lN_{𝐤_1}^l(1+N_{𝐤_2}^l)(1+N_{𝐤_3}^l)\}$$
$$\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\frac{\mathrm{d}𝐤_2}{(2\pi )^3}\frac{\mathrm{d}𝐤_3}{(2\pi )^3}(2\pi )^4\delta (\omega _𝐤^l+\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l)\delta (𝐤+𝐤_1𝐤_2𝐤_3)$$
$$\times w(𝐤,𝐤_1;𝐤_2,𝐤_3)\{N_{𝐤_1}^lN_{𝐤_2}^lN_{𝐤_3}^l+N_𝐤^lN_{𝐤_2}^lN_{𝐤_3}^lN_𝐤^lN_{𝐤_1}^lN_{𝐤_2}^lN_𝐤^lN_{𝐤_1}^lN_{𝐤_3}^l\}.$$
In writing this equation, we have used the fact that the probabilities of direct and reverse processes are equal and besides in the last line in the semiclassical regime we consider the soft modes to be strongly populated, i.e. $`(1+N_𝐤^l)(1+N_{𝐤_1}^l)N_𝐤^lN_{𝐤_1}^l+N_𝐤^l+N_{𝐤_1}^l`$ etc.. A similar Boltzmann equation for plasmons was studied intensively for the ordinary plasma (in the explicit expression for the function $`w(𝐤,𝐤_1;𝐤_2,𝐤_3)`$ can be found). The main purpose of this work is the derivation in the explicit form of the probability of plasmon-plasmon scattering for hot non-Abelian plasma. The function of $`w(𝐤,𝐤_1;𝐤_2,𝐤_3)`$ must satisfy the symmetry relations over permutation of arguments
$$w(𝐤,𝐤_1;𝐤_2,𝐤_3)=w(𝐤_2,𝐤_3;𝐤_,𝐤_1)=w(𝐤,𝐤_1;𝐤_3,𝐤_2)=w(𝐤_1,𝐤;𝐤_2,𝐤_3),$$
(2.26)
which are the consequence of the indistinguishable of the plasmons (recall that here, we discuss only colorless plasmons, i.e. those for wich the occupation number does not carry adjoint color indices).
In conclusion of this section we note some general properties of the processes of four-plasmon decays. Multiplying the r.h.s. of Eq. (2.25) in turn by $`\omega _𝐤^l`$ and $`𝐤`$, integrating with respect to the wavevector $`𝐤`$, and taking into account (2.26), it is easily checked that
$$\frac{\mathrm{d}𝐤}{(2\pi )^3}\omega _𝐤^lN_𝐤^l=const,𝒦\frac{\mathrm{d}𝐤}{(2\pi )^3}𝐤N_𝐤^l=const.$$
These relations are evident a consequence of the conservation laws of energy and momentum in the plasmon-plasmon scattering.
On the other hand, the total plasmon numbers in such decays should also be conserved, because the only process in which two plasmons decay into two others is permitted. Really, by integrating (2.25) over all $`𝐤`$-space and taking into account the relations (2.26), it is easy to verify that
$$𝒩\frac{\mathrm{d}𝐤}{(2\pi )^3}N_𝐤^l=const,$$
i.e. four-plasmon decay not changes the total number of plasmons.
## 3 The random-phase approximation
In this section let us briefly recall the main methodological notion used in study of the processes of nonlinear wave-interaction in a non-Abelian plasma within semiclassical approximation in .
We use the metric $`g^{\mu \nu }=diag(1,1,1,1)`$ and choose units such that $`c=k_B=1`$. The gauge field potentials are $`N_c\times N_c`$-matrices in a color space defined by $`A_\mu =A_\mu ^at^a`$ with $`N_c^21`$ Hermitian generators of the $`SU(N_c)`$ group in the fundamental representation. The field strength tensor $`F_{\mu \nu }=F_{\mu \nu }^at^a`$ with
$$F_{\mu \nu }^a=_\mu A_\nu ^a_\nu A_\mu ^a+gf^{abc}A_\mu ^bA_\nu ^c$$
(3.1)
obeys the Yang-Mills (YM) equation in a covariant gauge
$$_\mu F^{\mu \nu }(X)ig[A_\mu (X),F^{\mu \nu }(X)]\xi ^1^\nu ^\mu A_\mu (X)=j^\nu (X),$$
(3.2)
where $`\xi `$ is a gauge parameter. $`j^\nu `$ is a color current
$$j^\nu =gt^ad^4pp^\nu \mathrm{Tr}(T^af),$$
(3.3)
where $`T^a`$ are Hermitian generators of $`SU(N_c)`$ group in the adjoint representation $`((T^a)^{bc}=if^{abc},\mathrm{Tr}(T^aT^b)=N_c\delta ^{ab})`$. The distribution function of gluons $`f`$ satisfies the dynamical equation which in the semiclassical limit (when polarization effects are neglected ) is
$$p^\mu \stackrel{~}{𝒟}_\mu f+\frac{1}{2}gp^\mu \{_{\mu \nu },\frac{f}{p_\nu }\}=0,$$
(3.4)
where $`\stackrel{~}{𝒟}_\mu `$ is a covariant derivative acting as
$$\stackrel{~}{𝒟}_\mu =_\mu ig[𝒜_\mu (X),],$$
with \[ , \] and $`\{,\}`$ denoting the commutator and anticommutator, respectively, and $`𝒜_\mu `$, $`_{\mu \nu }`$ are defined as $`𝒜_\mu =A_\mu ^aT^a,_{\mu \nu }=F_{\mu \nu }^aT^a`$.
The distribution function $`f`$ can be decomposed into two parts: regular and random, where the latter is generated by spontaneous fluctuations in the plasma
$$f=f^R+f^T,$$
(3.5)
so that
$$f=f^R,f^T=0.$$
Here, angular brackets $``$ indicate a statistical ensemble of averaging. The initial values of parameters which characterize the collective degree of plasma freedom is a such statistical ensemble. For almost linear collective motion to be considered below this may be initial values of oscillation phases.
We also use the definition
$$A_\mu =A_\mu ^R+A_\mu ^T,A_\mu ^T=0.$$
(3.6)
The regular (background) part of the field $`A_\mu ^R`$ will be considered to be equal to zero. The condition for which the last assumption holds, will be considered closely in the next section.
By averaging equation (3.4) over statistical ensemble, we obtain the kinetic equation for the regular part of the distribution function of hard gluons $`f^R`$
$$p^\mu _\mu f^R=igp^\mu [𝒜_\mu ^T,f^T]\frac{1}{2}gp^\mu \{(_{\mu \nu }^T)_L,\frac{f^T}{p_\nu }\}\frac{1}{2}gp^\mu \{(_{\mu \nu }^T)_{NL},\frac{f^R}{p_\nu }\}$$
$$\frac{1}{2}gp^\mu \{(_{\mu \nu }^T)_{NL},\frac{f^T}{p_\nu }\}.$$
(3.7)
Here, the indices $`L`$ and $`NL`$ denote the linear and nonlinear parts with respect to field $`A_\mu ^a`$ of the strength tensor (3.1). The correlation functions on the r.h.s. of this equation are collision terms due to particle-wave interactions and describe the backreaction of the background state from the plasma waves.
We assume that the typical time the nonlinear relaxation for the oscillations is a small quantity relative to the time scale over which the distribution of hard transverse gluons $`f^R`$ vary substantially. Therefore we neglect by change of the regular part of the distribution function with space and time, assuming that this function is specified and describes the global equilibrium in the gluon plasma
$$f^Rf^0=2\frac{2\theta (p_0)}{(2\pi )^3}\delta (p^2)\frac{1}{\mathrm{e}^{(pu)/T}1},$$
(3.8)
where $`u_\mu `$ is the 4-velocity of the plasma. (Here, for convenience, we somewhat overdetermine the equilibrium distribution function of thermal gluons (2.1).)
We use the expansion in powers of the oscillations amplitude of the random function $`f^T`$ to investigate non-equilibrium processes in QGP, such that the excitation energy of waves is small quantity in relation to the total energy of the particles
$$f^T=\underset{n=1}{\overset{\mathrm{}}{}}f^{T(n)},$$
(3.9)
where $`f^{T(n)}`$ collects the contributions of the $`n`$-th power in $`A_\mu ^T`$. The expansion of a color current, corresponding to (3.9) has the form
$$j_\mu =j_\mu ^R+j_\mu ^T,j_\mu =j_\mu ^R,j_\mu ^T=\underset{n=1}{\overset{\mathrm{}}{}}j_\mu ^{T(n)},$$
(3.10)
where by the definition (3.3), we have
$$j_\mu ^{T(n)}=gt^ad^4pp_\mu \mathrm{Tr}(T^af^{T(n)}).$$
(3.11)
The regular part of a current vanishes for the global equilibrium gluon plasma.
Substituting the expansion (3.9) into (3.4), and collecting the terms of the same order in $`A_\mu ^T`$, we derive the system of equations
$$p^\mu _\mu f^{T(1)}=\frac{1}{2}gp^\mu \{(_{\mu \nu }^T)_L,\frac{f^R}{p_\nu }\},$$
(3.12)
$$p^\mu _\mu f^{T(2)}=igp^\mu ([𝒜_\mu ^T,f^{T(1)}][𝒜_\mu ^T,f^{T(1)}])$$
(3.13)
$$\frac{1}{2}gp^\mu (\{(_{\mu \nu }^T)_L,\frac{f^{T(1)}}{p_\nu }\}\{(_{\mu \nu }^T)_L,\frac{f^{T(1)}}{p_\nu }\})\frac{1}{2}gp^\mu \{(_{\mu \nu }^T)_{NL}(_{\mu \nu }^T)_{NL},\frac{f^R}{p_\nu }\},$$
$$p^\mu _\mu f^{T(3)}=igp^\mu ([𝒜_\mu ^T,f^{T(2)}][𝒜_\mu ^T,f^{T(2)}])\frac{1}{2}gp^\mu (\{(_{\mu \nu }^T)_L,\frac{f^{T(2)}}{p_\nu }\}$$
(3.14)
$$\{(_{\mu \nu }^T)_L,\frac{f^{T(2)}}{p_\nu }\})\frac{1}{2}gp^\mu (\{(_{\mu \nu }^T)_{NL},\frac{f^{T(1)}}{p_\nu }\}\{(_{\mu \nu }^T)_{NL},\frac{f^{T(1)}}{p_\nu }\})etc..$$
We rewrite the Yang-Mills equation (3.2), connecting a gauge field with a color current, in the following form
$$_\mu (F^{T\mu \nu })_L\xi ^1^\nu ^\mu A_\mu ^T+j^{T(1)\nu }=$$
(3.15)
$$=j_{NL}^{T\nu }+ig_\mu ([A^{T\mu },A^{T\nu }][A^{T\mu },A^{T\nu }])$$
$$+ig([A_\mu ^T,(F^{T\mu \nu })_L][A_\mu ^T,(F^{T\mu \nu })_L])+g^2([A_\mu ^T,[A^{T\mu },A^{T\nu }]][A_\mu ^T,[A^{T\mu },A^{T\nu }]]).$$
Here, on the l.h.s. we collect all linear terms with respect to $`A_\mu ^T`$ and we denote: $`j_{NL}^{T\nu }j^{T(2)\nu }+j^{T(3)\nu }+\mathrm{}.`$ It will be shown that at leading order in $`g`$ only, the first two terms in $`j_{NL}^{T\nu }`$ need to be kept.
It is not difficult to obtain the explicit form of the terms in the color current expansion (3.10) from the system of equations (3.12)-(3.14) and the relation (3.11). In the momentum representation and to the leading order in coupling constant up to the third-order terms, we have
$$j^{T(1)\mu }(k)=\mathrm{\Pi }^{\mu \nu }(k)A_\nu ^T(k),$$
(3.16)
where
$$\mathrm{\Pi }^{\mu \nu }(k)=g^2N_cd^4p\frac{p^\mu (p^\nu (k_p)(kp)_p^\nu )f^0}{pk+ip_0ϵ},ϵ+0$$
(3.17)
is the high temperature polarization tensor;
$$j^{T(2)a\mu }(k)=f^{abc}S_{k,k_1,k_2}^{(II)\mu \nu \lambda }(A_\nu ^b(k_1)A_\lambda ^c(k_2)A_\nu ^b(k_1)A_\lambda ^c(k_2))\delta (kk_1k_2)𝑑k_1𝑑k_2,$$
(3.18)
where
$$S_{k,k_1,k_2}^{(II)\mu \nu \lambda }=ig^3N_cd^4p\frac{p^\mu p^\nu p^\lambda }{pk+ip_0ϵ}\frac{(k_2_pf^0)}{pk_2+ip_0ϵ};$$
(3.19)
$$j^{T(3)a\mu }(k)=f^{abf}f^{fde}\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{(II)\mu \nu \lambda \sigma }(A_\nu ^b(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)A_\nu ^b(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)$$
$$A_\nu ^b(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2))\delta (kk_1k_2k_3)dk_1dk_2dk_3,$$
(3.20)
and
$$\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{(II)\mu \nu \lambda \sigma }=g^4N_cd^4p\frac{p^\mu p^\nu p^\lambda p^\sigma }{pk+ip_0ϵ}\frac{1}{p(k_1+k_2)+ip_0ϵ}\frac{(k_2_pf^0)}{pk_2+ip_0ϵ}.$$
(3.21)
For simplicity, hereafter we drop the superscript $`T`$ of a gauge field.
Furthermore, we rewrite the equation (3.15) in the momentum representation. By inserting the linear part of the random current (3.16) and nonlinear corrections (3.18) and (3.20) into Eq. (3.15), one finds
$$[k^2g^{\mu \nu }(1+\xi ^1)k^\mu k^\nu \mathrm{\Pi }^{\mu \nu }(k)]A_\nu ^b(k)$$
$$=f^{bcd}S_{k,k_1,k_2}^{\mu \nu \lambda }(A_\nu ^c(k_1)A_\lambda ^d(k_2)A_\nu ^c(k_1)A_\lambda ^d(k_2))\delta (kk_1k_2)𝑑k_1𝑑k_2$$
(3.22)
$$+f^{bcf}f^{fde}\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }(A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2))$$
$$\times \delta (kk_1k_2k_3)dk_1dk_2dk_3.$$
Here, $`S_{k,k_1,k_2}^{\mu \nu \lambda }S_{k,k_1,k_2}^{(I)\mu \nu \lambda }+S_{k,k_1,k_2}^{(II)\mu \nu \lambda },\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{(I)\mu \nu \lambda \sigma }+\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{(II)\mu \nu \lambda \sigma }`$. The functions $`S^{(II)}`$ and $`\mathrm{\Sigma }^{(II)}`$ are defined by expressions (3.19) and (3.21), respectively, and
$$S_{k,k_1,k_2}^{(I)\mu \nu \lambda }=ig(k^\nu g^{\mu \lambda }+k_2^\nu g^{\mu \lambda }k_2^\mu g^{\nu \lambda }),\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{(I)\mu \nu \lambda \sigma }=g^2g^{\nu \lambda }g^{\mu \sigma }.$$
(3.23)
These tensor structures are caused by self-action of a gauge field. They are defined by nonlinear terms on the r.h.s. of equation (3.15), which are not associated with a color current in QGP. In section 5 equation (3.22) will be formally solved by iteration.
In we introduce the correlation function of random oscillations
$$I_{\mu \nu }^{ab}(k^{},k)=A_\mu ^a(k^{})A_\nu ^b(k).$$
(3.24)
In order not to over burden equations by the symbol $`\mathrm{"}\mathrm{"}`$ we use a dagger $`\mathrm{"}\mathrm{"}`$ to denote complex conjugation. In thermal equilibrium, when the correlation function (3.24) in the coordinate representation depends only on the relative coordinates and time $`\mathrm{}X=X^{}X`$, we have
$$I_{\mu \nu }^{ab}(k^{},k)=I_{\mu \nu }^{ab}(k^{})\delta (k^{}k).$$
(3.25)
Off-equilibrium perturbations which are slowly varying in space and time lead to a delta function broadering, and $`I_{\mu \nu }^{ab}`$ depends on both arguments $`k`$ and $`k^{}`$.
Let us introduce $`I_{\mu \nu }^{ab}(k^{},k)=I_{\mu \nu }^{ab}(k,\mathrm{}k)`$, $`\mathrm{}k=k^{}k`$ with $`\mathrm{}k/k1`$ and insert the spectral intensity function in the Wigner form
$$I_{\mu \nu }^{ab}(k,x)=I_{\mu \nu }^{ab}(k,\mathrm{}k)\mathrm{e}^{i\mathrm{}kx}𝑑\mathrm{}k,$$
depending slowly on $`x`$.
Now we multiple equation (3.22) by $`A_\mu ^a(k^{})`$, subtract to it the complex-conjugated equation (with the replacement $`kk^{},ab`$) and average the equation using the formula (3.24). Furthermore we expand the polarization tensor into Hermitian and anti-Hermitian parts
$$\mathrm{\Pi }^{\nu \sigma }(k)=\mathrm{\Pi }^{H\nu \sigma }(k)+\mathrm{\Pi }^{A\nu \sigma }(k),\mathrm{\Pi }^{H\nu \sigma }(k)=\mathrm{\Pi }^{H\sigma \nu }(k),\mathrm{\Pi }^{A\nu \sigma }(k)=\mathrm{\Pi }^{A\sigma \nu }(k).$$
The term with $`\mathrm{\Pi }^A`$ corresponds to linear Landau damping. As was shown by Heinz and Siemens , that linear Landau damping for waves in QGP is absent and hence, this term vanishes. We expand the remaining terms on the l.h.s. in powers of $`\mathrm{}k`$ and keep only linear ones. This corresponds to $`\mathrm{𝑎𝑔𝑟𝑎𝑑𝑖𝑒𝑛𝑡}\mathrm{𝑒𝑥𝑝𝑎𝑛𝑠𝑖𝑜𝑛}`$ procedure, usually used in the derivation of kinetic equations. Multiplying the difference equation by $`\mathrm{e}^{i\mathrm{}kx}`$ and integrating over $`\mathrm{}k`$ with regard to
$$\mathrm{}k_\lambda I_{\mu \nu }^{ab}(k,\mathrm{}k)\mathrm{e}^{i\mathrm{}kx}𝑑\mathrm{}k=i\frac{I_{\mu \nu }^{ab}(k,x)}{x^\lambda },$$
we finally obtain the equation, which is a starting point for our further research
$$\frac{}{k_\lambda }[k^2g^{\mu \nu }(1+\xi ^1)k^\mu k^\nu \mathrm{\Pi }^{H\mu \nu }(k)]\frac{I_{\mu \nu }^{ab}}{x^\lambda }=$$
(3.26)
$$=idk^{}\{f^{bcd}S_{k,k_1,k_2}^{\mu \nu \lambda }A_\mu ^a(k^{})A_\nu ^c(k_1)A_\lambda ^d(k_2)\delta (kk_1k_2)dk_1dk_2$$
$$+f^{bcf}f^{fde}\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }(A_\mu ^a(k^{})A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)A_\mu ^a(k^{})A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2))$$
$$\times \delta (kk_1k_2k_3)dk_1dk_2dk_3(ab,kk^{},compl.conj.)\}.$$
## 4 Consistency with gauge symmetry
In this section we shall discuss the consistency of the approximation scheme which we use with the requirement of the non-Abelian gauge symmetry.
The initial dynamical equation (3.4) and the Yang-Mills equation (3.2) (without the gauge-fixing condition) transform covariantly under local transformations
$$\overline{A}_\mu (X)=h(X)(A_\mu (X)+\frac{i}{g}_\mu )h^{}(X),h(X)=\mathrm{exp}(i\theta ^a(X)t^a)$$
with the parameter $`\theta ^a(X)`$. We also have transformation of gluon distribution function
$$\overline{f}(p,X)=H(X)f(p,X)H^{}(X),$$
where $`H^{ab}(X)=\mathrm{Sp}[t^ah(X)t^bh^{}(X)]`$.
As is known (see, e.g. ), after the splitting (3.5), (3.6) the resulting equations left two symmetries: the $`\mathrm{𝑏𝑎𝑐𝑘𝑔𝑟𝑜𝑢𝑛𝑑}\mathrm{𝑔𝑎𝑢𝑔𝑒}\mathrm{𝑠𝑦𝑚𝑚𝑒𝑡𝑟𝑦}`$,
$$\overline{A}_\mu ^R(X)=h(X)(A_\mu ^R(X)+\frac{i}{g}_\mu )h^{}(X),\overline{A}_\mu ^T(X)=h(X)A_\mu ^T(X)h^{}(X),$$
(4.1)
and the $`\mathrm{𝑓𝑙𝑢𝑐𝑡𝑢𝑎𝑡𝑖𝑜𝑛}\mathrm{𝑔𝑎𝑢𝑔𝑒}\mathrm{𝑠𝑦𝑚𝑚𝑒𝑡𝑟𝑦}`$,
$$\overline{A}_\mu ^R(X)=0,\overline{A}_\mu ^T(X)=h(X)(A_\mu ^R(X)+A_\mu ^T(X)+\frac{i}{g}_\mu )h^{}(X).$$
(4.2)
The condition which we impose on a regular part of the gauge field $`A_\mu ^R`$ in the preceding section and the requirement that the statistical average of the fluctuation vanishes $`A_\mu ^T=0`$, break down both types of symmetry (4.1) and (4.2). Thus in the case of a gauge transformation (4.1) we obtain $`\overline{A}_\mu ^R0`$, and in the case of (4.2) we arrive at non-invariance of the constraint $`A_\mu ^T=0`$. Moreover, the introduced correlation function (3.24) also has an explicitly gauge non-covariant character. This leads to the fact that calculations in the previous section are gauge non-covariant, and therefore the value of these manipulations is doubtful.
Nevertheless, there is a special case, when the preceding (and following) conclusions are justified. This is a case of colorless fluctuation, where $`I_{\mu \nu }^{ab}(k,x)=\delta ^{ab}I_{\mu \nu }(k,x)`$. We can obtain a gauge-invariant equation for $`I_{\mu \nu }(k,x)`$ only in this restriction, in spite of the fact that the intermediate calculations spoil non-Abelian gauge symmetry of the initial equations (3.2)-(3.4).
In principle, we shall be able to maintain an explicit background gauge symmetry (4.1) at each step of our calculations, as has been done, for example, by Blaizot and Iancu for derivation of the Boltzmann equation describing the relaxation of ultrasoft color excitations. First of all we assume that $`A_\mu ^R0`$. Then as the gauge-fixing condition for the random field $`A_\mu ^T`$, we choose the background field gauge
$$𝒟_\mu ^R(X)A^{T\mu }(X)=0,𝒟_\mu ^R(X)_\mu igA_\mu ^R(X),$$
(4.3)
which is manifestly covariant with respect to gauge transformations of the background gauge field $`A_\mu ^R(X)`$. Lastly we define a gauge-covariant Wigner function as in $`[13,\mathrm{\hspace{0.17em}18}]`$
$$\stackrel{´}{I}_{\mu \nu }^{ab}(k,x)=\stackrel{´}{I}_{\mu \nu }^{ab}(s,x)\mathrm{e}^{iks}ds,sX_1X_2,x\frac{1}{2}(X_1+X_2),$$
where
$$\stackrel{´}{I}_{\mu \nu }^{ab}(s,x)U^{aa^{}}(x,x+\frac{s}{2})I^{a^{}b^{}}(x+\frac{s}{2},x\frac{s}{2})U^{b^{}b}(x\frac{s}{2},x),$$
instead of the usual Wigner function $`I_{\mu \nu }^{ab}(k,x)`$, whose ‘poor’ transformation properties follow from the initial definition $`I_{\mu \nu }^{ab}(X_1,X_2)=A_\mu ^{Ta}(X_1)A_\nu ^{Tb}(X_2)`$. The function $`U(x,y)`$ is the non-Abelian parallel transporter
$$U(x,y)=\mathrm{P}\mathrm{exp}\left\{ig_\gamma dz^\mu A_\mu ^R(z)\right\}.$$
The path $`\gamma `$ is the straight line joining $`x`$ and $`y`$.
The derivation of the kinetic equation for plasmons in this approach becomes quite cumbersome and non-trivial. For example, on the l.h.s. of the equations for a random part of the distribution (3.12)-(3.14), the covariant derivative $`𝒟_\mu ^R`$ will be used instead of the ordinary one $`_\mu `$. Besides, we cannot assume that the regular part of the distribution function is specified and equal to the Bose-Einstein distribution (3.8). It is necessary to also take into account their change using the kinetic equation (3.7) with the collision terms on the r.h.s. of (3.7), which describes the backreaction of the background distributions from the soft fluctuations. The correlators on the r.h.s. of equation (3.7) can be expressed in terms of the function $`\stackrel{´}{I}_{\mu \nu }^{ab}`$ and the distribution of hard transverse gluons $`f^R(p,X)`$ only.
However, if we restrict our consideration to the study of colourless excitations and replace the distribution function of hard gluons by its equilibrium value (3.8), then this leads to an effective vanishing of terms with mean field $`A_\mu ^R`$. This follows, for example, from the analysis of the derivation of the Boltzmann equation by Blaizot and Iancu . Therefore, the simplest way to derive the kinetic equations for soft colorless QGP excitations is to assume $`A_\mu ^R=0`$ and to use the prime gauge non-covariant correlator (3.24). In this case the background field gauge (4.3) is reduced to a covariant one. The resulting Boltzmann equation for colorless plasmons will be gauge invariant if all contributions to the probability of plasmon-plasmon scattering at the leading order in $`g`$ are taken into account (see also discussion in conclusion).
## 5 The interacting fields as the functions of free fields
Let us define approximate solution of equation (3.22) accurate up to third-order in the oscillations amplitude of the free field. For this purpose, it is convenient to write this equation in a more compact form.
We introduce the following notation
$$J^{T(2)a\mu }(k)f^{abc}S_{k,k_1,k_2}^{\mu \nu \lambda }(A_\nu ^b(k_1)A_\lambda ^c(k_2)A_\nu ^b(k_1)A_\lambda ^c(k_2))\delta (kk_1k_2)𝑑k_1𝑑k_2,$$
(5.1)
$$J^{T(3)a\mu }(k)f^{abf}f^{fde}\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }(A_\nu ^b(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)A_\nu ^b(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)$$
$$A_\nu ^b(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2))\delta (kk_1k_2k_3)dk_1dk_2dk_3.$$
(5.2)
The expressions (5.1) and (5.2) present the nonlinear color currents, including the self-action effects of gauge fields, in contrast to (3.18) and (3.20).
Using $`{}_{}{}^{}𝒟_{\mu \nu }^{}(k)`$ we denote the medium modified (retarded) gluon propagator, which in a covariant gauge has a form
$${}_{}{}^{}𝒟_{\mu \nu }^{}(k)=P_{\mu \nu }(k)^{}\mathrm{\Delta }^t(k)Q_{\mu \nu }(k)^{}\mathrm{\Delta }^l(k)+\xi D_{\mu \nu }(k)\mathrm{\Delta }^0(k),$$
(5.3)
where $`{}_{}{}^{}\mathrm{\Delta }_{}^{t,l}(k)=1/(k^2\mathrm{\Pi }^{t,l}(k)),\mathrm{\Pi }^t(k)=\frac{1}{2}\mathrm{\Pi }^{\mu \nu }(k)P_{\mu \nu }(k),\mathrm{\Pi }^l(k)=\mathrm{\Pi }^{\mu \nu }(k)Q_{\mu \nu }(k);\mathrm{\Delta }^0(k)=1/k^2`$. The Lorentz matrices in (5.3) are members of the basis
$$P_{\mu \nu }(k)=g_{\mu \nu }D_{\mu \nu }(k)Q_{\mu \nu }(k),Q_{\mu \nu }(k)=\frac{\overline{u}_\mu (k)\overline{u}_\nu (k)}{\overline{u}^2(k)},C_{\mu \nu }(k)=\frac{(\overline{u}_\mu (k)k_\nu +\overline{u}_\nu (k)k_\mu )}{\sqrt{2k^2\overline{u}^2(k)}},$$
$$D_{\mu \nu }=k_\mu k_\nu /k^2,\overline{u}_\mu (k)=k^2u_\mu k_\mu (ku).$$
(5.4)
Let us assume that we are in the rest frame of a heat bath, so that $`u_\mu =(1,0,0,0).`$
Using the above introduced functions, the equation (3.22) can be rewritten in the form
$${}_{}{}^{}𝒟_{}^{1\mu \nu }(k)A_\nu ^a(k)=J^{T(2)a\mu }(A,A)J^{T(3)a\mu }(A,A,A).$$
(5.5)
The nonlinear integral equation (5.5) is solved by the approximation scheme method - $`\mathrm{𝑡ℎ𝑒𝑤𝑒𝑎𝑘}\mathrm{𝑓𝑟𝑒𝑒}\mathrm{𝑓𝑖𝑒𝑙𝑑}\mathrm{𝑒𝑥𝑝𝑎𝑛𝑠𝑖𝑜𝑛}`$ (small perturbations). Discarding the nonlinear terms in $`A`$ on the r.h.s. of equation (5.5), we obtain in the first approximation
$${}_{}{}^{}𝒟_{}^{1\mu \nu }(k)A_\nu ^a(k)=0.$$
The solution of this equation, which we denote by $`A_\mu ^{(0)a}(k)`$ is the solution for free fields.
Further keeping the term, quadratic in the field on the r.h.s. of Eq. (5.5), we derive the equation
$${}_{}{}^{}𝒟_{}^{1\mu \nu }(k)A_\nu ^a(k)=J^{T(2)a\mu }(A^{(0)},A^{(0)}),$$
where on the r.h.s. we substitute free fields instead of interacting ones. The general solution of the last equation can be given in the form
$$A_\mu ^a(k)=A_\mu ^{(0)a}(k)^{}𝒟_{\mu \nu }(k)J^{T(2)a\nu }(A^{(0)},A^{(0)}).$$
This approximate solution was used in research of nonlinear plasmon damping in QGP .
The following term in the expansion of interacting fields is defined from equation
$${}_{}{}^{}𝒟_{}^{1\mu \nu }(k)A_\nu ^a(k)=J^{T(2)a\mu }(^{}𝒟J^{T(2)}(A^{(0)},A^{(0)}),A^{(0)})J^{T(2)a\mu }(A^{(0)},^{}𝒟J^{T(2)}(A^{(0)},A^{(0)}))$$
$$J^{T(3)a\mu }(A^{(0)},A^{(0)},A^{(0)}).$$
(5.6)
Using the explicit expressions for the currents (5.1) and (5.2), after cumbersome algebraic transformations, we obtain the form of the interacting field from the equation (5.6) with the accuracy required for our further calculations
$$A_\mu ^a(k)=A_\mu ^{(0)a}(k)^{}𝒟_{\mu \nu }(k)J^{T(2)a\nu }(A^{(0)},A^{(0)})^{}𝒟_{\mu \nu }(k)\stackrel{~}{J}^{T(3)a\nu }(A^{(0)},A^{(0)},A^{(0)}).$$
(5.7)
Here, the third-order color current on the r.h.s. is defined by the expression
$$\stackrel{~}{J}^{T(3)a\nu }(A^{(0)},A^{(0)},A^{(0)})f^{abf}f^{fde}\stackrel{~}{\mathrm{\Sigma }}_{k,k_1,k_2,k_3}^{\nu \lambda \sigma \rho }(A_\lambda ^{(0)b}(k_3)A_\sigma ^{(0)d}(k_1)A_\rho ^{(0)e}(k_2)$$
$$A_\lambda ^{(0)b}(k_3)A_\sigma ^{(0)d}(k_1)A_\rho ^{(0)e}(k_2))\delta (kk_1k_2k_3)dk_1dk_2dk_3,$$
(5.8)
where
$$\stackrel{~}{\mathrm{\Sigma }}_{k,k_1,k_2,k_3}^{\nu \lambda \sigma \rho }\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{\nu \lambda \sigma \rho }\frac{1}{2}^{}𝒟_{\delta \gamma }(k_1+k_2)(S_{k,k_3,k_1+k_2}^{\nu \lambda \delta }S_{k,k_1+k_2,k_3}^{\nu \delta \lambda })(S_{k_1+k_2,k_1,k_2}^{\gamma \sigma \rho }S_{k_1+k_2,k_2,k_1}^{\gamma \rho \sigma }),$$
(5.9)
and we take into account, that the third-order correlation function $`A^{(0)}A^{(0)}A^{(0)}`$ vanishes by virtue of the fact that $`A^{(0)}`$ represents the amplitude fully non-correlative gauge fields. The factor of $`\frac{1}{2}`$, in front of the second term on the r.h.s. of (5.9) arises from symmetrization with respect to permutation of the potentials $`A_\sigma ^{(0)d}(k_1)`$ and $`A_\rho ^{(0)e}(k_2)`$ in the expression (5.8). The current (5.8) may be interpreted as a certain third-order effective color current, in contrast to the initial ‘bare’ expression (5.2).
## 6 The correspondence principle
For the determination of the probability of plasmon-plasmon scattering in a gluon plasma the method developed in the theory of the nonlinear processes in electron-ion plasma and known as the correspondence principle , is usable. For the non-Abelian plasma this approach is especially effective in the temporal gauge, when we have closer correspondence with the electrodynamics of an ordinary plasma. The gist of this method is as follows.
The change in the plasmon numbers, caused by spontaneous processes of four-plasmon decays only, is
$$\left(\frac{N_𝐤^l}{t}+𝐕_𝐤^l\frac{N_𝐤^l}{𝐱}\right)_4^{sp}=\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\frac{\mathrm{d}𝐤_2}{(2\pi )^3}\frac{\mathrm{d}𝐤_3}{(2\pi )^3}$$
(6.1)
$$\times (2\pi )^4\delta (\omega _𝐤^l+\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l)\delta (𝐤+𝐤_1𝐤_2𝐤_3)w(𝐤,𝐤_1;𝐤_2,𝐤_3)N_{𝐤_1}^lN_{𝐤_2}^lN_{𝐤_3}^l.$$
This equation follows from equation (2.25) in the limit of small intensity $`N_𝐤^l0`$. In this case the change of energy of the longitudinal excitations is
$$\left(\frac{\mathrm{d}}{\mathrm{d}t}\right)_4^{sp}=\frac{\mathrm{d}𝐤}{(2\pi )^3}\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\frac{\mathrm{d}𝐤_2}{(2\pi )^3}\frac{\mathrm{d}𝐤_3}{(2\pi )^3}$$
(6.2)
$$\times (2\pi )^4\delta (\omega _𝐤^l+\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l)\delta (𝐤+𝐤_1𝐤_2𝐤_3)\omega _𝐤^lw(𝐤,𝐤_1;𝐤_2,𝐤_3)N_{𝐤_1}^lN_{𝐤_2}^lN_{𝐤_3}^l.$$
On the other hand the value $`(\mathrm{d}/\mathrm{d}t)_4^{sp}`$ represents the emitted radiant power of the longitudinal waves $`^l`$, which in turn is equal to the work done by the radiation field with the color current, creating it, in unit time
$$^l=d𝐱𝐄^a(𝐱,t)𝐉^a(𝐱,t)=\frac{\mathrm{d}\omega }{2\pi }\frac{\mathrm{d}\omega ^{}}{2\pi }\frac{\mathrm{d}𝐤}{(2\pi )^3}𝐄_{𝐤,\omega ^{}}^a𝐉_{𝐤,\omega }^a\mathrm{e}^{i(\omega ^{}\omega )t}=$$
(6.3)
$$=\frac{1}{2}\frac{\mathrm{d}\omega }{2\pi }\frac{\mathrm{d}\omega ^{}}{2\pi }\frac{\mathrm{d}𝐤}{(2\pi )^3}i\left(\frac{1}{\omega ^{}\epsilon ^l(\omega ^{},𝐤)}\frac{1}{\omega \epsilon ^l(\omega ,𝐤)}\right)\frac{k^ik^j}{𝐤^2}J_{𝐤,\omega }^{ai}J_{𝐤,\omega ^{}}^{aj}\mathrm{e}^{i(\omega ^{}\omega )t}.$$
Here, $`E^{ai}(𝐱,t)=A^{ai}(𝐱,t)/t`$ is chromoelectric field in the temporal gauge. The sign on the r.h.s. of (6.3) corresponds to the choice of sign in front of the current in the Yang-Mills equation (3.2). In conclusion of the last line (6.3) we take into account that the Fourier-component of a field $`𝐄_{𝐤,\omega }^a=𝐤E_{𝐤,\omega }^a/|𝐤|`$ is associated with $`𝐉_{𝐤,\omega }^a`$ by the Yang-Mills equation
$$E_{𝐤,\omega }^a=\frac{i}{\omega \epsilon ^l(\omega ,𝐤)}\frac{(𝐤𝐉_{𝐤,\omega }^a)}{|𝐤|}.$$
In order to define the probability of the four-plasmon decays, the correlation function on the r.h.s. of equation (6.3) has to contain terms of six-order in the free field $`A^{(0)}`$. The required sixth-order correlator yields the color current $`J^{T(3)ai}`$ (5.2) (more precisely, its expression in the temporal gauge). However here, it is also necessary to take into account the effects, that arise from iteration of the current $`J^{T(2)ai}`$ (5.1). Defining in this way all necessary contributions, making the correlation decoupling of the sixth-order correlators in terms of pairs and next expressing next $`A^{(0)}A^{(0)}`$ in terms of $`N^l`$, we obtain an emitted radiant power $`^l`$. Comparing $`^l`$ with (6.2), one identifies the required probability $`w(𝐤,𝐤_1;𝐤_2,𝐤_3)`$.
However this method encountered certain difficulties in deciding on the other gauges, e.g. the covariant gauge. Here, it is convenient for the definition of the probability of plasmon-plasmon scattering to start immediately from equation (3.26). Using the above-obtained expression (5.7) for the potentials of interacting fields, by a simple search we extract all the sixth-order correlators responsible for four-plasmon decays of the type (2.22). The rule used to choose the relevant terms is defined by the simple fact that when we make the correlation decoupling of the sixth-order correlators in the terms of the pairs (in order to define the product $`N_{𝐤_1}^lN_{𝐤_2}^lN_{𝐤_3}^l`$) as a factor, the $`\delta `$-functions arise in the form
$$\delta (\omega _𝐤^l+\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l)\delta (𝐤+𝐤_1𝐤_2𝐤_3).$$
(6.4)
Also the coefficient function preceding $`N_{𝐤_1}^lN_{𝐤_2}^lN_{𝐤_3}^l`$ must satisfy properties (2.26). As will be shown below, these conditions are sufficient to calculate the plasmon-plasmon scattering probability in the covariant gauge (this rather cumbersome and physically not quite transparent approach is suited for other gauges).
At first, we consider the contribution to the r.h.s. of the initial equation (3.26), associated with $`\mathrm{\Sigma }`$-functions. Here, for convenience of further reference we write it separately
$$idk^{}\{f^{bcf}f^{fde}\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }(A_\mu ^a(k^{})A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)A_\mu ^a(k^{})A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2))$$
$$\times \delta (kk_1k_2k_3)dk_1dk_2dk_3(ab,kk^{},compl.conj.)\}.$$
(6.5)
One can obtain the sixth-order terms in the free field by two waves. The first one is as follows. We substitute in turn the expression from the r.h.s. of (5.7), which contains only cubic terms in the potentials of free fields, instead of each potential of the interacting ones, i.e.
$$A_\mu ^a(k)^{}𝒟_{\mu \mu ^{}}(k)\stackrel{~}{J}^{T(3)a\mu ^{}}(A^{(0)},A^{(0)},A^{(0)}).$$
We replace the remaining potentials by the rule $`A_\nu ^c(k_3)A_\nu ^{(0)c}(k_3)`$, etc. The second way is to substitute the quadratic term in $`A^{(0)}`$ from the r.h.s. of (5.7) instead of any two potentials of the interacting fields, i.e.
$$A_\mu ^a(k)^{}𝒟_{\mu \mu ^{}}(k)J^{T(2)a\mu ^{}}(A^{(0)},A^{(0)}).$$
We replace the remaining potentials by free ones. It is necessary to look at all possible substitutions in both first and second ways.
The number of terms appearing can be cut if we note, that it is need to keep only such terms in intermediate expressions which contain the propagators $`{}_{}{}^{}𝒟_{\mu \mu ^{}}^{}(k)`$ and $`{}_{}{}^{}𝒟_{\mu \mu ^{}}^{}(k^{})`$. These propagators give the terms, proportional to
$$\delta (\mathrm{Re}\epsilon ^l(k))=\left(\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)_{\omega =\omega _𝐤^l}^1[\delta (\omega \omega _𝐤^l)+\delta (\omega +\omega _𝐤^l)],$$
(6.6)
i.e. the factor taken into account the existence of plasmons with wavevector $`𝐤`$ and energy $`\omega _𝐤^l`$, in spite of the fact that the number density of the plasmon $`N_𝐤^l`$ is explicitly absent.
Hence it follows that for the first way only the replacement in (6.5)
$$A_\mu ^a(k^{})^{}𝒟_{\mu \mu ^{}}^{}(k^{})\stackrel{~}{J}^{T(3)a\mu ^{}}(A^{(0)},A^{(0)},A^{(0)}),A_\nu ^c(k_3)A_\nu ^{(0)c}(k_3),\mathrm{}$$
gives desired contribution (similar for the conjugate term). This leads (6.5) to the expression
$$idk^{}\{f^{bcf}f^{fde}f^{ac^{}g}f^{gd^{}e^{}}{}_{}{}^{}𝒟_{\mu \mu ^{}}^{}(k^{})\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }\stackrel{~}{\mathrm{\Sigma }}_{k^{},k_1^{},k_2^{},k_3^{}}^{\mu ^{}\nu ^{}\lambda ^{}\sigma ^{}}(A_\nu ^{}^{(0)c^{}}(k_3^{})A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})$$
$$\times A_\nu ^{(0)c}(k_3)A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)A_\nu ^{}^{(0)c^{}}(k_3^{})A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})A_\nu ^{(0)c}(k_3)A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)$$
$$A_\nu ^{}^{(0)c^{}}(k_3^{})A_\nu ^{(0)c}(k_3)A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})$$
(6.7)
$$+A_\nu ^{}^{(0)c^{}}(k_3^{})A_\nu ^{(0)c}(k_3)A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2))$$
$$\times \delta (kk_1k_2k_3)\delta (k^{}k_1^{}k_2^{}k_3^{})\underset{i=1}{\overset{3}{}}dk_idk_i^{}(ab,kk^{},compl.conj.)\}.$$
In the second case, at first step it should be replaced by
$$A_\mu ^a(k^{})^{}𝒟_{\mu \mu ^{}}^{}(k^{})J^{T(2)a\mu ^{}}(A^{(0)},A^{(0)}).$$
This gives
$$idk^{}\{f^{bcf}f^{fde}f^{ab^{}c^{}}{}_{}{}^{}𝒟_{\mu \mu ^{}}^{}(k^{})\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }S_{k^{},k_1^{},k_2^{}}^{\mu ^{}\nu ^{}\lambda ^{}}(A_\nu ^{}^{(0)b^{}}(k_1^{})A_\lambda ^{}^{(0)c^{}}(k_2^{})A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)$$
$$A_\nu ^{}^{(0)b^{}}(k_1^{})A_\lambda ^{}^{(0)c^{}}(k_2^{})A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)A_\nu ^{}^{(0)b^{}}(k_1^{})A_\lambda ^{}^{(0)c^{}}(k_2^{})A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2)$$
$$+A_\nu ^{}^{(0)b^{}}(k_1^{})A_\lambda ^{}^{(0)c^{}}(k_2^{})A_\nu ^c(k_3)A_\lambda ^d(k_1)A_\sigma ^e(k_2))$$
(6.8)
$$\times \delta (k^{}k_1^{}k_2^{})\delta (kk_1k_2k_3)dk_1dk_2dk_3dk_1^{}dk_2^{}+(kk^{},ab,compl.conj.)\}.$$
By virtue of the stochasticity of gauge fields, the last term inside the parentheses vanishes. We replace the remaining potentials of the interacting fields by the rules
$$A_\nu ^c(k_3)^{}𝒟_{\nu \nu ^{}}(k_3)J^{T(2)c\nu ^{}}(A^{(0)},A^{(0)}),A_\lambda ^d(k_1)A_\lambda ^{(0)d}(k_1),A_\sigma ^e(k_2)A_\sigma ^{(0)e}(k_2),etc..$$
By inspecting the expression (6.8), one sees that, complex-conjugate and non-conjugate amplitudes of free fields enter by a non-symmetric fashion. As will be shown below, in this case $`\delta `$-functions (6.4) type expressing the energy and momentum conservation laws in the plasmon-plasmon scattering have not arisen.
By virtue of the fact that we restrict our consideration to the derivation of the kinetic equation for colorless excitations, i.e.
$$A^{(0)d}A^{(0)e}\delta ^{de},A^{(0)d^{}}A^{(0)e^{}}\delta ^{d^{}e^{}},$$
(6.9)
all terms inside the parentheses of the expression (6.7), excepting the first term (the sixth-order correlator) vanish because $`\delta ^{dc},\delta ^{d^{}c^{}},\mathrm{},`$ are contracted with antisymmetric structure constants. Further, it is necessary to decouple the averaging of six potentials into pairs. We define below, what the correlations decoupling is responsible for in the four-plasmon decay processes. By using the definition of the correlation function (3.24), (3.25), we have
$$A_\nu ^{}^{(0)c^{}}(k_3^{})A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})A_\nu ^{(0)c}(k_3)A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)$$
(6.10)
$$=A_\nu ^{}^{(0)c^{}}(k_3^{})A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})A_\lambda ^{(0)d}(k_1)A_\nu ^{(0)c}(k_3)A_\sigma ^{(0)e}(k_2)+\mathrm{}$$
$$=I_{\nu ^{}\lambda ^{}}(k_3^{})\delta ^{c^{}d^{}}\delta (k_3^{}+k_1^{})I_{\sigma ^{}\lambda }(k_2^{})\delta ^{e^{}d}\delta (k_2^{}k_1)I_{\nu \sigma }(k_3)\delta ^{ce}\delta (k_3+k_2)+\mathrm{}.$$
After substitution of the first term on the r.h.s. (6.10) into (6.7) and performing integration over $`dk_1dk_3dk_2^{}dk_3^{}`$ and elementary color algebra, we obtain
$$\delta ^{ab}N_c^2(i)\{^{}𝒟_{\mu \mu ^{}}^{}(k)\mathrm{\Sigma }_{k,k,k_2,k_1}^{\mu \nu \lambda \sigma }\stackrel{~}{\mathrm{\Sigma }}_{k,k_1^{},k,k_1^{}}^{\mu ^{}\nu ^{}\lambda ^{}\sigma ^{}}I_{\nu ^{}\lambda ^{}}(k_1^{})I_{\sigma ^{}\lambda }(k)I_{\nu \sigma }(k_2)dk_1^{}dk_2(compl.conj.)\}.$$
As can be seen from the last expression, the required $`\delta `$-functions (6.4) are not appeared and therefore this expression is not associated with the plasmon-plasmon scattering and it should be dropped. This is a general rule. The decomposition of averaging of free field amplitudes into the correlators containing the pair of complex-conjugate potentials or one of the nonconjugate potentials between the inside of the angular brackets (statistical averaging), does not give a contribution to the process of interest to us. For this reason, it is necessary to fully drop all contribution in the process of the plasmon-plasmon scattering defined by the expression (6.8), since making the correlation decoupling, the pair with complex-conjugate amplitudes or without conjugate necessarily arises. We write out decoupling of the sixth-order correlator, which gives a contribution to equation (6.1). Suppressing color and Lorentz indices and employing a condensed notion, $`A_1A_\mu ^a(k_1)`$, we have
$$A_3^{}^{}A_2^{}^{}A_1^{}^{}A_3A_2A_1=3\{A_3^{}^{}A_3A_1^{}^{}A_1A_2^{}^{}A_2+A_3^{}^{}A_3A_1^{}^{}A_2A_2^{}^{}A_1$$
$$+A_3^{}^{}A_1A_1^{}^{}A_3A_3^{}^{}A_2+A_3^{}^{}A_1A_1^{}^{}A_2A_2^{}^{}A_3+A_3^{}^{}A_2A_1^{}^{}A_3A_2^{}^{}A_1$$
$$+A_3^{}^{}A_2A_1^{}^{}A_1A_2^{}^{}A_3\}.$$
(6.11)
Now we consider the terms with $`S`$-functions on the r.h.s. of equation (3.26) and here, we also write them separately
$$idk^{}\{f^{bcd}S_{k,k_1,k_2}^{\mu \nu \lambda }A_\mu ^a(k^{})A_\nu ^c(k_1)A_\lambda ^d(k_2)\delta (kk_1k_2)dk_1dk_2$$
$$f^{acd}S_{k^{},k_1,k_2}^{\mu \nu \lambda }A_\mu ^b(k)A_\nu ^c(k_1)A_\lambda ^d(k_2)\delta (k^{}k_1k_2)dk_1dk_2\}.$$
(6.12)
According to the previous discussion, at a first step it is necessary to perform the replacement
$$A_\mu ^a(k^{})^{}𝒟_{\mu \mu ^{}}^{}(k^{})\stackrel{~}{J}^{T(3)a\mu ^{}}(A^{(0)},A^{(0)},A^{(0)}),$$
$$A_\mu ^b(k)^{}𝒟_{\mu \mu ^{}}(k)\stackrel{~}{J}^{T(3)b\mu ^{}}(A^{(0)},A^{(0)},A^{(0)}).$$
Furthermore we consequently replace the remaining two potentials of the interacting fields in the correlators (6.12) by $`A_\nu ^c(k_1)^{}𝒟_{\nu \nu ^{}}(k_1)J^{T(2)c\nu ^{}}(A^{(0)},A^{(0)}),A_\lambda ^d(k_2)A_\lambda ^{(0)d}(k_2)`$, etc. This automatically leads to symmetry of contribution in $`A^{(0)}`$ and $`A^{(0)}`$. As result we obtain
$$\frac{i}{2}\{^{}𝒟_{\mu \mu ^{}}^{}(k^{})f^{bcf}f^{fde}f^{ac^{}g}f^{gd^{}e^{}}(S_{k,k_3,k_1+k_2}^{\mu \nu \rho }S_{k,k_1+k_2,k_3}^{\mu \rho \nu })\stackrel{~}{\mathrm{\Sigma }}_{k^{},k_1^{},k_2^{},k_3^{}}^{\mu ^{}\nu ^{}\lambda ^{}\sigma ^{}}{}_{}{}^{}𝒟_{\rho \rho ^{}}^{}(k_1+k_2)$$
$$\times (S_{k_1+k_2,k_1,k_2}^{\rho ^{}\lambda \sigma }S_{k_1+k_2,k_2,k_1}^{\rho ^{}\sigma \lambda })(A_\nu ^{}^{(0)c^{}}(k_3^{})A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)A_\nu ^{(0)c}(k_3)$$
$$A_\nu ^{}^{(0)c^{}}(k_3^{})A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})A_\nu ^{(0)c}(k_3)A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)$$
(6.13)
$$A_\nu ^{}^{(0)c^{}}(k_3^{})A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)A_\nu ^{(0)c}(k_3)A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{}))$$
$$\times \delta (kk_1k_2k_3)\delta (k^{}k_1^{}k_2^{}k_3^{})\underset{i=1}{\overset{3}{}}dk_idk_i^{}(ab,kk^{},compl.conj.)\}.$$
Because of (6.9), in the last expression it is necessary to retain only the sixth-order correlator. Finally adding (6.7) and (6.13), we obtain the following equation, instead of (3.26)
$$\frac{}{k_\lambda }[k^2g^{\mu \nu }(1+\xi ^1)k^\mu k^\nu \mathrm{\Pi }^{H\mu \nu }(k)]\frac{I_{\mu \nu }^{ab}}{x^\lambda }=$$
$$=idk^{}\{f^{bcf}f^{fde}f^{ac^{}g}f^{gd^{}e^{}}{}_{}{}^{}𝒟_{\mu \mu ^{}}^{}(k^{})\stackrel{~}{\mathrm{\Sigma }}_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }\stackrel{~}{\mathrm{\Sigma }}_{k^{},k_1^{},k_2^{},k_3^{}}^{\mu ^{}\nu ^{}\lambda ^{}\sigma ^{}}$$
(6.14)
$$\times (A_\nu ^{}^{(0)c^{}}(k_3^{})A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{}))(A_\nu ^{(0)c}(k_3)A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2))$$
$$\times \delta (kk_1k_2k_3)\delta (k^{}k_1^{}k_2^{}k_3^{})\underset{i=1}{\overset{3}{}}dk_idk_i^{}(ab,kk^{},compl.conj.)\}.$$
## 7 HTL-amplitudes
Let us transform equation (6.14) into a suitable form for our further research. For this purpose we perform the symmetrization of the coefficient functions in the integrand on the r.h.s. of equation (6.14) over possible permutations of color and Lorentz indices, and arguments of potentials of gauge fields within of the two groups $`(A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)A_\nu ^{(0)c}(k_3))`$ and $`(A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})A_\nu ^{}^{(0)c^{}}(k_3^{}))`$ inside the statistical averaging angular brackets.
For example, for the first group this symmetrization leads to the expression
$$f^{bcf}f^{fde}\stackrel{~}{\mathrm{\Sigma }}_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)A_\nu ^{(0)c}(k_3)=\frac{1}{3!}\{f^{bcf}f^{fde}(\stackrel{~}{\mathrm{\Sigma }}_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }\stackrel{~}{\mathrm{\Sigma }}_{k,k_2,k_1,k_3}^{\mu \nu \sigma \lambda })$$
$$+f^{bdf}f^{fce}(\stackrel{~}{\mathrm{\Sigma }}_{k,k_3,k_2,k_1}^{\mu \lambda \nu \sigma }\stackrel{~}{\mathrm{\Sigma }}_{k,k_2,k_3,k_1}^{\mu \lambda \sigma \nu })+f^{bef}f^{fdc}(\stackrel{~}{\mathrm{\Sigma }}_{k,k_1,k_3,k_2}^{\mu \sigma \lambda \nu }\stackrel{~}{\mathrm{\Sigma }}_{k,k_3,k_1,k_2}^{\mu \sigma \nu \lambda })\}$$
$$\times A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)A_\nu ^{(0)c}(k_3)$$
$$=\frac{1}{3!}\{[(\stackrel{~}{\mathrm{\Sigma }}_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }\stackrel{~}{\mathrm{\Sigma }}_{k,k_2,k_1,k_3}^{\mu \nu \sigma \lambda })+(\stackrel{~}{\mathrm{\Sigma }}_{k,k_1,k_3,k_2}^{\mu \sigma \lambda \nu }\stackrel{~}{\mathrm{\Sigma }}_{k,k_3,k_1,k_2}^{\mu \sigma \nu \lambda })]f^{bcf}f^{fde}$$
(7.1)
$$+[(\stackrel{~}{\mathrm{\Sigma }}_{k,k_3,k_2,k_1}^{\mu \lambda \nu \sigma }\stackrel{~}{\mathrm{\Sigma }}_{k,k_2,k_3,k_1}^{\mu \lambda \sigma \nu })(\stackrel{~}{\mathrm{\Sigma }}_{k,k_1,k_3,k_2}^{\mu \sigma \lambda \nu }\stackrel{~}{\mathrm{\Sigma }}_{k,k_3,k_1,k_2}^{\mu \sigma \nu \lambda })]f^{bdf}f^{fce}\}A^{(0)d}_\lambda (k_1)A^{(0)e}_\sigma (k_2)A^{(0)c}_\nu (k_3).$$
In writing the last equality, we have used the relation between the structure constants
$$f^{bef}f^{fdc}=f^{bcf}f^{fde}f^{bdf}f^{fce}.$$
In a similar way we transform the coefficient $`f^{ac^{}g}f^{gd^{}e^{}}\stackrel{~}{\mathrm{\Sigma }}_{k^{},k_1^{},k_2^{},k_3^{}}^{\mu ^{}\nu ^{}\lambda ^{}\sigma ^{}}`$ in front of the second group of potentials of gauge fields.
The expression (7.1) is convenient because it enables us to rewrite the functions preceding the correlators on the r.h.s. of the kinetic equation (6.14) in terms of HTL-amplitudes . Actually, by the definition of the $`\stackrel{~}{\mathrm{\Sigma }}`$-function (5.9) the coefficient in front of $`f^{bcf}f^{fde}`$ equals
$$\mathrm{\Sigma }_{k,k_1,k_2,k_3}^{\mu \nu \lambda \sigma }\mathrm{\Sigma }_{k,k_2,k_1,k_3}^{\mu \nu \sigma \lambda }^{}𝒟_{\rho \alpha }(k_1+k_2)(S_{k,k_3,k_1+k_2}^{\mu \nu \rho }S_{k,k_1+k_2,k_3}^{\mu \rho \nu })(S_{k_1+k_2,k_1,k_2}^{\alpha \lambda \sigma }S_{k_1+k_2,k_2,k_1}^{\alpha \sigma \lambda })$$
$$+(\nu \sigma ,k_2k_3).$$
(7.2)
Furthermore, we use initial definitions (3.23), (3.19) and (3.21). We present the integration measure $`\mathrm{d}^4p`$ as $`\mathrm{d}p^0|𝐩|^2\mathrm{d}|𝐩|\mathrm{d}\mathrm{\Omega }`$, where $`\mathrm{d}\mathrm{\Omega }`$ is the angular measure. Using the definition of the equilibrium distributions (3.8) (for $`\mu =0`$) and taking into account
$$N_c\underset{\mathrm{}}{\overset{+\mathrm{}}{}}|𝐩|^2\mathrm{d}|𝐩|\underset{\mathrm{}}{\overset{+\mathrm{}}{}}p_0dp_0\frac{df^0(p_0)}{dp_0}=\frac{3}{4\pi }\left(\frac{\omega _{pl}}{g}\right)^2,$$
we perform the integral over $`\mathrm{d}p_0`$ and the radial integral over $`\mathrm{d}|𝐩|`$ in the expressions for $`S^{(II)}`$-function (3.19) and $`\mathrm{\Sigma }^{(II)}`$-function (3.21).
This enables us to present the expression (7.2) in the following form
$$g^2\{^{}\mathrm{\Gamma }^{\mu \sigma \lambda \nu }(k,k_2,k_1,k_3)^{}𝒟_{\rho \alpha }(k_1+k_2)^{}\mathrm{\Gamma }^{\mu \nu \rho }(k,k_3,k_1k_2)^{}\mathrm{\Gamma }^{\alpha \lambda \sigma }(k_1+k_2,k_1,k_2)$$
$$^{}𝒟_{\rho \alpha }(k_1+k_3)^{}\mathrm{\Gamma }^{\mu \sigma \rho }(k,k_2,k_1k_3)^{}\mathrm{\Gamma }^{\alpha \lambda \nu }(k_1+k_3,k_1,k_3)\}$$
$$g^2{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \sigma \lambda \nu }(k,k_2,k_1,k_3),$$
(7.3)
where
$${}_{}{}^{}\mathrm{\Gamma }_{}^{\mu \nu \lambda \sigma }(k,k_1,k_2,k_3)\mathrm{\Gamma }^{\mu \lambda \nu \sigma }+\delta \mathrm{\Gamma }^{\mu \nu \lambda \sigma }(k,k_3,k_1,k_2)$$
(7.4)
is the effective four-gluon vertex, which represents a sum of a bare four-gluon vertex
$$\mathrm{\Gamma }^{\mu \nu \lambda \sigma }=2g^{\mu \nu }g^{\lambda \sigma }g^{\mu \lambda }g^{\nu \sigma }g^{\mu \sigma }g^{\nu \lambda },$$
and a corresponding HTL-correction
$$\delta \mathrm{\Gamma }^{\mu \nu \lambda \sigma }(k,k_1,k_2,k_3)=m_D^2\frac{\mathrm{d}\mathrm{\Omega }}{4\pi }\frac{v^\mu v^\nu v^\lambda v^\sigma }{vk+iϵ}[\frac{1}{v(k+k_1)+iϵ}(\frac{\omega _2}{vk_2iϵ}$$
$$\frac{\omega _3}{vk_3iϵ})\frac{1}{v(k+k_3)+iϵ}(\frac{\omega _1}{vk_1iϵ}\frac{\omega _2}{vk_2iϵ})],(v^\mu =(1,𝐯),v^2=0);$$
$${}_{}{}^{}\mathrm{\Gamma }_{}^{\mu \nu \rho }(k,k_1,k_2)\mathrm{\Gamma }^{\mu \nu \rho }(k,k_1,k_2)+\delta \mathrm{\Gamma }^{\mu \nu \rho }(k,k_1,k_2)$$
(7.5)
is the effective three-gluon vertex. It also represents a sum of the bare three-gluon vertex
$$\mathrm{\Gamma }^{\mu \nu \rho }(k,k_1,k_2)=g^{\mu \nu }(kk_1)^\rho +g^{\nu \rho }(k_1k_2)^\mu +g^{\mu \rho }(k_2k)^\nu $$
(7.6)
and corresponding HTL-correction
$$\delta \mathrm{\Gamma }^{\mu \nu \rho }(k,k_1,k_2)=m_D^2\frac{\mathrm{d}\mathrm{\Omega }}{4\pi }\frac{v^\mu v^\nu v^\rho }{vk+iϵ}\left(\frac{\omega _2}{vk_2iϵ}\frac{\omega _1}{vk_1iϵ}\right),$$
(7.7)
$`m_D^2=3\omega _{pl}^2`$ is the Debye screening mass.
The polarization tensor (3.17) in this notation takes the form
$$\mathrm{\Pi }^{\mu \nu }(k)=m_D^2\left(g^{\mu 0}g^{\nu 0}\omega \frac{\mathrm{d}\mathrm{\Omega }}{4\pi }\frac{v^\mu v^\nu }{vk+iϵ}\right).$$
In a similar way, the coefficient in front of the product of structure constants $`f^{bdf}f^{fce}`$ in (7.1) may be presented as
$$g^2\{^{}\mathrm{\Gamma }^{\mu \lambda \nu \sigma }(k,k_1,k_3,k_2)^{}𝒟_{\rho \alpha }(k_2+k_3)^{}\mathrm{\Gamma }^{\mu \lambda \rho }(k,k_1,k_2k_3)^{}\mathrm{\Gamma }^{\alpha \nu \sigma }(k_2+k_3,k_3,k_2)$$
$$+^{}𝒟_{\rho \alpha }(k_1+k_3)^{}\mathrm{\Gamma }^{\mu \sigma \rho }(k,k_2,k_1k_3)^{}\mathrm{\Gamma }^{\alpha \lambda \nu }(k_1+k_3,k_1,k_3)\}$$
$$g^2{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \lambda \nu \sigma }(k,k_1,k_3,k_2).$$
(7.8)
Performing transformation of the coefficient preceding $`(A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})A_\nu ^{}^{(0)c^{}}(k_3^{}))`$ in a similar manner, we can cast the equation (6.14) in the following form
$$\frac{}{k_\lambda }[k^2g^{\mu \nu }(1+\xi ^1)k^\mu k^\nu \mathrm{\Pi }^{H\mu \nu }(k)]\frac{I_{\mu \nu }^{ab}}{x^\lambda }$$
(7.9)
$$=i\frac{g^4}{(3!)^2}dk^{}[\{f^{bcf}f^{fde}{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \sigma \lambda \nu }(k,k_2,k_1,k_3)+f^{bdf}f^{fce}{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \lambda \nu \sigma }(k,k_1,k_3,k_2)\}$$
$$\times ^{}𝒟_{\mu \mu ^{}}^{}(k^{})\{f^{ac^{}g}f^{gd^{}e^{}}{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu ^{}\sigma ^{}\lambda ^{}\nu ^{}}(k^{},k_2^{},k_1^{},k_3^{})+f^{ad^{}g}f^{gc^{}e^{}}{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu ^{}\lambda ^{}\nu ^{}\sigma ^{}}(k,k_1^{},k_3^{},k_2^{})\}$$
$$\times A_\nu ^{}^{(0)c^{}}(k_3^{})A_\lambda ^{}^{(0)d^{}}(k_1^{})A_\sigma ^{}^{(0)e^{}}(k_2^{})A_\nu ^{(0)c}(k_3)A_\lambda ^{(0)d}(k_1)A_\sigma ^{(0)e}(k_2)$$
$$\times \delta (kk_1k_2k_3)\delta (k^{}k_1^{}k_2^{}k_3^{})\underset{i=1}{\overset{3}{}}dk_idk_i^{}(ab,kk^{},compl.conj.)].$$
The r.h.s. of equation (7.9) has a non-trivial color structure that actually is well represented by non-trivial color structure of the initial dynamical equation (3.4). As will be shown bellow, this leads to a qualitative distinction between the elastic scattering probability $`w(𝐤,𝐤_1;𝐤_2,𝐤_3)`$ of colorless plasmons in a hot QCD plasma and a similar one of plasmons in a hot QED plasma .
At the end of this Section we present the identities analogous to the effective Ward one in hot gauge theories . It can be shown that the following equalities hold
$$k_\mu ^{}\mathrm{\Gamma }^{\mu \nu \lambda \sigma }(k,k_1,k_2,k_3)=^{}\mathrm{\Gamma }^{\nu \lambda \sigma }(k_1,k_2,k+k_3)^{}\mathrm{\Gamma }^{\nu \lambda \sigma }(k+k_1,k_2,k_3),$$
(7.10)
$$k_{1\nu }^{}\mathrm{\Gamma }^{\mu \nu \lambda \sigma }(k,k_1,k_2,k_3)=^{}\mathrm{\Gamma }^{\mu \lambda \sigma }(k+k_1,k_2,k_3)^{}\mathrm{\Gamma }^{\mu \lambda \sigma }(k,k_1+k_2,k_3)$$
(similar contractions with $`k_{2\lambda },k_{3\sigma }`$),
$$k_\mu ^{}\mathrm{\Gamma }^{\mu \nu \rho }(k,k_1,k_2)=^{}𝒟^{1\nu \rho }(k_1)^{}𝒟^{1\nu \rho }(k_2)$$
(7.11)
(similar contractions with $`k_{1\nu },k_{2\rho }`$). Here, $`{}_{}{}^{}𝒟_{}^{1\mu \nu }(k)=P^{\mu \nu }(k)^{}\mathrm{\Delta }^{1t}(k)+Q^{\mu \nu }(k)^{}\mathrm{\Delta }^{1l}(k)`$ is the inverse propagator without the gauge fixing term.
## 8 The kinetic equation for plasmons
We are now in a position to explicitly compute the probability of plasmon-plasmon scattering. First of all, we make the correlation decoupling of the sixth-order correlators on the r.h.s. of equation (7.9) in terms of the pair ones by the rule (6.11). After cumbersome calculations, the r.h.s. of equation (7.9) can be written as
$$g^4\mathrm{Im}(^{}𝒟_{\mu \mu ^{}}(k))\{f^{bcf}f^{fde}{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \sigma \lambda \nu }(k,k_2,k_1,k_3)+f^{bdf}f^{fce}{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \lambda \nu \sigma }(k,k_1,k_3,k_2)\}$$
$$\times \{f^{acg}f^{gde}{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu ^{}\sigma ^{}\lambda ^{}\nu ^{}}(k,k_2,k_1,k_3)+f^{adg}f^{gce}{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu ^{}\lambda ^{}\nu ^{}\sigma ^{}}(k,k_1,k_3,k_2)\}$$
(8.1)
$$\times I_{\nu \nu ^{}}(k_3)I_{\lambda \lambda ^{}}(k_1)I_{\sigma \sigma ^{}}(k_2)\delta (kk_1k_2k_3)dk_1dk_2dk_3.$$
In deriving of (8.1) we have used two relations, which satisfy $`{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \nu \lambda \sigma }(k,k_1,k_2,k_3)`$
$${}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \nu \lambda \sigma }(k,k_1,k_2,k_3)+^{}\stackrel{~}{\mathrm{\Gamma }}^{\mu \lambda \nu \sigma }(k,k_2,k_1,k_3)+^{}\stackrel{~}{\mathrm{\Gamma }}^{\mu \nu \sigma \lambda }(k,k_1,k_3,k_2)=0,$$
$${}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \nu \lambda \sigma }(k,k_1,k_2,k_3)=^{}\stackrel{~}{\mathrm{\Gamma }}^{\mu \sigma \lambda \nu }(k,k_3,k_2,k_1).$$
(8.2)
Their correctness may be verified by a direct calculation using known properties of HTL-amplitudes , entering into the definition of $`{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}_{}^{\mu \nu \lambda \sigma }(k,k_1,k_2,k_3)`$ (7.3). The first of the relations in (8.2) supplements, in the some sense, a similar relation for the HTL-correction $`\delta \mathrm{\Gamma }_4`$ to the bare four-gluon vertex, originally proposed by Frenkel and Taylor . The second one is a property of invariance of the function $`{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}`$, when the momenta order is reversed. Notice that order of the space-time indices in equation (8.2) is important.
Furthermore, if we restrict our consideration to the study of plasmons scattering by plasmons, then in the spectral decomposition of $`I_{\mu \mu ^{}}`$ into the orthonormal projectors $`P_{\mu \mu ^{}}`$ and $`Q_{\mu \mu ^{}}`$, it is necessary to keep only the longitudinal part, i.e. set
$$I_{\nu \nu ^{}}(k_3)=Q_{\nu \nu ^{}}(k_3)I^l(k_3),I_{\lambda \lambda ^{}}(k_1)=Q_{\lambda \lambda ^{}}(k_1)I^l(k_1),I_{\sigma \sigma ^{}}(k_2)=Q_{\sigma \sigma ^{}}(k_2)I^l(k_2).$$
(8.3)
In the propagator $`{}_{}{}^{}𝒟_{\mu \mu ^{}}^{}(k)`$ we also keep only the longitudinal part $`Q_{\mu \mu ^{}}(k)^{}\mathrm{\Delta }^l(k)`$. To determine the imaginary part of $`{}_{}{}^{}\mathrm{\Delta }_{}^{l}(k)=1/(k^2\mathrm{\Pi }^l(k))`$ we use the approximation (see e.g., Pustovalov and Silin in )
$$\frac{1}{k^2\mathrm{\Pi }^l(k)}\frac{\mathrm{P}}{\mathrm{Re}(k^2\mathrm{\Pi }^l(k))}i\pi sign(\mathrm{Im}(k^2\mathrm{\Pi }^l(k))\delta (\mathrm{Re}(k^2\mathrm{\Pi }^l(k)))=$$
$$=\frac{\mathrm{P}}{k^2\mathrm{Re}\epsilon ^l(k)}\frac{i\pi }{k^2}sign\omega \delta (\mathrm{Re}\epsilon ^l(k)),$$
(8.4)
which holds for a small $`\mathrm{Im}\epsilon ^l(k)`$. Here, we consider that because of (2.3), $`sign(\mathrm{Im}\epsilon ^l(k))=sign\omega `$. The symbol $`\mathrm{P}`$ denotes a principal value. We perceive the $`\delta `$-function of the real part of the longitudinal permeability, which appears in (8.4) in the ordinary sense (6.6). Substituting (6.6) into (8.4) and choosing $`\omega =\omega _𝐤^l>0`$ for definiteness, we obtain the required relation
$$\mathrm{Im}^{}𝒟_{\mu \mu ^{}}(k)\frac{\pi }{k^2}Q_{\mu \mu ^{}}(k)\left(\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)_{\omega =\omega _𝐤^l}^1\delta (\omega \omega _𝐤^l).$$
(8.5)
With allowing (8.3), (8.5) and performing the color algebra with the following identities
$$f^{bcf}f^{fde}f^{acg}f^{gde}=N_c^2\delta ^{ab},$$
$$f^{bcf}f^{fde}f^{adg}f^{gce}=\frac{1}{2}N_c^2\delta ^{ab},$$
(8.6)
we rewrite (8.1) in the form
$$\delta ^{ab}\pi g^4N_c^2\frac{1}{k^2}\left(\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)_{\omega =\omega _𝐤^l}^1\delta (\omega \omega _𝐤^l)\{|^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)|^2+|^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_1,k_3,k_2)|^2$$
$$+\mathrm{Re}(^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)^{}\stackrel{~}{\mathrm{\Gamma }}^{}(k,k_1,k_3,k_2))\}\frac{1}{\overline{u}^2(k)\overline{u}^2(k_1)\overline{u}^2(k_2)\overline{u}^2(k_3)}$$
(8.7)
$$\times I^l(k_1)I^l(k_2)I^l(k_3)\delta (kk_1k_2k_3)dk_1dk_2dk_3.$$
Here, we denote
$${}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_1,k_2,k_3)^{}\stackrel{~}{\mathrm{\Gamma }}^{\mu \nu \lambda \sigma }(k,k_1,k_2,k_3)\overline{u}_\mu (k)\overline{u}_\nu (k_1)\overline{u}_\lambda (k_2)\overline{u}_\sigma (k_3).$$
(8.8)
By virtue of the fact, that the spectral intensity functions $`I^l(k_i),i=1,2,3`$, entering into (8.7), satisfy the equations $`\mathrm{Re}\epsilon ^l(k_i)I^l(k_i)=0`$ , they have the following structure
$$I^l(k_i)=I_{𝐤_i}^l\delta (\omega _i\omega _{𝐤_i}^l)+I_{𝐤_i}^l\delta (\omega _i+\omega _{𝐤_i}^l),i=1,2,3,$$
(8.9)
where $`I_{𝐤_i}^l`$ are certain functions of the wavevectors $`𝐤_i`$.
We substitute (8.9) into (8.7) and perform the integration over the frequency $`\mathrm{d}\omega _i`$ with the help of the $`\delta `$-functions. The function $`I_k^l`$ in the l.h.s. of the kinetic equation (7.9) also has a structure of form (8.9). Here, we keep only a positive branch $`\omega >0`$ in agreement with (8.5). Furthermore performing integration (7.9) over $`\mathrm{d}\omega `$, where the r.h.s. of (7.9) has the form (8.7), as result we obtain
$$k^2\left(\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)_{\omega =\omega _𝐤^l}\frac{I_𝐤^l}{t}+(k^2)\left(\frac{\mathrm{Re}\epsilon ^l(k)}{𝐤}\right)_{\omega =\omega _𝐤^l}\frac{I_𝐤^l}{𝐱}$$
(8.10)
$$=\pi g^4N_c^2\frac{1}{k^2}\left(\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)_{\omega =\omega _𝐤^l}^1[Q(k,k_1,k_2,k_3)\delta (k+k_1+k_2+k_3)$$
$$+Q(k,k_1,k_2,k_3)\delta (kk_1k_2k_3)+Q(k,k_1,k_2,k_3)\delta (k+k_1+k_2k_3)$$
$$+Q(k,k_1,k_2,k_3)\delta (k+k_1k_2+k_3)+Q(k,k_1,k_2,k_3)\delta (kk_1+k_2+k_3)$$
$$+Q(k,k_1,k_2,k_3)\delta (k+k_1k_2k_3)+Q(k,k_1,k_2,k_3)\delta (kk_1+k_2k_3)$$
$$+Q(k,k_1,k_2,k_3)\delta (kk_1k_2+k_3)]_{onshell}I_{𝐤_1}^lI_{𝐤_2}^lI_{𝐤_3}^l\mathrm{d}𝐤_1\mathrm{d}𝐤_2\mathrm{d}𝐤_3,$$
where we have defined
$$Q(k,k_1,k_2,k_3)\frac{1}{\overline{u}^2(k)\overline{u}^2(k_1)\overline{u}^2(k_2)\overline{u}^2(k_3)}$$
(8.11)
$$\times \{|^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)|^2+|^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_1,k_3,k_2)|^2+\mathrm{Re}(^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)^{}\stackrel{~}{\mathrm{\Gamma }}^{}(k,k_1,k_3,k_2))\}.$$
In equation (8.10) the $`x`$-dependence of $`I_𝐤^l,I_{𝐤_1}^l`$, etc is understood, although not written explicitly.
The first term inside the square brackets on the r.h.s. of (8.10) describes the process of simultaneous fusion or emission of four plasmons in the plasma. Considering that the $`\delta `$-function has no support on the plasmon mass shell, its contribution to the kinetic equation vanishes. The second, third and fourth terms describe the decay of one plasmon into three or the fusion of three plasmons into one. As was mentioned in section 2, these processes are forbidden by the conservation law (2.23). Therefore these terms also vanish. The remaining terms describe the scattering of two plasmons off two plasmons. Three terms imply the existence of three possible channels of a given process, which change the plasmon numbers density $`N_𝐤^l`$
$$\mathrm{g}^{}+\mathrm{g}_1^{}\mathrm{g}_2^{}+\mathrm{g}_3^{},\mathrm{g}^{}+\mathrm{g}_2^{}\mathrm{g}_1^{}+\mathrm{g}_3^{},\mathrm{g}^{}+\mathrm{g}_3^{}\mathrm{g}_1^{}+\mathrm{g}_2^{}.$$
If we perform replacements $`𝐤_1𝐤_2`$ in the next to last term on the r.h.s. of (8.10), and $`𝐤_1𝐤_3`$ in the last term, then the r.h.s. of (8.10) can be presented as
$$\pi g^4N_c^2\frac{1}{k^2}\left(\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)_{\omega =\omega _𝐤^l}^1[Q(k,k_1,k_2,k_3)+Q(k,k_2,k_1,k_3)+Q(k,k_3,k_2,k_1)]_{onshell}$$
$$\times \delta (\omega _𝐤^l+\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l)\delta (𝐤+𝐤_\mathrm{𝟏}𝐤_\mathrm{𝟐}𝐤_\mathrm{𝟑})I_{𝐤_1}^lI_{𝐤_2}^lI_{𝐤_3}^l\mathrm{d}𝐤_1\mathrm{d}𝐤_2\mathrm{d}𝐤_3.$$
By using the definition of kernel $`Q`$ (8.11) and the properties (8.2), it can be shown, that
$$Q(k,k_2,k_1,k_3)+Q(k,k_3,k_2,k_1)=2Q(k,k_1,k_2,k_3).$$
Taking into account the last relation, it goes over from the function $`I_𝐤^l`$ to the number density of plasmons
$$N_𝐤^l=\left(k^2\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)_{\omega =\omega _𝐤^l}I_𝐤^l$$
and recovering the complete form for finite values $`N_𝐤^l`$, we obtain the required Boltzmann equation (2.25), where the probability of plasmon-plasmon scattering is defined by the following expression
$$w(𝐤,𝐤_1;𝐤_2,𝐤_3)=3\pi (2\pi )^5g^4N_c^2[\frac{1}{k^2k_1^2k_2^2k_3^2}\left(\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)^1\left(\frac{\mathrm{Re}\epsilon ^l(k_1)}{\omega _1}\right)^1$$
(8.12)
$$\times \left(\frac{\mathrm{Re}\epsilon ^l(k_2)}{\omega _2}\right)^1\left(\frac{\mathrm{Re}\epsilon ^l(k_3)}{\omega _3}\right)^1\{|^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)|^2+|^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_1,k_3,k_2)|^2$$
$$+\mathrm{Re}(^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)^{}\stackrel{~}{\mathrm{\Gamma }}^{}(k,k_1,k_3,k_2))\}\frac{1}{\overline{u}^2(k)\overline{u}^2(k_1)\overline{u}^2(k_2)\overline{u}^2(k_3)}]_{onshell}.$$
The result (8.12) is rather unexpected. As we can see from the expression (8.12), this probability does not reduced to the squared module of one scalar function, as this occurs in the Abelian plasma. Here, the scattering probability is defined by the squared module of two independent scalar functions and their interference<sup>3</sup><sup>3</sup>3It is easy to see from the expression (8.12), that the function $`w(𝐤,𝐤_1;𝐤_2,𝐤_3)`$ is positively definite., in spite of the fact that in this paper we restrict our consideration to a study of the nonlinear interaction of only colorless excitations. This radically distinguishes the Boltzmann equation (2.25), describing the effects of the collisions among colorless soft excitations from the corresponding Boltzmann equation including the effects of the collisions among colorless hard excitations . In the last case, the Boltzmann equation, corrected to color factors, fully coincides with the corresponding one in the Abelian plasma. This is a point which we find difficult to interpret and therefore additional analysis of this problem is required (see, also discussion in conclusion).
The scattering probability can be written in the form, which is manifestly symmetric under permutations of the external momenta $`𝐤,𝐤_1,𝐤_2`$ and $`𝐤_3`$
$$w(𝐤,𝐤_1;𝐤_2,𝐤_3)=\pi (2\pi )^5g^4N_c^2[\frac{1}{k^2k_1^2k_2^2k_3^2}\left(\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)^1\left(\frac{\mathrm{Re}\epsilon ^l(k_1)}{\omega _1}\right)^1\left(\frac{\mathrm{Re}\epsilon ^l(k_2)}{\omega _2}\right)^1$$
$$\times \left(\frac{\mathrm{Re}\epsilon ^l(k_3)}{\omega _3}\right)^1\{|^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)|^2+|^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_1,k_3,k_2)|^2+|^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_1,k_2,k_3)|^2+$$
$$\mathrm{Re}(^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)^{}\stackrel{~}{\mathrm{\Gamma }}^{}(k,k_1,k_3,k_2))+\mathrm{Re}(^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)^{}\stackrel{~}{\mathrm{\Gamma }}^{}(k,k_1,k_2,k_3))$$
$$+\mathrm{Re}(^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_1,k_3,k_2)^{}\stackrel{~}{\mathrm{\Gamma }}^{}(k,k_1,k_2,k_3))\}\frac{1}{\overline{u}^2(k)\overline{u}^2(k_1)\overline{u}^2(k_2)\overline{u}^2(k_3)}]_{onshell}.$$
This expression is suitable for checking the symmetry conditions (2.26), which is imposed on the plasmon-plasmon scattering probability.
Recall that the function $`{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}`$, which appears in the expression for probability (8.12), is defined by expressions (8.8) and (7.3). As we have shown in , the expression (8.8) (exactly, its part, independent on a gauge-parameter) is gauge invariant at least in a class of covariant and temporal gauges. With regard to the reasoning in section 4 this automatically leads to gauge invariance of the kinetic equation (2.25). The problem of the dependence of the plasmon-plasmon scattering probability (8.12) on a gauge parameter coming from the gauge fixing term in the gluon propagator, is discussed below.
Let us separate the terms with a gauge parameter from the expression $`{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)`$. By using the Ward identities (7.10) and (7.11), we have
$$\xi \{D_{\rho \alpha }(k_1+k_2)\mathrm{\Delta }^0(k_1+k_2)^{}\mathrm{\Gamma }^{\mu \nu \rho }(k,k_3,k_1k_2)^{}\mathrm{\Gamma }^{\alpha \lambda \sigma }(k_1+k_2,k_1,k_2)$$
$$+D_{\rho \alpha }(k_1+k_3)\mathrm{\Delta }^0(k_1+k_3)^{}\mathrm{\Gamma }^{\mu \sigma \rho }(k,k_2,k_1k_3)^{}\mathrm{\Gamma }^{\alpha \lambda \nu }(k_1+k_3,k_1,k_3)\}$$
$$\times \overline{u}_\mu (k)\overline{u}_\lambda (k_1)\overline{u}_\sigma (k_2)\overline{u}_\nu (k_3)$$
(8.13)
$$=\xi \{(\mathrm{\Delta }^0(k_1+k_2))^2(^{}𝒟^{1\mu \sigma }(k)^{}𝒟^{1\mu \sigma }(k_3))(^{}𝒟^{1\lambda \nu }(k_2)^{}𝒟^{1\lambda \nu }(k_1))$$
$$+(\mathrm{\Delta }^0(k_1+k_3))^2(^{}𝒟^{1\mu \nu }(k)^{}𝒟^{1\mu \nu }(k_2))(^{}𝒟^{1\lambda \sigma }(k_3)^{}𝒟^{1\lambda \sigma }(k_1))\}$$
$$\times \overline{u}_\mu (k)\overline{u}_\lambda (k_1)\overline{u}_\sigma (k_2)\overline{u}_\nu (k_3).$$
It is easily shown that expression on the r.h.s. of (8.13) vanish either because $`{}_{}{}^{}𝒟_{}^{1\mu \nu }(k)`$ is transverse, or by the definition of the mass-shall condition, i.e.
$$k_\mu ^{}𝒟^{1\mu \nu }(k)=0,^{}𝒟^{\mathrm{1\hspace{0.17em}0}\nu }(k)|_{\omega =\omega _𝐤^l}=0.$$
(8.14)
If we carry out the following replacements on the r.h.s. of the expressions (8.8), (7.3)
$$\overline{u}_\mu (k)\stackrel{~}{u}_\mu (k)\frac{k^2}{(ku)}(k_\mu u_\mu (ku)),$$
and
$${}_{}{}^{}𝒟_{\rho \alpha }^{}(k)^{}\stackrel{~}{𝒟}_{\rho \alpha }(k)=^{}𝒟_{\rho \alpha }(k)\left(\frac{\sqrt{2k^2\overline{u}^2}}{k^2(ku)}C_{\rho \alpha }(k)+\frac{\overline{u}^2(k)}{k^2(ku)^2}D_{\rho \alpha }(k)\right)^{}\mathrm{\Delta }^l(k)$$
$$\xi D_{\rho \alpha }(k)\mathrm{\Delta }^0(k)\xi _0\frac{k^2}{(ku)^2}D_{\rho \alpha }(k),$$
where $`\xi _0`$ is a gauge parameter in the temporal gauge, then it can be proved that a similar statement holds in the temporal gauge also.
Thus, the gauge-dependent parts of $`w(𝐤,𝐤_1;𝐤_2,𝐤_3)`$ drop out, since they are multiplied by the mass-shell factor. These factors are proportional to $`(\omega \omega _𝐤^l)`$. However, in the quantum case Baier et al observed that a naive calculation of the gluon damping rate in a covariant gauge appears to violate this consideration. The mass-shell factor is multiplied by the integral involving a power infrared divergence which is cut-off exactly on the scale $`(\omega \omega _𝐤^l)g^2T`$. This problem was careful discussed in consideration of the nonlinear Landau damping , when in (7.3) and consequently (8.13) it is necessary to set
$$k_1=k,k_2=k_1,k_3=k_1.$$
(8.15)
We have shown, that the gauge-dependent part of the nonlinear Landau damping rate vanishes on a mass-shell at least for the plasmons with zero momentum. In the case of the process of the plasmon-plasmon scattering considered here, the problem on the dependence of $`w(𝐤,𝐤_1;𝐤_2,𝐤_3)`$ on a gauge parameter is more subtle, since instead of (8.15), we have condition: $`k+k_1=k_2+k_3`$ only. Research into this nontrivial question goes beyond our present goal.
## 9 Lifetimes of plasmons
To calculate the lifetimes of colorless plasmons we first linearize the Boltzmann equation (2.25) (here, we assume that off-equilibrium fluctuation is perturbative small), writing the number density of plasmons as
$$N_𝐤^l=N_{eq}^l(𝐤)+\delta N_𝐤^l$$
where $`N_{eq}^l(𝐤)=(e^{\omega _𝐤^l/T^{}}1)^1`$ is the Planck distribution function and $`T^{}`$ is a certain constant, which can be interpreted as a plasmon gas temperature in the statistical equilibrium state. Then we find
$$\frac{\delta N_𝐤^l}{t}+𝐕_𝐤^l\frac{\delta N_𝐤^l}{𝐱}=$$
(9.1)
$$=\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\frac{\mathrm{d}𝐤_2}{(2\pi )^3}\frac{\mathrm{d}𝐤_3}{(2\pi )^3}(2\pi )^4\delta (\omega _𝐤^l+\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l)\delta (𝐤+𝐤_1𝐤_2𝐤_3)w(𝐤,𝐤_1;𝐤_2,𝐤_3)$$
$$\times \{\delta N_𝐤^l[N_{eq}^l(𝐤_2)N_{eq}^l(𝐤_3)(N_{eq}^l(𝐤_1)+1)N_{eq}^l(𝐤_1)(N_{eq}^l(𝐤_2)+1)(N_{eq}^l(𝐤_3)+1)]+$$
$$\delta N_{𝐤_2}^l[N_{eq}^l(𝐤_3)(N_{eq}^l(𝐤)+1)(N_{eq}^l(𝐤_1)+1)N_{eq}^l(𝐤)N_{eq}^l(𝐤_1)(N_{eq}^l(𝐤_3)+1)]+(𝐤𝐤_1,𝐤_2𝐤_3)\}.$$
This equation can be further simplified if we use the following parametrization for off-equilibrium fluctuation of the occupation number $`\delta N_𝐤^l`$
$$\delta N_𝐤^l\frac{\mathrm{d}N_{eq}^l(𝐤)}{\mathrm{d}\omega _𝐤^l}𝒲_𝐤^l=(1/T^{})N_{eq}^l(𝐤)(N_{eq}^l(𝐤)+1)𝒲_𝐤^l.$$
(9.2)
Let us introduce the momentum transfers, which carries the soft gluon exchanged in the collision of two plasmons, setting
$$𝐤_2=𝐤𝐪,𝐤_3=𝐤_1+𝐪.$$
(9.3)
Performing the integration over $`\mathrm{d}𝐤_3`$, replacing the integration over $`\mathrm{d}𝐤_2`$ by one with respect to momentum transfer and taking into account (9.2) and (9.3), finally we derive a linearized Boltzmann equation for $`𝒲_𝐤^l`$ function
$$\frac{𝒲_𝐤^l}{t}+𝐕_𝐤^l\frac{𝒲_𝐤^l}{𝐱}=\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\frac{\mathrm{d}𝐪}{(2\pi )^3}\frac{N_{eq}^l(𝐤_1)(N_{eq}^l(𝐤𝐪)+1)(N_{eq}^l(𝐤_1+𝐪)+1)}{(N_{eq}^l(𝐤)+1)}$$
$$\times w(𝐤,𝐤_1;𝐪)(2\pi )\delta (\omega _𝐤^l\omega _{𝐤𝐪}^l+\omega _{𝐤_1}^l\omega _{𝐤_1+𝐪}^l)\{𝒲_𝐤^l𝒲_{𝐤𝐪}^l+𝒲_{𝐤_1}^l𝒲_{𝐤_1+𝐪}^l\}.$$
(9.4)
Here, we goes over from the function $`w(𝐤,𝐤_1;𝐤_2,𝐤_3)`$ (8.12) to a new function $`w(𝐤,𝐤_1;𝐪)`$
$$w(𝐤,𝐤_1;𝐤_2,𝐤_3)|_{𝐤_2=𝐤𝐪,𝐤_3=𝐤_1+𝐪}w(𝐤,𝐤_1;𝐪).$$
Based on the exact form of the r.h.s. of equation (9.4), we define the lifetimes of the plasmon of momentum $`𝐤`$ as follows
$$\frac{1}{\tau _{pl}(𝐤)}=\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\frac{\mathrm{d}𝐪}{(2\pi )^3}\frac{N_{eq}^l(𝐤_1)(N_{eq}^l(𝐤𝐪)+1)(N_{eq}^l(𝐤_1+𝐪)+1)}{(N_{eq}^l(𝐤)+1)}$$
$$\times w(𝐤,𝐤_1;𝐪)(2\pi )\delta (\omega _𝐤^l\omega _{𝐤𝐪}^l+\omega _{𝐤_1}^l\omega _{𝐤_1+𝐪}^l).$$
(9.5)
Here, the integrand has a more involved structure in comparison with a similar expression for the case of the lifetimes of the hard transverse gluon . The reason of this fact is that the momentum of the soft quasiparticles (plasmons) becomes of the same order as momentum transfers and under these conditions one must take into account the non-trivial character of the frequency dependence $`\omega `$ of the collective gluon modes on momentum $`𝐤`$ and vertex corrections.
It is not difficult to estimate the order of the expression (9.5) at the momentum scale $`gT`$. Considering the plasmon gas in thermal equilibrium with hard particles from the heat bath, i.e. $`T^{}T`$, and using the definition of $`w`$ function, we obtain
$$\frac{1}{\tau _{pl}}g^3N_c^2T.$$
(9.6)
However, obtaining a numerical factor of proportionality in (9.6) is a complicated problem even for limiting case of $`𝐤=0`$-mode. Here, we restrict our consideration to following general remark, which somewhat simplifies the matter.
In proving the gauge invariance of the nonlinear Landau damping rate, we have shown that the $`{}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}`$ function (8.8), (7.3) entering into the definition of the scattering probability (8.12) can be introduced in its simplest form
$${}_{}{}^{}\stackrel{~}{\mathrm{\Gamma }}(k,k_2,k_1,k_3)=$$
$$k^2k_1^2k_2^2k_3^2\{^{}\mathrm{\Gamma }^{0000}(k,k_2,k_1,k_3)^{}𝒟_{\rho \alpha }(kk_3)^{}\mathrm{\Gamma }^{00\rho }(k,k_3,k_1k_2)^{}\mathrm{\Gamma }^{\alpha 00}(k_1+k_2,k_1,k_2)$$
$$^{}𝒟_{\rho \alpha }(kk_2)^{}\mathrm{\Gamma }^{00\rho }(k,k_2,k_1k_3)^{}\mathrm{\Gamma }^{\alpha 00}(k_1+k_3,k_1,k_3)\}|_{onshell}.$$
However, even if the last expression is accounted for, after performing the integration over solid angles in $`\delta \mathrm{\Gamma }^{0000}`$ and $`\delta \mathrm{\Gamma }^{00\rho }`$ HTL-amplitudes (see, e.g., ), we obtain expressions which are very cumbersome. For this reason passage to the limit $`|𝐤|0`$ is non-trivial.
The estimate (9.6) shows, that at the momentum scale $`gT`$ a ”collision” damping rate of the soft longitudinal mode is suppressed by a power of $`g`$ relative to a value of the nonlinear Landau damping rate ($`g^2T`$) . Therefore at the soft momentum scale it can be neglected by the influence of the plasmon interactions among themselves on the relaxation dynamics of a plasma excitations.
As was mentioned in Sec. 2, the process of the nonlinear Landau damping leads to an effective pumping of energy across the spectrum towards small wavenumbers. By the effect of pumping, all plasmons will tend to concentrate near small $`|𝐤|=|𝐤_0|0`$. However, phase space, within which the plasmons are occupied, proportional to $`|𝐤_0|^3`$, will also be very small (a similar state, when a great many plasmons with small wavenumbers are mainly concentrated in a small phase volume, in the theory of strong turbulent electron-ion plasma this is sometimes called a plasmon condensate ). By virtue of this fact intensive collisions between plasmons arise and this scattering process is described by the Boltzmann equation (2.25). This leads to a ”throw out” of plasmons from region of small $`|𝐤|`$ and thus a suppression of the increase of the $`𝐤=0`$-mode. In mathematical language this denotes that for definite values of the momentum, the magnitude of the off-equilibrium fluctuation $`\delta N_𝐤^l`$ becomes as large as $`N_{eq}^l(𝐤)`$, and therefore a linearization of the Boltzmann equation (2.25) is no longer valid, like the estimation (9.6), corrected within the framework of this approximation. In the region $`|𝐤|gT`$, where collisions among plasmons start to play a role, it is necessary to solve the exact nonlinear integro-differential equation, whose r.h.s. is considered as both the process of the scattering plasmons off thermal particles, and the process of four-plasmon decays, i.e.
$$\frac{N_𝐤^l}{t}+𝐕_𝐤^l\frac{N_𝐤^l}{𝐱}=3(\omega _{pl}/g)^2\frac{\mathrm{d}𝐤_1}{(2\pi )^3}(\omega _𝐤^l\omega _{𝐤_1}^l)\mathrm{Q}(𝐤,𝐤_1)N_𝐤^lN_{𝐤_1}^l$$
(9.7)
$$+\frac{\mathrm{d}𝐤_1}{(2\pi )^3}\frac{\mathrm{d}𝐤_2}{(2\pi )^3}\frac{\mathrm{d}𝐤_3}{(2\pi )^3}(2\pi )^4\delta (\omega _𝐤^l+\omega _{𝐤_1}^l\omega _{𝐤_2}^l\omega _{𝐤_3}^l)\delta (𝐤+𝐤_1𝐤_2𝐤_3)$$
$$\times w(𝐤,𝐤_1;𝐤_2,𝐤_3)(N_{𝐤_1}^lN_{𝐤_2}^lN_{𝐤_3}^l+N_𝐤^lN_{𝐤_2}^lN_{𝐤_3}^lN_𝐤^lN_{𝐤_1}^lN_{𝐤_2}^lN_𝐤^lN_{𝐤_1}^lN_{𝐤_3}^l),$$
where
$$\mathrm{Q}(𝐤,𝐤_1)=(2\pi )^4N_c\left(\frac{\mathrm{Re}\epsilon ^l(k)}{\omega }\right)_{\omega =\omega _𝐤^l}^1\left(\frac{\mathrm{Re}\epsilon ^l(k_1)}{\omega _1}\right)_{\omega _1=\omega _{𝐤_\mathrm{𝟏}}^l}^1$$
$$\times \frac{\mathrm{d}\mathrm{\Omega }}{4\pi }\delta (\omega _𝐤^l\omega _{𝐤_1}^l𝐯(𝐤𝐤_1))|^c(𝐤,𝐤_1)+^{}(𝐤,𝐤_1)+^{}(𝐤,𝐤_1)|^2.$$
The amplitudes $`^c,^{}`$ and $`^{}`$ have the following forms
$$^c(𝐤,𝐤_1)=\frac{g^2}{|𝐤||𝐤_1|}\frac{1}{\omega _𝐤^l\omega _{𝐤_1}^l}\frac{(\mathrm{𝐤𝐯})(𝐤_1𝐯)}{\omega _𝐤^l(\mathrm{𝐤𝐯})},$$
$$^{}(𝐤,𝐤_1)=\frac{g^2}{|𝐤||𝐤_1|}\frac{((𝐤𝐤_1)𝐯)}{(𝐤𝐤_1)^2}$$
$$\times \left(\frac{(kk_1)^2{}_{}{}^{}\mathrm{\Delta }_{}^{l}(kk_1)}{\omega \omega _1(\omega \omega _1)^2}\delta \mathrm{\Gamma }^{ijl}(k,k_1,k+k_1)k^ik_1^j(kk_1)^l\right)_{\omega =\omega _𝐤^l,\omega _1=\omega _{𝐤_1}^l},$$
$$^{}(𝐤,𝐤_1)=\frac{g^2}{|𝐤||𝐤_1|}\frac{([𝐧(𝐤𝐤_1)]𝐯)}{𝐧^2(𝐤𝐤_1)^2}$$
$$\times \left(\frac{{}_{}{}^{}\mathrm{\Delta }_{}^{t}(kk_1)}{\omega \omega _1}^{}\mathrm{\Gamma }^{ijl}(k,k_1,k+k_1)k^ik_1^j[𝐧(𝐤𝐤_1)]^l\right)_{\omega =\omega _𝐤^l,\omega _1=\omega _{𝐤_1}^l}.$$
Here, $`𝐧[\mathrm{𝐤𝐤}_1]`$. We present the amplitudes $`^c,^{}`$ and $`^{}`$ in their simplest form, defined by temporal gauge.
A qualitative analysis of the solution behavior of a similar equation in the case of electron-ion plasma with a rather crude approximation of the kernel $`\mathrm{Q}(𝐤,𝐤_1)`$ and the probability of plasmon-plasmon scattering $`w(𝐤,𝐤_1;𝐤_2,𝐤_3)`$, was performed by Kovrizhnykh . He has shown, in particular, that really in the kinematic regime of small wavenumbers $`|𝐤|`$, the process of plasmon-plasmon scattering starts to play a dominant role, that prevents an increase of the $`𝐤=0`$-mode. In the case of the the non-Abelian plasma, an analogous qualitative investigation of the solution of equation (9.7) presents difficulties because of the above-mentioned complexity of the integrands on the r.h.s. of equation (9.7), and therefore it requires some additional approximations, allowing these functions to be made more visible and suitable for numerical calculations.
## 10 Conclusion
Within the framework of the semiclassical kinetic theory of a hot gluon plasma we have obtained the transport equation for the plasmons, taking into account four-plasmon decay. The probability of plasmon-plasmon scattering at the leading order in the coupling constant is derived. This is defined with the help of three-gluon and four-gluon effective vertices, and the effective propagator only, as the probability of the process of scattering of the plasmon by hard QGP particles . It is proved that the scattering probability is gauge-independent at least within a class of covariant and temporal gauges.
In this paper we have restricted our consideration to a derivation of the Boltzmann equation for the simplest collective Bose-modes of hot gluon plasma, colorless plasmons. This fact was reflected in deciding on a trivial color structure of the plasmon occupation number: $`N_𝐤^{lab}\delta ^{ab}`$. It is clear that such a choice of trivial color structure is completely neglected by the fundamental property of plasmons in the non-Abelian plasma, such as the availability of color charge. This is a very important difference from the Abelian case, where the plasmons do not carry electric charge. While the initial state of the problem can have a trivial color structure, the dynamics of these collective modes can change the color of the plasmons. In light of these remarks, here we are dealing with an Abelianized version of more complicated colored plasmon dynamics. It is a first step towards a full description of plasmon transport properties to derive the Boltzmann equation which takes into account the precession of the color charge of the plasmon. It is probably only in the construction of such a complete equation that we gain a reasonable explanation of the structure of the scattering probability (8.12).
The scheme of the derivation of the Boltzmann equation for colorless plasmons, proposed in this paper, admits the straightforward generalization to the case, when the number density of the plasmon $`N_𝐤^{lab}`$ has a non-trivial color structure, e.g. such as $`N_𝐤^{lab}=(T^c)^{ab}N_𝐤^{lc}`$. In this case, as was discussed in Sec. 4, it is necessary to assume that the regular (background) part of the field $`A_\mu ^R(X)`$ is different from zero. The density matrix, effective propagator, three-gluon, four-gluon vertices become the functionals of the background field $`A_\mu ^R(X)`$ , which assumed vanishing at $`X_0\mathrm{}`$. This significantly complicates the problem of construction of the required kinetic equation. It can somewhat simplify the problem if it is considered that the fluctuation $`\delta N_𝐤^{lab}`$ is a small perturbation in relation to the equilibrium value $`\delta ^{ab}N_{eq}^l(𝐤)`$, where $`N_{eq}^l(𝐤)`$ is the Planck distribution, and restrict the consideration to the linearized Boltzmann equation for colored plasmons. The exact consideration of a given question will be subject of separate research.
There is an independent test of the validity of the derived Boltzmann equation (2.25) (exactly, its linearized version (9.4)). As was mentioned in introduction in the fundamental derivation of the Boltzmann equation for a high temperature Yang-Mills plasma was proposed by Blaizot and Iancu within the CTP formalism framework. Their derivation relies on a gauge-covariant gradient expansion of the Kadanoff-Baym equations for the gluon two-point function. The Boltzmann equation has emerged as the quantum transport equation at leading order in $`g`$ for the gauge-covariant fluctuation $`\delta \stackrel{´}{G}`$ of a hard gluon propagator.
Besides the above-mentioned paper, the Kadanoff-Baym equations for the off-equilibrium propagator of the soft gluon $`𝒟_{\mu \nu }(X,Y)`$, which are formally identical to those for a hard gluon propagator $`G_{\mu \nu }(X,Y)`$, are written out. These equations are used in only to deduce the relation between the off-equilibrium gauge-covariant fluctuation $`\stackrel{´}{𝒟^<}(k,x)`$ and the gauge-covariant fluctuation of the leading-order soft gluon polarization tensor $`\stackrel{´}{\delta \mathrm{\Pi }^<}(k,x)`$, and the problem of self-interactions of the soft fields is not considered here. However, in principal there is nothing to forbid the use of these equations for research in the dynamics of the soft fields and construction of the relevant transport equation within the framework of the scheme, suggested by Blaisot and Iancu . Here, by $`\stackrel{´}{\delta \mathrm{\Pi }^<}(k,x)`$ we mean the fluctuation of the next-to-leading order of the soft gluon self-energy, involving three- and four-gluon off-equilibrium vertices with soft external lines. The relevant effects of self-interaction of the soft fields are encoded in these functions. Here, we note that in the real-time formalism, technical complications resulting from the doubling of degrees of freedom are arisen. However in recent years a suitable technique allowing one effectively to work not only with the single-particle propagator, but also with three- and four-point functions, has been developed . This greatly simplifies calculations in real time and enables us to derive the transport equation for the soft gluons, in particular, the plasmons, directly from the underlying quantum field theory and compare it with the equation, obtained in this paper in the context of the semiclassical approximation. This formal scheme is also needed to specify the limits of validity of the semiclassical kinetic approach to the research of the processes of nonlinear interaction of the soft fields in a hot QCD plasma.
## Acknowledgement
We are very grateful to Stanislaw Mrówczyński for useful correspondence. This work was supported in part by Grant INTAS (No. 2000-15) and Grant for Young Scientist of Russian Academy of Sciences (No. 1999-80).
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# Radio Sources in Low-Luminosity Active Galactic Nuclei. I. VLA Detections of Compact, Flat-Spectrum Cores
## 1 Introduction
There is increasing evidence that a large fraction of the nuclei of nearby galaxies show many similarities with powerful active galactic nuclei (AGN); these objects are termed low-luminosity active galactic nuclei (LLAGN; active galaxies with nuclear L$`{}_{H\alpha }{}^{}`$ 10<sup>40</sup> erg s<sup>-1</sup>; Ho, Filippenko & Sargent 1997a, hereafter H97a). These similarities include broad H$`\alpha `$ lines (Ho et al. 1997b), broader H$`\alpha `$ lines in polarized emission than in total emission (Barth et al., 1999), nuclear UV sources (Maoz et al., 1995; Barth et al., 1998) and water vapor megamasers (Braatz et al., 1997). However, the emission-line spectra of LLAGNs (i.e. low-luminosity Seyferts, LINERs, and “transition” nuclei \[nuclei with spectra intermediate between those of LINERs and H II regions\]), can also be modeled in terms of photoionization by hot, young stars (Terlevich & Melnick, 1985; Filippenko & Terlevich, 1992; Shields, 1992), by collisional ionization in shocks (Koski & Osterbrock, 1976; Fosbury et al., 1978; Heckman, 1980; Terlevich et al., 1992; Dopita & Sutherland, 1995), or by aging starbursts (Alonso-Herrero et al., 1999). Thus, it is not clear that accretion onto a black hole powers all LLAGNs.
How does one distinguish accretion-powered LLAGNs from LLAGNs powered by hot stars or supernova shocks? Broad H$`\alpha `$ lines and bright unresolved optical or UV sources are ambiguous indicators because they can all be produced in starburst models (see Terlevich et al. 1992), and a search for broader polarized H$`\alpha `$ emission is currently feasible in only a few of the brightest LLAGNs (Barth et al., 1999). Further, all these indicators are highly dependent on viewing geometry and obscuration and on the signal to noise of the observations, a problem exacerbated by the low optical and UV luminosities of LLAGNs. Their observation demands both high signal-to-noise spectra and high spatial resolution to separate weak nuclear emission lines from the starlight of the host galaxy.
One well-known property of some powerful AGNs is a compact, flat-spectrum nuclear radio source, usually interpreted as the synchrotron self-absorbed base of the jet which fuels larger-scale radio emission. Astrophysical jets are known to be produced in systems undergoing accretion onto a compact object (e.g. Pringle, 1993; Blandford, 1993), so such compact radio sources in galactic nuclei may reasonably be considered a signature of an AGN. Much theoretical work (e.g. Begelman, Blandford & Rees, 1984; Lovelace & Romanova, 1996; Falcke & Biermann, 1999) has been devoted to this disk-jet relationship in the case of galactic nuclei and it has been suggested that scaled-down versions of AGN jets can produce flat-spectrum radio cores in LLAGNs (Falcke, 1996; Falcke & Biermann, 1996). Nuclear radio emission with an inverted spectrum is also expected from an advection-dominated accretion flow (ADAF; e.g. Narayan et al., 1998), a possible form of accretion onto a black hole at low accretion rates (Rees et al., 1982), or from bremsstrahlung, cyclotron and synchrotron emission from plasma accreting quasi-spherically onto a black hole (Melia, 1994). From the observational perspective, Heckman (1980) showed that LINER nuclei tend to be associated with a compact radio source, and compact, flat-spectrum radio cores are known to be present in many ‘normal’ E/S0 galaxies (Sadler et al., 1989; Wrobel & Heeschen, 1991; Slee et al., 1994). Flat-spectrum radio cores are, however, uncommon in normal spirals or Seyfert galaxies (Ulvestad & Wilson, 1989; Vila et al., 1990; Sadler et al., 1995). Flat-spectrum radio sources can also result through thermal emission from optically-thin ionized gas or through free-free absorption of non-thermal radio emission, a process which probably occurs in compact nuclear starbursts (Condon et al., 1991). The brightness temperature, T<sub>b</sub>, in such starbursts is limited to log \[T$`_b`$(K)\] $``$ 5 (Condon et al., 1991). Thus it is necessary to show that T<sub>b</sub> exceeds this limit before accretion onto a black hole can be claimed as the power source.
In this paper, we present a high-frequency (15 GHz or 2 cm), high-resolution ($``$ 0$`\stackrel{}{\mathrm{.}}`$15) survey of LLAGNs with the Very Large Array (VLA, Thompson et al., 1980). This resolution is high enough to isolate nuclear emission from that of the host galaxy, and the radiation is unaffected by the obscuration present at UV or optical wavelengths, and less affected by free-free absorption than observations at longer cm wavelengths. Further, large samples can be quickly studied, as deep radio maps are achievable in as little as 15–20 min per object. Higher resolution follow-up observations are then required in order to eliminate the possibility that the radio cores are thermal in origin, as discussed above. We have therefore embarked on a program to observe a large number of LLAGNs at high resolution with the VLA and the Very Long Baseline Array (VLBA, Napier et al., 1994) in order to identify accretion-powered nuclei, to test the predictions of ADAF and jet models, and to characterize the presence and structure of radio jets on sub-pc to hundred-pc scales. In this paper, we report on the results of the first stage of this program — VLA observations of 48 LLAGNs which have been extensively observed at other wavebands. Preliminary results of this project have been published in Falcke et al. (1997, 1998, 1999) and Nagar et al. (1999b).
## 2 Sample, Observations and Data Reduction
A total of 48 galaxies were selected from the magnitude-limited Palomar spectroscopic survey of 486 nearby bright galaxies (Ho, Filippenko, & Sargent 1995). All 48 galaxies were originally selected, using preliminary spectral classification, to have nuclei with “pure” LINER or “transition object” (composite LINER + H II) characteristics. The spectral type of 8 of the nuclei was somewhat uncertain, and these galaxies were later reclassified as Seyferts (H97a). Thus, our sample of 48 LLAGNs consists of 22 “pure” LINERs, 18 “transition objects,” and 8 low-luminosity Seyferts. The 48 objects were not selected by well-defined criteria; they were individually chosen for study at multiple wavelengths because they were particularly “interesting” LLAGNs. Our sample is the subject of a rigorous observational campaign involving radio observations at 2 cm (reported here), 3.6 cm and 6 cm (Ho et al., 2000), Hubble Space Telescope (HST) optical (Peng et al., 1998) and UV imaging (Maoz et al., 1995; Barth et al., 1998) as well as spectroscopy (Maoz et al., 1998), and X-ray imaging by the Chandra observatory.
All 48 LLAGNs in the sample, plus two additional objects – NGC 4449 (whose nuclear spectrum is that of an H II region) and NGC 5756 (a southern spiral not in the Palomar spectroscopic survey) – were observed by the VLA over three runs on 1996 October 7, 11 and 18, while the array configuration was being changed from ‘D’ to ‘A’ (see Thompson et al. 1980 for a description of the VLA configurations). The weather was fair for all runs with clear skies and temperatures between 0 and 23°C. Two 8 min observations of each galaxy were sandwiched between three 2 min observations of a nearby phase calibrator, so that the total time on source was typically 16 min. For the October 7 run, there were 9 antennas at A-configuration pads plus 18 antennas at D-configuration pads. Five D-configuration antennas were taken out of the array at various times during the run to be moved to their A-configuration pads; they did not start observing again until after our run was completed. A sixth D-configuration antenna was removed from the array for painting. Of the nine antennas at A-configuration pads, the data from the three most recently moved were unusable, since the antenna positions were poorly determined. For the October 11 run, 25 of the 27 antennas were on their A-configuration pads, but the data from the four most recently moved antennas were unusable as their positions were poorly determined. For the October 18 run, all except one of the antennas were at their A-configuration pads, and all had reasonably well-determined positions.
The data were calibrated using the AIPS software, following the standard procedures outlined in the AIPS cookbook. Since positions for recently-moved antennas were not well determined at the time of observations, antenna positions listed in the data files were updated with the corrections to the positions found by VLA operators over the few weeks after the observations. We found that using all position corrections up to those determined on November 23, 1996, led to the best calibrator phase solutions for most of the recently-moved antennas. Other bad data were removed before the phase solutions, derived from the phase calibrators, were applied to the target galaxies. On the October 7 and October 18 runs, the source 1328+307 (3C 286) was observed twice and used to set the flux-density scale to that of Baars et al. (1977). On the shorter October 11 run, an observation of 0404+768 was used to set the flux-density scale (we used a 2 cm flux-density of 1.46 Jy for this source). The phase calibrator flux-densities were bootstrapped to that of 0404+768, and the bootstrapped values were found to be in good agreement with those listed in the VLA flux-density database (part of the “VLA calibrator manual”; available on-line at http://www.nrao.edu).
The maps were made with AIPS task IMAGR. For the October 7 run, we imaged each object twice, once using only the antennas on D-configuration pads, which resulted in $``$5″-resolution maps, and once using only those baselines which included at least one antenna on an A-configuration pad, which resulted in $``$0$`\stackrel{}{\mathrm{.}}`$15 resolution maps. For the other two runs, a single map was made for each source, using all antennas. For sources stronger than about 3 mJy, we were able to iteratively self-calibrate the data so as to increase the signal-to-noise ratio in the final map. Phase-only self-calibration was found to give good results, with little further improvement from a combined amplitude & phase self-calibration. The typical r.m.s. noise in the maps prior to self-calibration was 0.1 mJy. Many of the maps of the undetected nuclei show 4 to 5 $`\sigma `$ peaks at various locations in the map. We therefore conservatively use 10 times the r.m.s. noise in the map as the upper limit for the flux density of the undetected sources, so the sensitivity of our survey is about 1.1 mJy.
On the October 11 run, NGC 266, NGC 404, and NGC 7217 were also observed at 0.7 cm and 3.6 cm, with the VLA split into two subarrays which observed at different frequencies (not all antennas had 0.7 cm receivers). These data were reduced as described above. An observation of 0404+768 was used to set the flux-density scale at 0.7 cm. We set the 0.7 cm flux density of 0404+768 to 578.2 mJy for the AC IF’s and 575.3 for the BD IF’s, after finding, from the VLA flux-density database, that these were the flux densities on 1997 February 24, the date closest to our run that the source was observed (the flux density stayed between 500 and 600 mJy for the following $``$300 days). The flux density of the phase calibrators, 0055+300 and 2201+315, were bootstrapped to that of 0404+768, and the results were consistent with the values listed in the flux-density database. At 3.6 cm, the observation of 3C 48 could not be used to set the flux-density scale, as the four innermost antennas (and almost every alternate antenna on each arm) were observing at 0.7 cm, so that none of the 3.6 cm interferometer baselines was small enough for an accurate flux-density calibration solution (the VLA calibrator manual specifies that A-configuration observations of this calibrator should be restricted to 0–40 k$`\lambda `$). We therefore set the flux density of 2201+315 (the phase calibrator for NGC 7217) to 2.51 Jy in both AC and BD IF’s at 3.6 cm, as this was its measured value on 1997 January 17, the date closest to our run (after this date, the flux density of this source rose steadily to 3.5 Jy over the next $``$300 days and then stayed roughly constant). The flux density of 0055+300 (the phase calibrator for NGC 266 and NGC 404) was then bootstrapped from that of 2201+315. The result, S(0055+300) = 0.82 Jy, is consistent with the values listed in the flux-density database for this source. Nevertheless, the 3.6 cm and 0.7 cm flux densities of the three galaxies are subject to additional uncertainty.
All data were calibrated and imaged independently by two of us (NMN and HF), yielding consistent results (all corrections for antenna positions were also independently done). We are therefore confident of the results, despite the uncertainties caused by the uncertain antenna positions.
## 3 Results
### 3.1 Detection Rate
The results of the VLA observations are listed in Table 1, with columns explained in the footnotes. The radio positions for the 22 detected objects are limited by the positional accuracy of the phase calibrators, which is typically 2–10 milli-arcsec (mas), and by the accuracy of the Gaussian fit to the source brightness distribution, which depends on the signal to noise ratio of the source detection. The overall accuracy should typically be better than $``$50 mas for all target sources detected in A-configuration, and better than 500 mas for target sources detected in only D-configuration. We have compared the radio positions derived here with optical positions from Cotton, Condon, & Arbizzani (1999), which were measured from the digital sky survey with typical 1$`\sigma `$ accuracy 1$`\stackrel{}{\mathrm{.}}`$5–2$`\stackrel{}{\mathrm{.}}`$5 in each of right ascension and declination. For all but two of the 2 cm detected nuclei, the radio and optical positions agree to within 1″–3$`\stackrel{}{\mathrm{.}}`$5 (see Table 1). In NGC 185 and NGC 4550, the offsets between radio and optical positions are significant, and as discussed in Appendix A, it is possible that the radio sources we detected in these two galaxies are not related to the galaxy nucleus.
Assuming that the radio sources towards NGC 185 and NGC 4550 are in fact associated with their nuclei, the detection rate of 2 cm radio cores in the A-configuration maps of our LLAGN sample is 37% (18 of 48), at a 10$`\sigma `$ detection limit of 1.1 mJy. The rate is more striking if we consider the detection rates for low-luminosity Seyferts (75%; 6 of 8) and for LINERs (50%; 11 of 22). On the other hand, only one of the 18 observed ‘transition’ nuclei was detected. The morphological type of the galaxy with the detected transition nucleus, NGC 5866, is one of the earliest of the transition objects observed. An additional four objects were detected in D- but not A-configuration maps; we do not consider these four as nuclei with compact 2 cm radio cores. In comparison, Wrobel & Heeschen (1991) detected 30% (64 of 210) of all E/S0 galaxies with D $`<`$ 40 Mpc and m<sub>B</sub> $`<`$ 14 mag in a 6 cm VLA survey with resolution 5″ and 10$`\sigma `$ detection limit 1 mJy; the lower resolution and lower frequency make their survey more prone to detecting emission from sources unrelated to the central engine.
The radio luminosities of most of the detected 2 cm cores are between 10<sup>19.3</sup> and 10<sup>22</sup> W Hz<sup>-1</sup> (Fig. 1), similar to the luminosities seen in ‘normal’ E/S0 galaxies (Sadler et al., 1989). It is notable, however, that a significant fraction of the detected 2 cm compact cores are in spiral galaxies (Fig. 1).
The 0.7 cm maps of NGC 404, NGC 266, and NGC 7217 were, as expected, noisy, and all three galaxies were undetected at a 10$`\sigma `$ limit of 10 mJy. At 3.6 cm, NGC 266 and NGC 7217 were detected as unresolved sources with flux densities 3.1 mJy and 1.2 mJy, respectively, while NGC 404 was not detected at a 10$`\sigma `$ upper limit of 0.9 mJy.
### 3.2 Compact, Flat-Spectrum Cores
None of the objects show reliable extended structure in either A- or D-configuration maps. This is not surprising as the high resolution may resolve out most extended emission, which is expected to be weak at the high frequency observed. For a few sources, a Gaussian fit (with peak flux-density, major and minor axes as free parameters) to the nuclear radio source brightness distribution results in a Gaussian size slightly larger than the beam size so that the peak flux-density is slightly smaller than the total flux-density (see Table 1); however, the deconvolved source sizes are much smaller than half the beam size, so we consider these sources as unresolved. Of the 8 objects which are detected in both our A- and D-configuration maps, NGC 2655 is the only one with significantly more flux density at 5″ resolution than at 0$`\stackrel{}{\mathrm{.}}`$15 resolution. Thus, most of the detected 2 cm nuclear radio sources are compact at the 0$`\stackrel{}{\mathrm{.}}`$15 resolution (typically 15–25 pc) of our survey. The implied brightness temperatures for the 2 cm compact nuclei are T<sub>b</sub> $``$ 10<sup>3.0-4.0</sup> K except for NGC 4278 and NGC 6500 which have T<sub>b</sub> $``$ 10<sup>4.5</sup> K, and NGC 185, NGC 4548, NGC 4550, NGC 4636 and NGC 5033 which have T<sub>b</sub> $``$ 10<sup>2.5-3.0</sup> K.
We have supplemented our 2 cm observations with radio data at other wavelengths from the literature, such as the 6 cm radio survey of Wrobel & Heeschen (1991) and the 20 cm VLA FIRST survey catalog (White et al., 1997, hereafter FIRST), in order to derive non-simultaneous spectral indices for the 2 cm detected nuclei; details of the spectral index determination are listed in Appendix A for each galaxy detected at 2 cm. While the actual value of the spectral index is uncertain, given the resolution mismatch and the non-simultaneity of the observations, we can determine whether the core is flat spectrum ($`\alpha `$0.3; S$`{}_{\nu }{}^{}\nu ^\alpha `$), or steep spectrum ($`\alpha <`$0.3), as noted in column 13 of Table 1. Since compact flat-spectrum radio sources are often variable, the use of non-simultaneous data at two frequencies can cause some intrinsically flat-spectrum radio sources to be misclassified as steep-spectrum sources. However, since extended steep-spectrum radio sources are not variable, the use of non-simultaneous data at two frequencies should rarely cause intrinsically steep-spectrum radio sources to be misclassified as flat-spectrum sources. Also, since the resolution is better at the higher frequency, resolution effects will tend to steepen the measured spectrum if extended emission is present. Thus we are confident of the reality of the flat spectra obtained.
Fifteen of the 18 objects detected in our A-configuration 2 cm maps show a flat-spectrum radio core, and one has an undetermined radio spectrum. The remaining two objects detected in the A-configuration 2 cm maps, NGC 2655 and NGC 4636, show evidence for the presence of a steep-spectrum radio core, although NGC 4636 has “jet”-like steep-spectrum nuclear radio extensions (Stanger & Warwick, 1986) which may dominate over any flat-spectrum nuclear component within the beam of the radio maps. In the case of NGC 2655, the inferred steep-spectrum of the nucleus is consistent with the presence of extended 2 cm radio emission.
The radio morphology of the LLAGNs detected at 2 cm is different from what is seen in “classical” Seyferts (Ulvestad & Wilson, 1989; Nagar et al., 1999a), where the nuclear radio emission is usually dominated by steep-spectrum “jet” or double radio components on tens to hundreds of pc scales, and any pc-scale radio core may potentially be hidden by free-free absorption at cm wavelengths (e.g., Ulvestad et al., 1981; Gallimore et al., 1999). The high incidence of compact, flat-spectrum radio cores in the spiral galaxies in our sample is unusual (e.g., Vila et al., 1990; Sadler et al., 1995). In fact, the high incidence of flat-spectrum radio cores in low-luminosity Seyferts is also unusual, given that flat-spectrum compact radio cores are found in only about 10% of “classical” Seyferts (de Bruyn & Wilson, 1978; Ulvestad & Wilson, 1989; Sadler et al., 1995).
### 3.3 The Relationship Between Radio and Emission-line Luminosities
Correlations between the emission-line and radio luminosities of active galaxies are well known, and are found in both Seyfert (de Bruyn & Wilson, 1978) and radio galaxies (Baum & Heckman, 1989). The \[O III\] $`\lambda `$5007 line is usually used in assessing these correlations, but this line is weak in LINERs and we use instead \[O I\] $`\lambda `$6300, which also has the advantage of being uncontaminated by emission from H II regions. In our LLAGN sample, the 2 cm radio power appears closely correlated with the nuclear \[O I\] $`\lambda `$6300 luminosity (Fig. 2a and 2b). Testing the statistical significance of this apparent correlation is not straightforward as there exist both upper limits (line not detected) and highly uncertain values (non-photometric data) for the \[O I\] $`\lambda `$6300 luminosities, along with upper limits to the 2 cm radio powers. Eight of the LLAGNs in the sample have non-photometric \[O I\] $`\lambda `$6300 flux densities. We treated these eight flux densities as photometric measurements when using the bivariate tests from the ASURV statistical software package (Lavalley et al., 1992, hereafter ASURV). This approximation should not bias the results significantly as the uncertainties in the flux density are expected to be much less than an order of magnitude. When this is done, tests on the whole sample indicate a correlation between the radio and \[O I\] $`\lambda `$6300 luminosities at the 97% significance level. Deleting NGC 185 (the point at the bottom left corner of Fig. 2a), for which the association of the radio source with the nucleus is uncertain (Section 3.1), lowers the significance level by only 1%.
### 3.4 The 2 cm Non-detections
Do the LLAGNs undetected at 2 cm truly lack radio cores, or do a significant fraction of them have radio cores which are weaker than those in the detected LLAGNs because of, for example, a lower bulge luminosity (Sadler et al., 1989) or a later morphological type? The distributions of host galaxy bulge luminosity and distance (as listed in H97a) are more or less similar for the detected and undetected LINERs. The 2 cm compact core LINERs have a median host galaxy morphological type (T<sub>median</sub>= $``$1.5) which is earlier than, and a median \[O I\] $`\lambda `$6300 luminosity (Log L$`_{[OI]}`$(median) = 31.5 Watts) which is higher than, that for LINERs without compact 2 cm cores (T<sub>median</sub>= 2.0, Log L$`_{[OI]}`$(median) = 31.1 Watts), though the difference is not statistically significant. Given the correlation between emission-line and radio luminosity, it is possible that deeper radio imaging will find compact radio cores in the 2 cm undetected LINERs. The host galaxies of the two undetected Seyferts have bulge luminosities in the B band among the lowest of the eight Seyferts in our sample. There is no clear difference between the morphological types of the detected and undetected Seyferts.
Seventeen of the non-detected nuclei are “transition objects.” Statistical tests from the ASURV package indicate that the distributions of the 2 cm radio power of LINERs and transition objects in the sample are different at a confidence level of 99.3–99.5%. While the numerical value of the significance level is questionable given that there is only one detection among the transition objects, arbitrarily changing the four transition objects with the lowest radio flux-density upper limits into radio detections still results in a difference between transition object and LINER radio powers at the 92–98% significance level. While the transition nuclei have somewhat lower median bulge luminosities and generally later morphological types compared to the detected LINERs, the statistical significance of these differences is not high. Not unexpectedly, the transition objects in our sample do appear to have weaker \[O I\] $`\lambda `$6300 luminosities as compared to the LINERs and Seyferts (Fig. 2). If the non-photometric \[O I\] $`\lambda `$6300 flux measurements are considered to be accurate, statistical tests from the ASURV package indicate that in our sample the \[O I\] $`\lambda `$6300 luminosities of the transition objects are lower than that of LINERs and Seyferts at the $``$99.8% significance level. It remains to be seen whether this luminosity difference is large enough to explain the striking difference in the detection rates. Surprisingly, the one transition nucleus detected at 2 cm, NGC 5866 (the rightmost circle in Fig. 2b), has one of the lowest \[O I\] $`\lambda `$6300 luminosities among the transition nuclei.
### 3.5 Other AGN indicators
We find compact 2 cm radio cores in both type 1 and type 2 nuclei, that is, those with and without visible broad H$`\alpha `$ emission. Furthermore, some type 1 nuclei in our sample (two Seyferts and four LINERs) do not have a compact radio core at our detection limits. Among LINERs, 64% (7 of 11) of type 1 nuclei host compact radio cores, as compared to only 36% (4 of 11) of type 2 nuclei.
Additional evidence supports the accretion-powered nature of the nuclei with compact, flat-spectrum 2 cm cores. NGC 4579 is one of the very few LLAGNs with a detected Fe K$`\alpha `$ line (Terashima et al. 1998, 2000); such lines are believed to originate through fluorescence in an accretion disk. The optical–UV spectrum of NGC 4579 has been studied in detail by Barth et al. (1996), who show that the broad-line region in this object is detectable in a number of permitted transitions. The broad-band spectral energy distribution of NGC 4579 (Ho, 1999) has been interpreted in the context of an ADAF model (Quataert et al., 1999). NGC 4278 and NGC 6500 are both known to harbor pc-scale radio cores studied with Very Long Baseline Interferometry (VLBI) techniques (Jones, 1984; Jones et al., 1981; Falcke et al., 2000), and NGC 4203 (Shields et al., 2000), NGC 4278 and NGC 4636 (Magorrian et al., 1998) show kinematic evidence for massive dark objects at their centers. Maoz et al. (1998) find that the nuclei of NGC 404, NGC 4569, and NGC 5055, all undetected at 2 cm, contain massive stars which probably generate sufficient ionizing photons to power the optical emission lines, so it is not necessary to invoke an accretion-powered source. On the other hand, they find that the nuclei of M 81, NGC 4594, and NGC 6500 (and perhaps NGC 4579) exhibit a severe “deficit” in ionizing photons (i.e. apparently insufficient photons shortward of the Lyman limit, as inferred through extrapolation of the observed UV spectrum, to power the optical and UV emission lines), but their X-ray/UV ratios are two orders of magnitude higher than those of the three stellar-powered LINERs. Maoz et al. (1998) concluded that, in these four objects a component that emits primarily in the extreme-UV may be the main photoionizing source. We observed and strongly detected two of these nuclei (NGC 4579 and NGC 6500) at 2 cm, and the other two are known to have pc-scale radio cores (Bartel et al., 1982; Graham et al., 1981). Finally, the lack of clear evidence for X-ray emission from an AGN in NGC 404, NGC 4111, NGC 4192, and NGC 4569 (Terashima et al., 1999), galaxies undetected at 2 cm, is consistent with an absence of an accretion-powered source.
### 3.6 Minor-axis Jets or Outflows
A significant number of the LLAGNs with 2 cm compact radio cores in our sample show extended 20 cm radio emission in 5″ resolution FIRST maps, and/or 1$`\stackrel{}{\mathrm{.}}`$5–4″ maps from other sources (see Appendix A). We therefore compared the P.A. of this radio extension, when available, with that of the host galaxy major axis for all LLAGNs in our sample which host 2 cm compact radio cores. The relevant P.A.’s are listed in Table 2, and a histogram of the difference in the two P.A.’s, $`\delta `$, is shown in Fig. 3. The single transition object (shown in grey) shows extended emission along the host galaxy major axis, so it is likely that the radio emission is from the galaxy disk. On the other hand, LINERs and low-luminosity Seyferts with compact 2 cm radio cores appear to favor extended radio emission along the galaxy disk minor axis, although the small number of objects prohibits any meaningful statistical tests. The extent of this minor-axis radio emission is typically only 0$`\stackrel{}{\mathrm{.}}`$5 to 3″ (50–350 pc for most galaxies in the sample), and could trace the inner parts of wide-angle, minor-axis outflows, as seen in classical Seyferts (Colbert et al., 1996), or low-luminosity collimated “jets” from the central engine.
## 4 Discussion
We have detected unresolved 2 cm nuclear radio emission in 75% of low-luminosity Seyferts and 50% of LINERs. At least 83% of these nuclear radio sources have flat radio spectra. At this point it is not clear whether the 2 cm non-detected LINERs and low-luminosity Seyferts really do not have compact radio cores, or whether their compact cores simply have lower radio luminosities or are hidden by, for example, free-free absorption. In powerful AGNs, compact flat-spectrum nuclear radio cores are widely accepted as the signature of synchrotron emission from the base of a synchrotron self-absorbed jet (see the review by Zensus 1997). The compact, flat-spectrum radio cores we detect in LLAGNs could be scaled-down versions of powerful AGNs (e.g. Falcke et al., 1999), or could trace emission from an ADAF (Narayan et al., 1998), or material undergoing quasi-spherical accretion (Melia, 1992). Alternatively, the radio cores could be produced by thermal bremsstrahlung from ionized gas or by synchrotron radiation which has suffered free-free absorption. Such emission could originate from an active nucleus or nuclear star forming regions. The lower limits to the brightness temperatures, T$`{}_{}{}^{2cm}{}_{b}{}^{}`$, of the 2 cm compact core objects in our sample are 10<sup>2.5-4.5</sup> K, so the VLA observations alone cannot rule out the notion that the radio emission originates in star formation regions. However, most objects have brightness temperatures too high to be optically-thin thermal emission from a gas at 10<sup>4</sup> K. The expected value of T<sub>b</sub> is $``$ 10<sup>4</sup> K ($`\nu /\nu _\tau )^{2.1}`$, where $`\nu `$ is the observing frequency and $`\nu _\tau `$ is the frequency at which the source becomes optically thick. Many of our objects have a flat spectrum between 2 cm and 20 cm suggesting $`\nu _\tau `$ $`<`$ 1.4 GHz (20 cm). If this is the case, then thermal emission must have T$`{}_{b}{}^{2cm}<10^4(20/2)^{2.1}`$ K i.e. $`<`$ 90 K, lower than all of our measurements. Further, higher resolution observations of 10 of the detected objects in this sample reveal brightness temperatures $``$ 10<sup>8</sup> K, strongly suggesting that synchrotron self-absorption is responsible for the flat spectra and arguing against a thermal origin (Falcke et al., 2000). Thus, it is very likely that at least 50% of the LINERs and low-luminosity Seyferts in our sample contain accretion-powered nuclei. Several independent factors, described in Section 3.5, support our proposition that the presence of a compact, flat-spectrum (as determined from VLA observations) radio core is related to the presence of accretion onto a massive black hole in LLAGNs.
The detection rate of radio cores in transition objects in our sample is dramatically low (1 of 18). The significant difference between transition and “pure” LINER detection rates, combined with the lack of any significant difference between their host galaxy properties, suggests one of the following. 1) The central engines of transition objects are intrinsically different from those of “pure” LINERs, or 2) Transition objects are truly composite nuclei with a LINER and an H II component (e.g. Ho et al., 1993), in which case both the LINER component and its radio emission are much weaker in transition objects than in “pure” LINERs. Using the luminosity in the \[O I\] $`\lambda `$6300 line as an indicator of the luminosity of the LINER component, this second premise is supported by the results that the \[O I\] $`\lambda `$6300 and 2 cm luminosities are correlated at the $``$97% significance level, and that the \[O I\] $`\lambda `$6300 luminosity of the transition objects is lower than that of the LINERs and Seyferts in the sample at the $``$99.8% significance level (Fig. 2).
Finally, if the extended 20 cm radio emission, which appears preferentially aligned with the host galaxy minor axis in a subset of 2 cm compact core LLAGNs (Fig. 3), does trace jets from the central engine, there are interesting implications for the orientation of the central engines of AGN. Similar analyses of “classical” Seyferts (Ulvestad & Wilson, 1984; Schmitt et al., 1997; Clarke et al., 1998; Nagar & Wilson, 1999) have shown that the distribution of P.A.<sub>radio</sub> – P.A.$`_{galaxymajoraxis}`$ is more or less uniform between 0° and 90°. One explanation for this uniform distribution is that radiation from the central engine illuminates the inner accretion disk inducing a “radiative instability” (Pringle, 1997) which causes large warps in the inner disk. In this case, even if the outer accretion disk shares the same axis of rotation as the host galaxy disk, the jet (presumably launched along the axis of the inner accretion disk) may then be at a random angle with respect to the axis of the galaxy disk. If the radiation instability is indeed the correct explanation, then in LLAGNs, where the radiation field of the central engine is much weaker, the disk is expected to be less prone to the radiation instability (see e.g. equation 3.2 and Section 3.2 of Pringle 1997), so that the jet axis is more likely to be along the host galaxy rotation axis. This is precisely what we find. Alternately, the jet power in LLAGNs may be low enough that buoyancy dominates jet motions on the scale (typically 50–350 pc) of the extended emission, redirecting the jet along the minor axis. Our further surveys for 2 cm radio cores in LLAGNs and our completed observations of 17 nearby bright LLAGNs at 0$`\stackrel{}{\mathrm{.}}`$3 resolution with the Multi-Element Radio Linked Interferometer Network (MERLIN) at 20 cm will potentially reveal more on this subject.
## 5 Conclusions
Our detection of compact, flat-spectrum radio cores in about 50% of low-luminosity Seyferts and LINERs in a sample of 48 low-luminosity AGNs, when combined with VLBA observations (Falcke et al., 2000), suggests that at least half of all low-luminosity Seyferts and LINERs are accretion powered. Given the sensitivity limit of our survey, the true incidence of radio cores is likely to be higher. The 2 cm radio power is significantly correlated with the \[O I\] $`\lambda `$6300 luminosity for LLAGNs in our sample. The lower detection rate (at $``$92% significance) of compact radio cores in “transition” nuclei suggests that either these nuclei are intrinsically different from “pure” LINERs, or their “pure” LINER component is of lower luminosity, with correspondingly lower luminosity radio cores. The latter interpretation is favored by the correlation between the radio and \[O I\] $`\lambda `$6300 luminosities coupled with the lower \[O I\] $`\lambda `$6300 luminosities of the “transition” nuclei as compared to the LINERs.
The presence or absence of a detected broad H$`\alpha `$ line is not a good indicator of the presence or absence of a compact, flat-spectrum 2 cm radio core. A significant number of the LINERs and low-luminosity Seyferts which contain 2 cm compact radio cores show evidence, at low resolution (1″–5″) and frequency (1.4 GHz or 20 cm), for extended radio emission along the galaxy disk minor axis. This radio emission may trace a wide-angle outflow or a weak, highly-collimated jet along the disk rotation axis. If the latter is true, it lends support to the idea that it is the “radiative instability” which causes warps in the nuclear accretion disks of more luminous Seyfert galaxies.
Finally, we note that the data presented here are the initial results of a larger program to study a well-defined sample of LLAGNs at high resolution with the VLA and VLBA; results of these will appear in future papers in this series.
This research has been supported by NSF grant AST 9527289 and by NASA grant NAG 81027. The work of HF is funded by DFG grant Fa 358/1-1 & 2. The work of L. C. H. is partly funded by NASA grant NAG 5-3556, and by NASA grants GO-06837.01-95A and AR-07527.02-96A from the Space Telescope Science Institute (operated by AURA, Inc., under NASA contract NAS5-26555).
## Appendix A Notes on Individual Objects
Specific notes on all galaxies detected at 2 cm and on some undetected galaxies are listed here.
#### NGC 185
The 2 cm source we detect is offset 14$`\stackrel{}{\mathrm{.}}`$5 west and 4″ north of the optical position determined by Cotton et al. (1999) from digital sky survey plates. The 1$`\sigma `$ error of the optical position determination is estimated to be 2$`\stackrel{}{\mathrm{.}}`$7 in each of right ascension and declination. The galaxy is diffuse, $``$ 12′ x 10′ in extent, and has a prominent dust lane to the north-west which may bias the optical position determination. Overlaying the optical and radio images shows that the 2 cm radio source lies at the south-west tip of the dust lane, and the offset from the optical center does appear to be more than can be explained by dust lane obscuration. Further, the radio source has a very low luminosity (10<sup>16.7</sup> Watts Hz<sup>-1</sup>) which can be explained by a source in the galaxy stellar disk. Hummel (1980) lists a 5$`\sigma `$ upper limit of 10 mJy in a 23″ resolution map at 20 cm, and we find a 2 cm flux density of 0.8 mJy in our 0$`\stackrel{}{\mathrm{.}}`$15 map. It is therefore possible that the source is flat spectrum, but a measurement at a resolution closer to that of the 2 cm map is required for a definitive result, and for now we consider this source to have an undetermined spectral index.
#### NGC 266
Our simultaneous 3.6 cm and 2 cm observations on 1996 October 11, with resolutions 0$`\stackrel{}{\mathrm{.}}`$29 x 0$`\stackrel{}{\mathrm{.}}`$24, and 0$`\stackrel{}{\mathrm{.}}`$15 x 0$`\stackrel{}{\mathrm{.}}`$13, detected an unresolved nuclear source with flux densities 3.1 mJy and 4.1 mJy, respectively. The nucleus therefore has an inverted radio spectrum, $`\alpha _2^{3.6}`$ 0.5. This source was, however, not detected at a 10$`\sigma `$ upper limit of 10 mJy in our simultaneous 0.7 cm observation. NGC 266 has been identified with source NVSS J004947+321637 by the NASA Extragalactic Database (NED; Helou et al., 1991) i.e. it has been detected in the 20 cm NRAO VLA sky survey (NVSS), with a peak flux-density of 8.9 mJy/beam at a resolution of 45″, though the 20 cm position is offset 44″ from our 2 cm position. In any case, 8.9 mJy may be taken as an upper limit to the 20 cm flux density.
#### NGC 404
Our simultaneous 3.6 cm, 2 cm, and 0.7 cm observations on 1996 October 11 did not detect the nucleus at 10$`\sigma `$ upper limits of 0.9 mJy, 1.3 mJy, and 10 mJy, respectively.
#### NGC 2655
Comprehensive radio continuum and optical emission-line imaging, as well as optical spectra, are presented by Keel & Hummel (1988). Their VLA 6 cm map shows a roughly E-W symmetric extension at 0$`\stackrel{}{\mathrm{.}}`$5 resolution (with a peak flux-density of 36 mJy/beam from a component with half-power size 0$`\stackrel{}{\mathrm{.}}`$4 x 0$`\stackrel{}{\mathrm{.}}`$5). Combining this peak flux-density with our 2 cm peak flux-densities of 6 mJy/beam and 13.4 mJy/beam at resolutions of 0$`\stackrel{}{\mathrm{.}}`$15 and 5″, respectively, results in a non-simultaneous spectral index, $`\alpha _2^6`$ = $``$0.7 to $``$1.6, for the radio emission within 0$`\stackrel{}{\mathrm{.}}`$5 of the nucleus. Keel & Hummel (1988) find S$`{}_{\nu }{}^{}\nu ^{0.67}`$ for the nucleus, from 4″ resolution non-simultaneous VLA maps at 6 and 20 cm. They also find a secondary 6 cm source in P.A. 117°, with a diffuse bridge of radio emission joining it to the nuclear component. The morphology, kinematics, and line ratios of the surrounding emission-line gas, suggest that this secondary 6 cm source and bridge are related to nuclear ejecta (Keel & Hummel, 1988). Clearly, most of the radio emission in this object is extended, and if the emission from the unresolved radio nucleus is related to an accreting black hole, then its steep radio spectrum suggests that it is more likely that the emission is from synchrotron-emitting jets, rather than from a compact radio core. The host galaxy P.A. is not defined for this galaxy in the Uppsala General Catalogue of Galaxies (Nilson, 1973, hereafter UGC) or the Third Reference Catalogue of Bright Galaxies (de Vaucouleurs et al., 1991, hereafter RC3) as the galaxy is nearly round.
#### NGC 2681
FIRST lists a peak 20 cm flux density of 3.8 mJy/beam at 5″ resolution, with extended emission in P.A. 35°. We did not detect this object at 2 cm. The UGC and RC3 do not list a P.A. for this galaxy as it is almost round.
#### NGC 2787
Heckman et al. (1980) find 20 cm and 6 cm peak flux-densities of $`<`$9 mJy/beam and 9 mJy/beam, at resolutions of 4″ and 1$`\stackrel{}{\mathrm{.}}`$7, respectively, so the nucleus has a flat or inverted spectrum between 20 cm and 6 cm, $`\alpha _6^{20}`$ 0. Combining the 6 cm peak flux-density with our 2 cm peak flux-densities of 7 mJy/beam and 8.9 mJy/beam at resolutions of 0$`\stackrel{}{\mathrm{.}}`$15 and 5″, respectively, results in a non-simultaneous spectral index, $`\alpha _2^6`$ between $``$0.2 and 0, for all emission within the central 1$`\stackrel{}{\mathrm{.}}`$7.
#### NGC 3147
Vila et al. (1990) find S$`{}_{}{}^{peak}{}_{20cm}{}^{}`$ = 10.6 mJy/beam and S$`{}_{}{}^{peak}{}_{6cm}{}^{}`$ = 8.55 mJy/beam at resolutions of $``$1$`\stackrel{}{\mathrm{.}}`$2 and $``$1$`\stackrel{}{\mathrm{.}}`$5, respectively, with both maps showing only a compact core. Consistent with this, Laurent-Muehleisen et al. (1997) find a nuclear peak flux-density of 9 mJy/beam at 6 cm at a resolution of 0$`\stackrel{}{\mathrm{.}}`$4. Combining this 6 cm flux density with our 2 cm peak flux-density (8 mJy/beam), the non-simultaneous spectral index of the emission in the central 0$`\stackrel{}{\mathrm{.}}`$4 is flat, $`\alpha _2^6`$ 0.
#### NGC 3169
Hummel et al. (1987) find a core flux-density of 8.2 mJy in the central 2″ of their 1$`\stackrel{}{\mathrm{.}}`$2 resolution, 20 cm map, and also extended emission in P.A. 120°, more or less along the minor axis of the host galaxy. Combining their 20 cm peak flux-density with our measured 2 cm peak flux-densities of 6.8 mJy/beam and 10 mJy/beam at resolutions of 0$`\stackrel{}{\mathrm{.}}`$15 and 5″, respectively, results in a non-simultaneous spectral index, $`\alpha _2^{20}`$, between $``$0.1 and 0.1, for the radio emission within the central 2″.
#### NGC 4111
FIRST lists this object’s peak 20 cm flux density as 5.8 mJy/beam at 5″ resolution. The emission is extended in P.A. 70°, along the minor axis of the host galaxy. We did not detect this object at 2 cm.
#### NGC 4143
\- Wrobel & Heeschen (1991) find a 6 cm flux density of 6.7 mJy at 5″ resolution. The similar resolution (5″) FIRST map shows a peak flux-density of 5 mJy/beam at 20 cm, with extended emission in P.A. 10°, which is not along the minor axis of the host galaxy. Combining the FIRST peak flux-density with our 2 cm peak flux-densities of 3.3 mJy/beam and 5.3 mJy/beam at resolutions of 0$`\stackrel{}{\mathrm{.}}`$15 and 5″, respectively, results in a non-simultaneous spectral index, $`\alpha _2^{20}`$ between $``$0.2 and 0, for all emission within the central 5″. Ho et al. (1997b) note that broad H$`\beta `$ may also be present in this object.
#### NGC 4203
Wrobel (1991) derived a 6 cm flux density of 12.5 mJy at $``$5″ resolution, and FIRST lists a peak 20 cm flux density of 6.9 mJy/beam (at a resolution of 5″), implying that the core has a highly inverted spectrum between 20 cm and 6 cm, $`\alpha _6^{20}`$ 0.5. Our D-configuration 2 cm flux density of 12.3 mJy, suggests a non-simultaneous spectral index $`\alpha _2^6`$ 0. The extended emission in the FIRST map is in P.A. 90°, more or less along the minor axis of the host galaxy.
#### NGC 4216
Hummel et al. (1987) observed this object with 1$`\stackrel{}{\mathrm{.}}`$3 resolution at 20 cm and did not detect any emission in the central 2″ of the galaxy, at a 3$`\sigma `$ upper limit of 0.5 mJy.
#### NGC 4278
Jones (1984) have observed this object using VLBI and find the core flux density is 180 mJy and 190 mJy at 18 cm and 6 cm, on size scales $`<`$5 mas and $`<`$1 mas, respectively. Wilkinson et al. (1998) find a 3.6 cm peak flux-density of 153 mJy/beam at 200 mas resolution, and we find a 2 cm peak flux-density of 88.3 mJy/beam at 150 mas resolution, so the core may be flat between 18 and 6 cm and then steeper down to 2 cm. Alternatively, the lower flux density at 2 cm may be due to source variability. The 20 cm emission detected in FIRST is extended in P.A. 166°, but the UGC and RC3 do not list a host galaxy P.A. as the galaxy is almost round.
#### NGC 4419
Hummel et al. (1987) find a core flux-density of 32 mJy in the central 2″ of their 1$`\stackrel{}{\mathrm{.}}`$3 resolution 20 cm map, and extended emission in P.A. 77°. Combining this flux density with our 2 cm D-configuration peak flux-density of 5.4 mJy/beam at 5″ resolution, results in a non-simultaneous spectral index, $`\alpha _2^{20}`$0.8. This object was not detected in our 0$`\stackrel{}{\mathrm{.}}`$15 resolution, 2 cm A-configuration map.
#### NGC 4449
FIRST lists a 20 cm peak flux-density of 1.7 mJy/beam at 5″ resolution, with extended emission in P.A. 65°, more or less along the host galaxy major axis. We did not detect this object at 2 cm.
#### NGC 4527
Vila et al. (1990) find S$`{}_{}{}^{peak}{}_{20cm}{}^{}`$ = 4.3 mJy/beam and S$`{}_{}{}^{peak}{}_{6cm}{}^{}`$ = 1.3 mJy/beam at resolutions of $``$1$`\stackrel{}{\mathrm{.}}`$3 and $``$1$`\stackrel{}{\mathrm{.}}`$4, respectively (a spectral index, $`\alpha _6^{20}`$, of about $``$1.1). The 20 cm and 6 cm maps are both highly extended (P.A. 65°) along the major axis of the galaxy, with an edge brightened morphology, suggestive of an annulus (Vila et al., 1990). We did not detect this galaxy at 2 cm.
#### NGC 4548
Hummel et al. (1987) find a 3$`\sigma `$ upper limit of 2.5 mJy at 20 cm, over the central 2″ of their 1$`\stackrel{}{\mathrm{.}}`$3 resolution VLA map. Since we detect a flux density of 1.2 mJy at 2 cm, the non-simultaneous spectral index of the core is $`\alpha _2^{20}`$ $``$0.3.
#### NGC 4550
The 2 cm source we detect is offset 9$`\stackrel{}{\mathrm{.}}`$3 west and 5$`\stackrel{}{\mathrm{.}}`$5 north of the optical position determined by Cotton et al. (1999) from digital sky survey plates. The 1$`\sigma `$ error of the optical position determination is estimated to be 2$`\stackrel{}{\mathrm{.}}`$3 in each of the R.A. and Dec. The galaxy does not show any obvious sign of a peculiar morphology, and the east-west galaxy extent is $``$1′, so the radio-optical offset does appear significant. A $``$3″ resolution 11 cm map (Haynes & Sramek, 1980) shows an unresolved radio core at a position between our radio detection and the optical nucleus. Wrobel & Heeschen (1991) find the 6 cm emission over the central 5″ is $`<`$ 1 mJy. We measure a 2 cm peak flux-density of 0.7 mJy/beam, so the source is probably flat spectrum, $`\alpha _2^6`$ $``$0.3.
#### NGC 4565
Hummel et al. (1987) find a core 20 cm flux density of 1.2 mJy in the central 2″ of their 1$`\stackrel{}{\mathrm{.}}`$3 resolution VLA map. FIRST lists a peak flux-density of 1.5 mJy/beam at 20 cm in a 5″ resolution map, with extended emission in P.A. 36°, along the minor axis of the host galaxy. Combined with our 2 cm peak flux-density of 3.7 mJy/beam at 0$`\stackrel{}{\mathrm{.}}`$15 resolution, the non-simultaneous spectral index is $`\alpha _2^{20}`$ 0.5. This galaxy has the lowest broad H$`\alpha `$ luminosity found in any known active nucleus (Ho et al. 1997b).
#### NGC 4579
Hummel et al. (1987) find a core flux-density of 25.1 mJy in the central 2″ of their 1$`\stackrel{}{\mathrm{.}}`$3 resolution, 20 cm VLA map, with extended emission in P.A. $``$135°. The radio core was found to have an inverted spectrum between 13 cm and 3.6 cm with the Parkes-Tidbinbilla 275 km interferometer (S$`{}_{\nu }{}^{}`$ $`\nu ^{0.2}`$; Sadler et al. 1995). Van der Hulst, Crane & Keel (1981) found the 6 cm emission to be extended in P.A. 134°, with extent 0$`\stackrel{}{\mathrm{.}}`$4$`\pm `$0$`\stackrel{}{\mathrm{.}}`$2 and flux density 39.9 mJy.
#### NGC 4636
We derive a peak 2 cm flux density of 1.6 mJy/beam for this galaxy. Stanger & Warwick (1986) have imaged this object at 4″ resolution at 6 cm and 20 cm. Both maps show extended jet-like structure in P.A. $``$ 45°, more or less along the minor axis of the host galaxy. They found a peak 20 cm flux density of 11.4 mJy/beam, and a steep spectrum ($`\alpha `$ = $``$0.6) over the jets and at the core. The steep-spectrum jets may dominate the core emission within the central 4″ beam, so the core could potentially have a flat spectrum at higher resolution. In the absence of higher resolution maps, we list this object as a steep spectrum source in Table 1.
#### NGC 5005
Vila et al. (1990) find S$`{}_{}{}^{peak}{}_{20cm}{}^{}`$ = 14.6 mJy/beam and S$`{}_{}{}^{peak}{}_{6cm}{}^{}`$ = 2.7 mJy/beam at resolutions of $``$1$`\stackrel{}{\mathrm{.}}`$1 and $``$1$`\stackrel{}{\mathrm{.}}`$0, respectively (a spectral index of $`\alpha _6^{20}`$ = $``$1.8). The source is resolved at both wavelengths, in P.A. $``$ 135° on a $``$5″ scale, again more or less along the minor axis of the host galaxy. Hummel et al. (1985) find a 20 cm peak flux-density of 70 mJy/beam at 14″ resolution, with extended emission (about 2′ extent) in P.A. 67°, along the major axis of the galaxy. FIRST also lists this source with a peak 20 cm flux density of 35.9 mJy/beam (at 5″ resolution), with extended emission in P.A. 124°. The relatively steep spectral slope between 20 and 6 cm, and the extended radio morphology at these frequencies, is consistent with our 2 cm non-detection.
#### NGC 5033
FIRST lists a peak flux-density of 1.4 mJy/beam at 20 cm at a resolution of 5″. The inner radio emission in the FIRST map is extended in P.A. 91°, along the minor axis of the host galaxy, and there is an outer diffuse component which is more aligned with the host galaxy major axis. Combining the FIRST peak flux-density with our 2 cm peak flux-density (1.4 mJy/beam) results in a non-simultaneous spectral index $`\alpha _2^{20}`$ 0.
#### NGC 5055
The 20 cm FIRST map has peak flux-density 1.5 mJy/beam, and shows extended emission in P.A. 42°. We did not detect this galaxy at 2 cm.
#### NGC 5273
FIRST lists a peak flux-density of 2.6 mJy/beam at 20 cm; the extended radio emission (P.A. 180°) is along the major axis of the host galaxy disk. We did not detect this galaxy at 2 cm.
#### NGC 5866
Condon et al. (1990) find a peak 20 cm flux density of 9.5 mJy/beam in a 5″ resolution map made with the VLA. This is somewhat different from the value listed in FIRST - 20 cm peak flux-density of 15.4 mJy/beam at a resolution of 5″. Wrobel & Heeschen (1991) derived a flux density of 7.4 mJy from their 5″ resolution 6 cm map. We measured a peak 2 cm flux density of 7.1 mJy/beam at 0$`\stackrel{}{\mathrm{.}}`$15 resolution, so the core has a spectral index, $`\alpha _2^6`$ 0. The extended flux density in the FIRST map is in P.A. 132°, along the major axis of the host galaxy, not unexpected in this transition object. This is the only transition object in the sample in which we detected a 2 cm radio core.
#### NGC 6500
The spectrum of the core of this spiral is known to be flat on scales of $``$20 mas (Jones et al., 1981). The 1$`\stackrel{}{\mathrm{.}}`$2 resolution 20 cm map of Unger et al. (1989) shows extended emission in P.A. 140°, along the minor axis of the galaxy, suggestive of an outflow. Their higher resolution (0$`\stackrel{}{\mathrm{.}}`$3) 20 cm MERLIN map shows an extension in P.A. $``$70°.
#### NGC 7217
Our simultaneous 3.6 cm observation on 1996 October 11, with resolution 0$`\stackrel{}{\mathrm{.}}`$27 x 0$`\stackrel{}{\mathrm{.}}`$25, detected an unresolved nuclear source with flux density 1.2 mJy. Since the 3.6 cm detection gives us an accurate position for the radio core, we use a 5$`\sigma `$ upper limit of 0.6 mJy at 2 cm.
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